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warning/0004/math0004130.html	ar5iv	text	# Some results of regularity for Severi varieties of projective surfaces ## Introduction. Nodal curves play a central role in the subject of singular curves. The definition of the Severi variety of irreducible, nodal curves on any smooth, projective surface is standard. For a given effective divisor $`CDiv(S)`$, let $`C`$ denote the linear system associated to the line bundle $`\text{O}_S(C)Pic(S)`$. If we suppose that the generic element of $`C`$ is a smooth, irreducible curve, it makes sense to consider the subscheme $`V_{C,\delta }`$ of $`C`$, which parametrizes all curves $`C^{}C`$ that are irreducible and have only $`\delta `$ nodes as singular points. It is well known that such a subscheme is locally closed in the projective space $`C`$. In , Anhang F, Severi studied some properties of the variety $`V_{d,g}`$, defined as the closure of the locus consisting of reduced and irreducible plane curves, of geometric genus $`g`$, in the projective space parametrizing plane curves of degree $`d`$. $`V_{d,g}`$ contains, as an open dense subscheme, the locus $`V_{d,\delta }`$ corresponding to irreducible curves having only $`\delta `$ nodes as singularities. $`V_{d,g}`$ is classically known as the Severi variety of plane curves of given genus and degree. He proved that, for every $`d3`$ and $`0\delta \left(\begin{array}{c}d1\\ 2\end{array}\right)`$, $`V_{dL,\delta }`$ is non-empty and everywhere smooth of codimension $`\delta `$ in $`dL`$, where $`L`$ denotes a line in $`^2`$. Only after more than 60 years, Harris completed, in , Severi’s proof of the irreducibility of the Severi variety $`V_{d,g}`$ of the projective plane, by showing that the dense open subset $`V_{d,\delta }`$ is connected. With abuse of terminology, we shall use in the sequel the word Severi variety for the locally closed subscheme $`V_{C,\delta }`$ of the linear series $`C`$ on a projective surface $`S`$. In recent times, there have been many results on this subject and in many directions. In fact, one may study several problems concerning Severi varieties. Existence problems are covered by recent investigations in the case of Del Pezzo surfaces or K3 surfaces ; on the other hand, we have a complete computation of the degree of Severi varieties in the plane, treated in . There are some results even on the irreducibility problem, contained in , for general surfaces of $`^3`$ of degree $`d8`$. A natural approach to the dimension problem is to use deformation theory of nodal curves. Apart from the Severi classical result, whose proof can be extended to some rational or ruled surfaces and to K3 surfaces, information about regularity of $`V_{C,\delta }`$ on a surface $`S`$ of general type can be obtained by studying suitable rank 2 vector bundles on $`S`$. The first who used such approach were Chiantini and Sernesi. In they found an upper bound on $`\delta `$, ensuring that the family $`V_{C,\delta }`$ is smooth of codimension $`\delta `$ in the projective space $`C`$. Their proof focused on surfaces such that $`K_S`$ is ample and $`C`$ is a divisor which is numerically equivalent to $`pK_S`$, $`p^+`$ and $`p2`$. An improvement of this result is given in . The authors generalized this approach in two directions. In fact, they allowed arbitrary singularities and they weakened the assumption of $`K_S`$ being ample, so that $`S=^2`$ is included. Their assumptions are: $`C`$, $`CK_S`$ ample divisors and $`C^2K_S^2`$; moreover, some numerical hypotheses are made, which imply, in the case of nodes, that $`(C2K_S)^2>0`$ and $`C(C2K_S)>0`$. In this paper, we give a purely numerical criterion to prove the regularity of $`V_{C,\delta }`$ , provided that $`\delta `$ is less than a suitable upper bound. More precisely, we prove the following Theorem Let $`S`$ be a smooth, projective surface and $`C`$ be a smooth, irreducible divisor on $`S`$. Suppose that: 1. $`(C2K_S)^2>0`$ and $`C(C2K_S)>0`$; 2. either $$\begin{array}{cccc}K_S^2>4\hfill & if& C(C2K_S)8,\hfill & or\\ K_S^20\hfill & if& 0<C(C2K_S)<8.\hfill & \end{array}$$ 3. $`CK_S0`$; 4. $`H(C,K_S)<4(C(C2K_S)4)`$, where $`H(C,K_S)`$ is the $`Hodgenumber`$ (see def. 2) of $`C`$ and $`S`$; 5. $$\begin{array}{cccc}\delta \frac{C(C2K_S)}{4}1\hfill & if& C(C2K_S)8,\hfill & or\\ \delta <\frac{C(C2K_S)+\sqrt{C^2(C2K_S)^2}}{8}\hfill & if& 0<C(C2K_S)<8.\hfill & \end{array}$$ Then, if $`C^{}C`$ is a reduced, irreducible curve with only $`\delta `$ nodes as singular points and if $`N`$ denotes the $`0`$-dimensional scheme of nodes in $`C^{}`$, in the above hypotheses $`N`$ imposes independent conditions to $`C`$, i.e. the Severi variety $`V_{C,\delta }`$ is smooth of codimension $`\delta `$ at $`C^{}`$. We shall give some examples which show that such a result really generalizes the ones recalled before. Moreover, we can obtain some results on the Severi varieties of surfaces in $`^3`$ which are elements of a component of the Noether-Lefschetz locus, consisting of surfaces which contain a line. This is related to some results contained in , where the question of algebraic hyperbolicity for surfaces $`S`$ in $`^3`$ and in $`^4`$ is treated. In the sequel, we shall work in the category of $`schemes`$. We will denote by $``$ the linear equivalence of divisors, whereas $`_{num}`$ shall denote the numerical equivalence of divisors. ## Preliminaries Let $`S`$ be a projective, non-singular algebraic surface and $`D`$ a complete linear system on $`S`$ whose general member is a smooth, irreducible curve. If $`p_a(D)`$ denotes the $`arithmeticgenus`$ of $`D`$ then, by the adjunction formula, $$p_a(D)=\frac{D(D+K_S)}{2}+1.$$ For a given $`\delta 1`$, suppose that $`V_{D,\delta }`$ is non-empty. Let $`CV_{C,\delta }`$ and let $`N`$ be the scheme consisting of the $`\delta `$ nodes of $`C`$. The $`geometricgenus`$ of $`C`$ is $`g=p_g(C)=p_a(C)\delta .`$ We know that the Zariski tangent space of $`D`$ at the point $`[C]`$, parametrizing $`C`$, is isomorphic to $$H^0(S,\text{O}_S(D))/<C>,$$ whereas the Zariski tangent space to $`V_{D,\delta }`$ at $`[C]`$ is $$T_C(V_{D,\delta })H^0(S,\text{I}_N(D))/<C>,$$ where $`\text{I}_N\text{O}_S`$ denotes the ideal sheaf of the subscheme $`N`$ of $`S`$ (see, for example, ). The relative obstruction space is a subspace of $`H^1(S,\text{I}_N(D))`$. In particular, $`N`$ imposes independent conditions to $`D`$ if and only if $$dim(V_{D,\delta })=dimT_C(V_{D,\delta })=dimD\delta $$ at $`C`$. In this case, $`V_{D,\delta }`$ is smooth of the $`expecteddimension`$ at \[$`C`$\]. The component containing \[$`C`$\] will be called regular at \[$`C`$\]. Otherwise, it is said to be a superabundant component. The regularity property is very strong, since it implies that the nodes of $`C`$ can be independently smoothed (see or ). ###### Definition 1. Let $`S`$ be a smooth, projective surface and $`CDiv(S)`$. $`C`$ is said to be a nef divisor if $`CF0`$, for each effective divisor $`F`$. A nef divisor $`C`$ is called $`big`$ if $`C^2>0`$. We recall that, by the Kleiman criterion (see ), $`C`$ is nef if and only if it is in the closure of the ample divisor cone of $`S`$. ###### Definition 2. Let $`S`$ be a smooth, projective surface and $`CDiv(S)`$. We shall denote by $`H(C,K_S)`$ the Hodge number of $`C`$ and $`S`$, defined by $$H(C,K_S):=(CK_S)^2C^2K_S^2.$$ The Index theorem (see , page 120) ensures us that this number is a non-negative real number. ## The main result In the previous section we have recalled all definitions and properties needed to prove our principal result. We are now able to state the following ###### Theorem 1. Let $`S`$ be a smooth, projective surface and $`C`$ be a smooth, irreducible divisor on $`S`$. Suppose that: 1. $`(C2K_S)^2>0`$ and $`C(C2K_S)>0`$; 2. $$\begin{array}{ccccc}(i)\hfill & K_S^2>4\hfill & if& C(C2K_S)8,\hfill & or\\ (ii)\hfill & K_S^20\hfill & if& 0<C(C2K_S)<8.\hfill & \end{array}$$ 3. $`CK_S0`$; 4. $`H(C,K_S)<4(C(C2K_S)4)`$, where $`H(C,K_S)`$ is the $`Hodgenumber`$ of $`C`$ and $`S`$ (def.2); 5. $$\begin{array}{ccccc}(i)\hfill & \delta \frac{C(C2K_S)}{4}1\hfill & if& C(C2K_S)8,\hfill & or\\ (ii)\hfill & \delta <\frac{C(C2K_S)+\sqrt{C^2(C2K_S)^2}}{8}\hfill & if& 0<C(C2K_S)<8.\hfill & \end{array}$$ Then, if $`C^{}C`$ is a reduced, irreducible curve with only $`\delta `$ nodes as singular points and if $`N`$ denotes the $`0`$-dimensional scheme of nodes in $`C^{}`$, in the above hypotheses $`N`$ imposes independent conditions to $`C`$, i.e. the Severi variety $`V_{C,\delta }`$ is smooth of codimension $`\delta `$ at $`C^{}`$. ###### Proof. For the sake of simplicity we will write $`K`$, instead of $`K_S`$, to denote a canonical divisor of $`S`$. By contradiction, assume that $`N`$ does not impose independent conditions to $`C`$. Let $`N_0N`$ be a minimal 0-dimensional subscheme of $`N`$ for which this property holds and let $`\delta _0=N_0`$. This means that $`H^1(S,\text{I}_{N_0}(C))0`$ and that $`N_0`$ satisfies the Cayley-Bacharach condition (see, for example, ). Therefore, a non-zero element of $`H^1(\text{I}_{N_0}(C))`$ corresponds to a non-trivial rank 2 vector bundle $`EExt^1(\text{I}_{N_0}(CK),\text{O}_S)`$; so, one can consider the following exact sequence (1) $$0\text{O}_SE\text{I}_{N_0}(CK)0.$$ This implies that $$c_1(E)=CK,c_2(E)=\delta _0\delta ,$$ i.e. (2) $$c_1(E)^24c_2(E)=(CK)^24\delta _0.$$ We now want to compute $`(\text{2})`$ in cases $`5.(i)`$ and $`5.(ii)`$. In the first one, $$(CK)^24\delta _0(CK)^24\delta =C^22CK4+4+K^24\delta K^2+4>0,$$ by $`2(i)`$. In the other case, using $`5.(ii)`$ and the Index Theorem, $$(CK)^24\delta _0(CK)^24\delta =C^22CK+K^24\delta >K^20,$$ since we supposed $`2(ii)`$. In both cases, the vector bundle $`E`$ is $`Bogomolovunstable`$ (see or ), i.e. there exist $`M,BDiv(S)`$ and a 0-dimensional scheme $`Z`$ (possibly empty) such that (3) $$0\text{O}_S(M)E\text{I}_Z(B)0$$ holds and $`(MB)N(S)^+`$. We recall that $`N(S)^+`$ denotes the ample divisor cone of $`S`$. This means that $`(MB)^2`$ $`>0,`$ $`(MB)H`$ $`>0,`$ $`Hampledivisor.`$ The exact sequence $`(\text{3})`$ ensures us that $`H^0(E(M))0`$. If we consider the tensor product of the exact sequence $`(\text{1})`$ by $`\text{O}_S(M)`$, we get (5) $$0\text{O}_S(M)E(M)\text{I}_{N_0}(CKM)0.$$ We state that $`H^0(\text{O}_S(M))=0`$; otherwise, $`M`$ would be an effective divisor, therefore $`MH>0`$ for each ample divisor $`H`$. From $`(\text{3})`$, it follows that $`c_1(E)=M+B`$, so, by $`(\text{1})`$ and $`(4)`$, (6) $$MB=2MC+KN(S)^+.$$ Thus, for every ample divisor $`H`$, (7) $$MH>\frac{(CK)H}{2}.$$ Furthermore, from hypotheses $`1.`$ and $`3.`$, it immediately follows that $`C(CK)>0`$ and $`C^2>0`$. Indeed, $`C(CK)=C(C2K)+CK>0`$ and $`C^2>CK0`$. Since $`C`$ is irreducible, this implies that $`C`$ is a nef divisor; from $`(\text{7})`$ and from Kleiman’s criterion, we get (8) $$MC\frac{(CK)C}{2}.$$ It follows that $`MC<0`$ so, since $`C`$ is nef, $`M`$ can not be effective. If we consider the cohomological exact sequence associated to $`(\text{5})`$, we deduce that there exists a divisor $`\mathrm{\Delta }CKM`$ s.t. $`N_0\mathrm{\Delta }`$. If the irreducible nodal curve $`C^{}C`$, whose sets of nodes is $`N`$, were component of $`\mathrm{\Delta }`$, then $`MK`$ would be an effective divisor. By applying $`(\text{8})`$ and by using the fact that $`C(CK)>0`$ and hypothesis $`3.`$, one determines $$C^{}(MK)=C(MK)=CKCMCK\frac{(CK)C}{2}=$$ $$=\frac{(C+K)C}{2}=\frac{KC}{2}\frac{C^2}{2}<CK0,$$ which contradicts the effectiveness of $`MK`$, since $`C`$ is nef. Bezout’s theorem implies that (9) $$C^{}\mathrm{\Delta }=C^{}(CKM)2\delta _0.$$ On the other hand, taking $`M`$ maximal, we may further assume that the general section of $`E(M)`$ vanishes in a $`2`$-codimensional locus $`Z`$ of $`S`$. Thus, $`c_2(E(M))=deg(Z)0`$. By standard computations on Chern classes, we obtain $$c_2(E(M))=c_2(E)+M^2+c_1(E)(M)=\delta _0+M^2M(CK),$$ which implies (10) $$\delta _0M(CKM).$$ By applying the Index theorem to the divisor pair $`(C,\mathrm{\hspace{0.33em}2}MC+K)`$, we get (11) $$C^2(2MC+K)^2(C(CK)2C(CKM))^2.$$ From $`(\text{9})`$ and from the fact that $`C(CK)`$ is positive, it follows that (12) $$C(CK)2C(CKM)C(CK)4\delta _0.$$ We observe that the left side member of (12) is non-negative, since $`C(CK)2C(CKM)=C(2MC+K)`$, where $`C`$ is effective and, by (6), $`2MC+KN(S)^+`$. Thus, (12) still holds when we square both its members and, together with (11), this gives (13) $$C^2(2MC+K)^2(C(CK)4\delta _0)^2.$$ On the other hand, using $`(\text{10})`$, we get $$(2MC+K)^2=4(M\frac{(CK)}{2})^2=$$ $$=(CK)^24(CKM)M(CK)^24\delta _0,$$ i.e. (14) $$(2MC+K)^2(CK)^24\delta _0.$$ Putting together $`(\text{13})`$ and $`(\text{14})`$, we get (15) $$F(\delta _0):=16\delta _0^24C(C2K)\delta _0+(CK)^2C^2K^20.$$ Summarizing, the assumption on $`N`$, stated at the beginning, implies $`(\text{15})`$<sup>1</sup><sup>1</sup>1We remark that, in the case of nodes, this condition is the same of ; moreover, their hypotheses $`(1.2)`$ and $`(1.3)`$ coincide in the case of nodes and become $`F(\delta _0)<0`$.. We want to show that our numerical hypotheses hold if and only if the opposite inequality is satisfied. To this aim, observe that the discriminant of the equation $`F(\delta _0)=0`$ is $`16C^2(C2K)^2`$, so, by hypotheses $`1.`$ and $`3.`$, it is positive. The inequality $`F(\delta _0)<0`$ is verified iff $`\delta _0(\alpha (C,K),\beta (C,K))`$, where $$\alpha (C,K)=\frac{C(C2K)\sqrt{C^2(C2K)^2}}{8}and$$ $$\beta (C,K)=\frac{C(C2K)+\sqrt{C^2(C2K)^2}}{8};$$ we have to show that, with our numerical hypotheses, $`\delta _0(\alpha (C,K),\beta (C,K))`$. From $`5.`$, it immediately follows that $`\delta _0<\beta (C,K)`$, since, as we shall see in the sequel, the bound in $`5.(i)`$ is smaller than $`\beta (C,K)`$. Note also that $`\alpha (C,K)0`$. Indeed, if $`\alpha (C,K)<0`$, then $`C(C2K)<\sqrt{C^2(C2K)^2}`$, which contradicts the Index theorem, since $`C(C2K)>0`$. Observe that $`\alpha (C,K)<1`$ if and only if (16) $$C(C2K)8<\sqrt{C^2(C2K)^2}$$ To simplify the notation, we put $`t=C(C2K)`$ so that $`(\text{16})`$ becomes (17) $$t8<\sqrt{t^24H(C,K)}.$$ Two cases can occur. If $`t8<0`$, there is nothing to prove since the right side member of (17) is always positive. Note, before proceeding to consider the other case, that in this situation we want that $`\beta (C,K)>1`$ in order to have at least an effective positive integral value for the number of nodes; but $`\beta (C,K)>1`$ if and only if $`(0<)8t<\sqrt{t^24H(C,K)}`$. By squaring both members of the previous inequality, we get $`4H(C,K)<16t64`$, which is our hypothesis 4.; so the upper-bound for $`\delta `$ is surely greater than 1. Moreover, the expression for such bound is the one in $`5.(ii)`$ and it can not be written in a better non-trivial form. On the other hand, if $`t80`$, by squaring both members of $`(\text{17})`$, we get $`H(C,K)<4(C(C2K)4)`$, which is our hypothesis 4.. Therefore, $`\alpha (C,K)<1`$; moreover, the condition $`\beta (C,K)>1`$ is trivially satisfied, since it is equivalent to $`t8>\sqrt{C^2(C2K)^2}`$. From $`(\text{16})`$, we can write $$\frac{C(C2K)+C(C2K)8}{8}<\frac{C(C2K)+\sqrt{C^2(C2K)^2}}{8},$$ so we can replace the bound $`\delta <\beta (C,K)`$ with the more ”accessible” one $`\delta \frac{C(C2K)}{4}1`$, which is the bound in $`5.(i)`$. Observe that $$\frac{C(C2K)+C(C2K)8}{8}<\frac{C(C2K)+\sqrt{C^2(C2K)^2}}{8}$$ $$\frac{C(C2K)}{4},$$ so, it is not correct to directly write $`\delta <\frac{C(C2K)}{4}`$. Therefore, $`5.(i)`$ is the right approximation. In conclusion, our numerical hypotheses contradict (15), therefore the assumption $`h^1(\text{I}_N(C))`$ $`0`$ leads to a contradiction. ∎ ###### Remark 1. 1) The previous theorem gives purely numerical conditions to deduce some informations about Severi varieties of smooth projective surfaces. In the next section, we shall discuss some interesting examples of projective surfaces for which the theorem can be easily applied. More precisely, we will consider smooth surfaces in $`^3`$ which are general elements of a component of the Noether-Lefschetz locus; for example, surfaces of general type, of degree $`d5`$, which contain a line. Our result obviously generalizes the one of Chiantini and Sernesi. In their case, since $`C_{num}pK_S`$, $`p`$ and $`p2`$, we always have $`\alpha (C,K_S)=0`$ and $`\beta (C,K_S)=\frac{p(p2)}{4}K_S^2`$; this depends on the fact that $`H(pK_S,K_S)=0`$, for every $`p`$. With the further hypotheses that $`p^+`$, $`p`$ odd, and that the Neron-Severi group of $`S`$ is $`NS(S)[K_S]`$, they proved that one can take $`\delta =\frac{(p1)^2}{4}K_S^2`$. These bounds are sharp, at least for the general quintic surface in $`^3`$. Furthermore, as recalled in Remark 1.2 in , in the case of rational or ruled surfaces (for which $`CK_S<0`$) or $`K3`$ surfaces (for which $`CK_S=0`$) if $`C`$ is base-point-free the argument for $`S=^2`$ can be repeated without changes, since the line bundle $`N_\phi `$, on the normalization of the nodal curve $`C^{}C`$, is non-special. Our result focuses on cases in which $`CK_S0`$ (see hypothesis $`3.`$), where the previous approach fails. 2) One can immediately deduce that when the Hodge number is zero, i.e. when we are considering a divisor pair such that $`(CK)^2=C^2K^2`$, then in the previous proof we find $`\alpha (C,K)=0`$ and $`\beta (C,K)=\frac{C(C2K)}{4}`$. 3)Theorem 1 generalizes, in the case of nodes, the result in . This will be clear after having considered the following examples. Examples: 1) Let $`S^3`$ be a general smooth quartic. We have, $`\text{O}_S(K_S)\text{O}_S`$. Let $`H`$ denote the plane section and $`D`$ be the generic element of $`2H`$. From Bertini’s theorem, $`D`$ is smooth and irreducible. If $`\pi :\stackrel{~}{S}S`$ denotes the blow-up of $`S`$ in a point $`PS`$ and $`E`$ the associated exceptional divisor, then $`K_{\stackrel{~}{S}}E`$, i.e. the canonical divisor of the blown-up surface is linearly equivalent to the exceptional divisor. Thus, $`C2\pi ^{}(H)`$ can not be ample, since $`CK_{\stackrel{~}{S}}=0`$; so, the first hypothesis in does not hold. Nevertheless, observe that the generic element of $`C`$ is smooth and irreducible. Moreover, $`C2K_{\stackrel{~}{S}}2\pi ^{}(H)2E`$ so that $`(C2K_{\stackrel{~}{S}})^2=12`$, $`C(C2K_{\stackrel{~}{S}})=16`$, $`CK_{\stackrel{~}{S}}=0`$, $`K_{\stackrel{~}{S}}^2=1`$, $`H(C,K_{\stackrel{~}{S}})=16`$ and $`4(C(C2K_{\stackrel{~}{S}})4)=48`$. Since we are in the situation $`5.(i)`$, we get $`\delta \frac{16}{4}1=3`$, i.e. on $`\stackrel{~}{S}`$, if $`V_{2\pi ^{}(H),\delta }\mathrm{}`$ and if $`\delta 3`$, then it is everywhere smooth of the expected dimension. 2) Let $`S`$ be a smooth quintic surface in $`^3`$ which contains a line $`L`$. Denote by $`\mathrm{\Gamma }S`$ a plane quartic which is coplanar to $`L`$, so that $`\mathrm{\Gamma }HL`$ ($`H`$ denotes the plane section). Thus, $$H^2=5,HL=1,L^2=3,H\mathrm{\Gamma }=4,\mathrm{\Gamma }^2=0and\mathrm{\Gamma }L=4.$$ Choose $`C3H+L`$, so that $`C`$ contains the curves which are residue to $`\mathrm{\Gamma }`$ in the complete intersection of $`S`$ with the smooth quartic surfaces of $`^3`$ containing $`\mathrm{\Gamma }`$. $`3H+L`$ is base-point-free and not composed with a pencil, since $`(3H+L)L=0`$ and $`3H`$ is an ample divisor. By Bertini’s theorems, its general member is smooth and irreducible; but $`C`$ and $`CK_S`$ can not be both either ample or, even, nef divisors. In fact, $`CL=0`$ and $`(CK_S)L=(2H+L)L=1`$. Therefore, the result in can not be applied. Neverthless, $`CK_S=C(C2K_S)=H(C,K_S)=16`$, $`(C2K_S)^2=3`$, $`K_S^2=5`$, $`4(C(C2K_S)4)=48`$ and, since $`C(C2K_S)>8`$, $`\delta \frac{16}{4}1=3`$. Thus, if $`3H+L`$ contains some nodal, irreducible curves, then, if $`\delta 3`$, $`V_{3H+L,\delta }`$ is everywhere smooth of the expected dimension. ## Some results on surfaces in $`^3`$ which contain a line We now consider a class of examples to which our result can be easily applied. We shall focus on surfaces of $`^3`$ containing a line. Such approach can be generalized to surfaces belonging to other components of the Noether-Lefschetz locus. Firstly, let $`S^3`$ be a smooth quintic and $`LS`$ a line. Since $`p_a(L)=p_g(L)=0`$, by the adjunction formula and by the fact that $`K_SH`$ we get $`L^2=3`$. As before, $$K_S^2=5,LH=1,L^2=3.$$ We are interested on some results of regularity for Severi varieties of curves on $`S`$, which are residue to the line $`L`$ in the complete intersection of $`S`$ with a general surface of degree $`a`$ passing through the line. Thus $`CaHL`$ on $`S`$. By straightforward computations, we get $$deg(C)=(aHL)H=5a1,$$ $$p_a(C)=\frac{5a^2+3a}{2}1.$$ We want to find conditions on $`a`$ in order to apply our result. (i) $`C`$ has a smooth and irreducible general member: For the smoothness, we have to prove that $`aHL`$ is base-point-free and not composed with a pencil. Since $`a1`$, $`aHL=(a1)H+HL`$. If $`a2`$, the linear system $`(a1)H`$ can not have fixed intersection points on $`L`$. We can restrict ourselves to treat the behaviour of $`HL`$ on $`L`$. If $`HL`$ admitted fixed points on $`L`$, each of those points should be a tangent point for $`S`$ and the general plane of $`^3`$ passing through the line. This would imply that $`S`$ is a singular surface in such points, which contradict the hypothesis. Moreover, $`HL`$ can not be composed with a pencil, since $`HL+LH`$. For the irreducibility, we can use the fact that $`C`$ and $`L`$ are linked in $`^3`$ (see ). In fact, this implies that $`C`$ is $`projectivelynormal`$ in $`^3`$, i.e. if we consider the exact sequence $$0\text{I}_{C/^3}(\rho )\text{O}_^3(\rho )\text{O}_C(\rho )0,$$ then $`H^1(\text{I}_{C/^3}(\rho ))=0`$, for each $`\rho `$. By choosing $`\rho =0`$, we get $`H^0(\text{I}_{C/^3})=H^1(\text{I}_{C/^3})=0`$ so $`H^0(\text{O}_C)H^0(\text{O}_^3)`$. This proves that $`C`$ is a connected curve; since we have already proven its smoothness, then the general member is also irreducible. (ii) Numerical hypotheses: By simple computations, one observes that all the numerical conditions in Theorem 1 simultaneously hold if $`a4`$. We can completely generalize the previous procedure to the case of a smooth surface of degree $`d6`$ which contains a line $`L`$. For the detailed computations, the reader is referred to . Let $`S^3`$ be such a surface and $`CaHL`$, so that $$deg(C)=ad1,$$ $$p_a(C)=\frac{ad(a+d)2ad(4a+1)+3}{2}.$$ Moreover, $`L^2=2d`$, since $`K_S(d4)H`$ and $`LH=1`$. For the smoothness and the irreducibility of the general member of $`C`$ one can use the previous argument. Now, $`K_S^2=(d4)^2d24`$, because $`d6`$. It is not difficult to compute (see for details) that, for $`6d7`$, all the numerical hypotheses in Theorem 1 simultaneously hold if $`a2d6`$ (note that, for $`d=5`$ we obtained $`a4`$ so that $`d=5,\mathrm{\hspace{0.33em}6},\mathrm{\hspace{0.33em}7}`$ behave in the same way). On the other hand, for $`d8`$, the condition on the Hodge number (i.e. hypothesis $`4.`$ in Theorem 1) determines a bound on $`a`$ which is bigger than the one determined by the other conditions, i.e. $`2d6`$. Indeed, condition $`4.`$ holds if and only if $$4a^2d8a(d^24d+1)(d^410d^3+33d^244d+56)>0.$$ By solving this inequality, we find (18) $$a>d4+\frac{1}{d}+\frac{1}{2}\sqrt{d^36d^2+d+28+\frac{24}{d}+\frac{4}{d^2}}.$$ It is a straightforword computation to find that the right side member of $`(\text{18})`$ is bigger than $`2d6`$ when $`d8`$. Therefore, in this case, all the numerical conditions of Theorem 1 simultaneously hold if $`(\text{18})`$ holds. In order to find a better expression for such a lower-bound on $`a`$, we observe that $$\sqrt{d^36d^2+d+28+\frac{24}{d}+\frac{4}{d^2}}<\sqrt{d^36d^2+d+32},$$ since $`d8`$. We are looking for a real number $`b`$ such that $`\sqrt{d^36d^2+d+32}\sqrt{(d\sqrt{d}b)^2}`$. For such a value, we have (19) $$2b\sqrt{d}6d1+\frac{b^232}{d}.$$ Moreover, $`(\text{18})`$ becomes (20) $$ad3+\frac{d}{2}\sqrt{d}\frac{b}{2}.$$ Obviously, the right side member of $`(\text{20})`$ must be greater than $`2d6`$ for $`d8`$. Observe that this happens if and only if (21) $$d\sqrt{d}>2d+b6.$$ Therefore, putting $`\phi (d):=d\sqrt{d}2d+6`$, from $`(\text{21})`$ we have $`b<\phi (d)`$. The function $`\phi (d)`$ is monotone increasing for $`d2`$ so, to find a uniform bound on $`b`$ for all the cases in $`d8`$, it is sufficient to consider $`b<\phi (8)`$, i.e. $`b12`$. By taking into account $`(\text{19})`$, we find that in all cases a good choice is $`b=9`$. Thus, $`(\text{18})`$ can be replaced by $`ad3+\frac{d\sqrt{d}9}{2}`$. Analogous computations show that, when $`d5`$, only condition 2.(i) can occur, i.e. $`C(C2K_S)8`$. Therefore the expression for the bound on the number of nodes is the one in 5.(i), which is $$\delta \frac{a^2d2a(d^24d+1)+2d15}{4}.$$ Now, by summarizing all we have observed up to now, we are able to state the following ###### Proposition 1. Let $`S`$ be a smooth surface in $`^3`$ of degree $`d5`$, which contains a line $`L`$. Consider on $`S`$ the linear series $`aHL`$, with 1. $`a2d6`$, if $`5d7`$; 2. $`ad3+\frac{d\sqrt{d}9}{2}`$, if $`d8`$. (We denote by $`x`$ the round-up of the real number $`x`$, i.e. the smallest integer which is bigger than or equal to $`x`$). Suppose, also, that for a given integer $`\delta `$ the Severi variety $`V_{aHL,\delta }`$ is non-empty. Then, if $$\delta \frac{a^2d2a(d^24d+1)+2d15}{4},$$ the Severi variety is everywhere smooth of the expected dimension. ###### Remark 2. We want to point out that our results, in a certain sense, agree with what is stated in . In fact the authors proved the following result. Theorem Let $`D`$ be a reduced curve in $`^3`$ and $`s`$, $`d`$ be two integers such that $`ds+4`$. Moreover, suppose that i)there exists a surface $`Y^3`$ of degree $`s`$ which contains $`D`$; ii) the general element of the linear system $`\text{O}_Y(dHD)`$ is smooth and irreducible. Denote by $`S`$ a general surface of $`^3`$, of degree $`d`$, containing $`D`$. Thus, $`S`$ does not contain reduced, irreducible curves $`CD`$ of geometric genus $`g<1+deg(C)\frac{(ds5)}{2}`$. In particular, if $`ds+6`$ and $`p_g(D)2`$, $`S`$ is algebraically hyperbolic. In the case of our proposition, $`S`$ is a surface of degree $`d5`$ and $`D=L`$, such that $`L^2=2d`$. Thus, we can consider $`s=1`$, i.e. $`Y`$ is a plane containing the line $`L`$ and $`\text{O}_Y(dHL)=\text{O}_^2(d1)`$ which has a smooth and irreducible general element. Therefore, if there exists a curve $`C`$ of a given degree, then $$p_g(C)2+\frac{(d6)}{2}deg(C)=2+\frac{(d6)}{2}CH.$$ If, moreover, $`C`$ is a nodal curve, then $$\delta =p_a(C)p_g(C)\frac{C^2+CK_S}{2}+12\frac{(d6)}{2}CH=$$ $$\frac{C^2}{2}+\frac{(d4)}{2}CH\frac{(d6)}{2}CH1=\frac{C^2}{2}+CH1.$$ On the other hand, since in such cases, when all our hypotheses are satisfied, $`C(C2K_S)8`$, $`5(i)`$ determines $$\delta \frac{C(C2dH+8H)}{4}1=\frac{C^2}{4}\frac{(d4)}{2}CH1.$$ Observe that $`\frac{C^2}{4}\frac{(d4)}{2}CH1\frac{C^2}{2}+CH1`$ if and only if $`\frac{C^2}{4}+CH(\frac{d}{2}1)0`$. Since $`d5`$ and since $`C`$ is big and nef (consequence of $`1.`$ and $`3.`$), this latter inequality is always strictly verified. This means that our bounds on $`\delta `$ are in the range of values, for the number of nodes, that are necessary for the existence of such a curve.
warning/0004/hep-th0004091.html	ar5iv	text	# References Effective Quarks and Their Interactions<sup>1</sup><sup>1</sup>1Talk presented by D. McMullan Emili Bagan<sup>2</sup><sup>2</sup>2email: bagan@ifae.es Physics Department Brookhaven National Laboratory Upton NY 11973 USA Robin Horan<sup>3</sup><sup>3</sup>3emails: rhoran@plymouth.ac.uk, mlavelle@plymouth.ac.uk and dmcmullan@plymouth.ac.uk , Martin Lavelle<sup>3</sup> and David McMullan<sup>3</sup> School of Mathematics and Statistics The University of Plymouth Plymouth, PL4 8AA UK > Abstract: This talk will summarise the progress we have made in our programme to both characterise and construct charges in gauge theories. As an application of these ideas we will see how the dominant glue surrounding quarks, which is responsible for asymptotic freedom, emerges from a constituent description of the interquark potential. Introduction The phenomenological confidence in the existence of coloured hadronic constituents is in marked contrast to the theoretical uncertainties associated with attempts to describe such charges. Although there are many models of partonic and effective degrees of freedom in the literature, none have yet emerged directly from the fundamental gauge theoretic description of the strong interactions — QCD. The source of the difficulty in directly extracting these effective degrees of freedom from the underlying gauge theory is a particular example of the basic dichotomy we all face in QCD: the degrees of freedom that make up the QCD Lagrangian and successfully probe the ultra-violet regime are not related in any obvious way to the large scale, infra-red degrees of freedom that describe the observed hadrons or their constituents. The picture that has emerged<sup>4</sup><sup>4</sup>4For a nice discussion of the stunning experimental results achieved over the past decade see . from studies of deep inelastic scattering, and more recently from diffractive processes, is that, as we probe smaller and smaller sub-hadronic scales, we go from wholly hadronic degrees of freedom to constituent structures (quarks). In their turn, these constituents have properties (for example, their mass) that run as shorter distance scales are probed. As such they are not viewed as fundamental fields but as composites made up in some way from the matter and gluonic degrees of freedom that enter the QCD Lagrangian. There is an immediate theoretical problem for any coloured constituent: how does such a quark or gluon have a well defined colour given that, as is shown in the picture below, it is made up of some mixture of apparently coloured partonic degrees of freedom? In order to answer this fundamental question and hence introduce our approach to the construction of such (colour) charged degrees of freedom, we need to recall how and when colour can be defined in QCD. Colour and Gauge Invariance The very structure of QCD as a gauge theory tells us that physical fields must be invariant under local<sup>5</sup><sup>5</sup>5In QCD we actually expect the stronger statement to emerge that physical fields are invariant under all gauge transformations: both local and global, in other words colour charges are supposed to be confined in (hadronic) colour singlets. We will see at the end of this article how this stronger result actually emerges directly from our approach to the constructions of colour charges. gauge transformations. This apparent truism generates though an immediate problem when we look at the colour charge itself: $$Q^a=d^3x(J_0^a(x)f_{bc}^aE_i^b(x)A_i^c(x)).$$ (1) This is clearly not gauge invariant! So can we talk about colour in any meaningful way? The answer can be seen to be yes, when we recognise that the question we should be asking is whether the colour charge is gauge invariant when restricted to physical (i.e., gauge invariant) states. On these states, the non-abelian version of Gauss’ law implies that $$Q^a=\frac{1}{g}d^3x_iE_i^a(x).$$ (2) Under a gauge transformation $`E_i^aT^aU^1E_i^aT^aU`$ so that $$Q^aT^a\frac{1}{g}d^3x_i(U^1E_i^aT^aU).$$ (3) We can now write this as the surface integral of the chromo-electric flux in any given direction and hence we see that, on gauge invariant states, the colour charge transforms as $$Q^aT^a\frac{1}{g}\underset{R\mathrm{}}{lim}_{S^2}𝑑\underset{¯}{s}U^1\underset{¯}{}U.$$ (4) Hence the colour charge will be gauge invariant under local gauge transformations if, at spatial infinity, we have $`UU_{\mathrm{}}`$ where $`U_{\mathrm{}}`$ lies in the centre of SU(3). Continuity then tells us that this group element will be a constant independent of the direction taken to spatial infinity. This imposes a $`_3`$ (triality) structure on possible charged states. However, we will only concern ourselves with the zero triality sector in this talk where $`U_{\mathrm{}}`$ is the identity. To summarise the above discussion, we have seen that in order to be able to define a coloured object, such as a quark, we require that it must be gauge invariant and also that allowed gauge transformations must be restricted as above. We now want to study the construction of charged fields. Construction of Charges In QCD we have the gauge transformations $`A_\mu (x)`$ $``$ $`U^1(x)A_\mu (x)U(x)+{\displaystyle \frac{1}{g}}U^1(x)_\mu U(x)`$ $`\psi (x)`$ $``$ $`U^1(x)\psi (x)`$ From the non-triviality of these transformations we see that neither of these fields have a well defined colour and hence they cannot be identified with observed gluonic or quark degrees of freedom. The generic form for a charged (matter) field is given by a process we call dressing the matter to give the product $$h^1(x)\psi (x),$$ (5) where, under a gauge transformation, the dressing transforms as $$h^1(x)h^1(x)U(x).$$ (6) This is the minimal condition we must impose on the dressing in order for the corresponding charged matter field to be gauge invariant and hence having a well defined colour . But does this really means we have a physical field? Is gauge invariance alone enough? Consider the stringy gauge invariant $`e^+e^{}`$ state $$\|\overline{\psi }(x)\mathrm{exp}(ie_x^y𝑑wA(w))\psi (y)$$ (7) This is gauge invariant, but is it physical? For static matter, the energy is the expectation value of the Hamiltonian $`\frac{1}{2}d^3z(E^2(z)+B^2(z))`$. This yields for the potential energy of this state the confining potential $$V(xy)e^2\|xy\|\delta ^2(0).$$ (8) Given that we are here dealing with four dimensional QED, it is clear that the original stringy state cannot be accepted as a physical configuration! Indeed it is an infinitely excited state and what we would want to construct is the ground state for the system. To do this in a systematic fashion, we need a further condition apart from gauge invariance. The question which we now address is: what condition on the dressing gives a stable charged particle? In order to motivate this extra condition, we first consider $`\phi `$ to be a heavy field which creates a particle at the point $`x`$ with a given 4-velocity $`u`$. The field must be constant along the trajectory of a particle moving with this velocity which leads to the equations of motion: $$u^\mu _\mu \phi (x)=0.$$ (9) If $`\phi `$ is now a heavy gauged field then its equation of motion becomes $$u^\mu D_\mu \phi (x)=0.$$ (10) A physical heavy coloured field can only emerge from dressing this field, $`\varphi =h^1\phi `$. If this is truly a heavy field, then it should furthermore satisfy the equation $$u^\mu _\mu \varphi (x)=0.$$ (11) This means that the dressing must satisfy the dressing equation $$u^\mu _\mu (h^1)=gh^1u^\mu A_\mu .$$ (12) One can in fact show that this equation applies to any theory with massive charges . Electric Charges It is instructive to first study the dressing process in the simpler case of QED. Our two inputs into the construction are then, as we have just seen, gauge invariance $`h^1h^1\mathrm{e}^{ie\theta }`$ and our kinematical requirement (the dressing equation) $$u^\mu _\mu (h^1)=ieh^1u^\mu A_\mu $$ (13) The great advantage of QED is that we can explicitly solve these equations to obtain $$h^1=\mathrm{e}^{ieK}\mathrm{e}^{ie\chi },$$ (14) where $`K(x)={\displaystyle _\mathrm{\Gamma }}𝑑\mathrm{\Gamma }(\eta +v)^\mu {\displaystyle \frac{^\nu F_{\nu \mu }}{𝒢}},`$ (15) $`\chi (x)={\displaystyle \frac{𝒢A}{𝒢}},`$ (16) with $`\eta =(1,\underset{¯}{0})`$, $`v=(0,\underset{¯}{v})`$, $`𝒢^\mu =(\eta +v)^\mu (\eta v)^\mu `$ and where $`\mathrm{\Gamma }`$ is the trajectory of the particle. We thus see that the dressing has two structures: a gauge dependent part, $`\chi `$, which makes the whole charge gauge invariant, and is thus in some sense minimal, and a further gauge invariant part, $`K`$ which is needed (together with the precise form of $`\chi `$) to satisfy the dressing equation. These structures are reflected in physical calculations: $`\chi `$ removes soft divergences in QED calculations and, as we will see, in QCD generates the anti-screening interaction responsible for asymptotic freedom. $`K`$ removes the phase divergences in the on-shell Green’s functions of QED<sup>6</sup><sup>6</sup>6The cancellation of the various IR divergences will be presented in the talk by Martin Lavelle.. For greater insight into these structures, let us consider the specific case of a static charge, $`v=0`$. This is: $$h^1\psi (x)=\mathrm{e}^{ieK}\mathrm{e}^{ie\chi }\psi (x)$$ (17) with now $`K(x)={\displaystyle _{\mathrm{}}^{x^0}}𝑑t{\displaystyle \frac{^\nu F_{\nu 0}}{^2}}`$ (18) $`\chi (x)={\displaystyle \frac{_iA_i}{^2}}`$ (19) and where $$\frac{1}{^2}f(t,\underset{¯}{x}):=\frac{1}{4\pi }d^3y\frac{f(t,\underset{¯}{y})}{\|\underset{¯}{x}\underset{¯}{y}\|}$$ (20) The non-locality of any description of a charge is manifest here. The minimal part part of the electromagnetic cloud around a static charge was first found by Dirac . The additional structure does not affect the electric field of the charge and was therefore not picked up by Dirac’s original argument. We will now show that such charged (dressed) matter is free at large times and that we can so recover a particle description . To see why this is important, we now recall that Kulish and Faddeev showed that, at large times the matter field is not free, but rather becomes<sup>7</sup><sup>7</sup>7For more details of this method and a refinement of their work, see the talk by Robin Horan $$\psi (x)\frac{d^3p}{(2\pi )^3}\frac{D(p,t)}{\sqrt{2E_p}}\left\{b(𝒑,s)u^s(𝒑)e^{ipx}+\mathrm{}\right\},$$ (21) where $`D`$ is a distortion factor. This implies that there is no particle picture. Of course since the coupling does not asymptotically vanish $`\psi `$ is not gauge invariant even at large times and so we should not expect to relate it to a physical particle! However, when we extract the annihilation operator for our dressed field we obtain $$b(q)\{1+e_{\mathrm{soft}}\frac{d^3k}{(2\pi )^3}(\frac{Va}{Vk}\frac{qa}{qk})\mathrm{e}^{itkq/E_q}\mathrm{c}.\mathrm{c}.\}+O(e^2),$$ (22) with $`V^\mu =(\eta +v)^\mu (\eta v)kk^\mu `$. There are two corrections now: the usual one from the interactions of the matter field and another one from the dressing. It is easy to show that at the right point on the mass shell, $`q=m\gamma (\eta +v)`$, these distortions cancel! We thus see that our dressed matter asymptotically corresponds to free fields and we regain a particle picture. Colour Charges After this construction of abelian charges, we now want to proceed to the non-abelian theory. We recall that the minimal static dressing in QED was: $`\mathrm{exp}(ie\chi )`$, with $`\chi =_iA_i/^2`$. This vanishes in Coulomb gauge and this observation lets us generalise this dressing to QCD where it may be extended to an arbitrary order in $`g`$ (see the Appendix of . Indeed we can also extend to non-static charges, but in the application that follows we require static quarks. In QCD we write the dressing as a perturbative expansion $$\mathrm{exp}(ie\chi )\mathrm{exp}(g\chi ^aT^a)h^1$$ (23) with $`g\chi ^aT^a=(g\chi _1^a+g^2\chi _2^a+g^3\chi _3^a+\mathrm{})T^a`$ The dressing gauge argument mentioned above implies $$\chi _1^a=\frac{_jA_j^a}{^2};\chi _2^a=f^{abc}\frac{_j}{^2}\left(\chi _1^bA_j^c+\frac{1}{2}(_j\chi _1^b)\chi _1^c\right)$$ (24) etc. This can be extended to all orders in the coupling. However, we will return below to the question of non-perturbative solutions. The Interquark Potential We now want to study the interaction energy of the ground state in the presence of a matter field and its antimatter equivalent . To construct this we can either dress a quark and an antiquark separately or dress a single meson in which there are no gauge invariant constituents. We will now dress the quark fields and study the potential between them: whether or not this gives the correct interaction energy, i.e., the potential, is a test of the validity of a constituent picture. We recall from our discussion of the ‘stringy’ state that this is a sensitive test. Our procedure is as follows: as sketched above we first extend our expression for the minimally dressed quark to higher orders in perturbation theory. We then take such minimally dressed quark/antiquark states, $`\overline{\psi }(y)h(y)h^1(y^{})\psi (y^{})\|0`$, and sandwich the Hamiltonian, $$H=\frac{1}{2}(E_i^aE_i^a+B_i^aB_i^a)d^3x$$ (25) between them. Using the standard equal-time commutators $$[E_i^a(x),A_j^b(y)]_{\mathrm{e}t}=i\delta ^{ab}\delta (𝒙𝒚),$$ (26) we can then calculate the potential. The lowest order result, i.e., at order $`g^2`$, is just the Coulomb potential: $$V^{g^2}(r)=\frac{g^2NC_F}{4\pi r},$$ (27) where $`r`$ is the separation of the matter fields. This is of course just QED with coloured icing. What about QCD with non-abelian ingredients? Well at order $`g^4`$ we need to calculate the minimal static dressing to order $`g^3`$. This can be done with the above mentioned efficient algorithm. A relatively simple calculation then yields for the potential at $`g^4`$ $$V^{g^4}(r)=\frac{g^4}{(4\pi )^2}\frac{NC_FC_A}{2\pi r}4\mathrm{log}(\mu r).$$ (28) What does this tell us about our dressed state? We recall that the QCD potential may be extracted from a Wilson loop as follows: $$V(r)=\underset{t\mathrm{}}{lim}\frac{1}{it}\mathrm{log}0\|\mathrm{Tr}\text{P}\mathrm{exp}\left(g𝑑x_\mu A_a^\mu T^a\right)0$$ (29) At order $`g^4`$ this yields $$V(r)=\frac{g^2C_F}{4\pi r}\left[1+\frac{g^2}{4\pi }\frac{C_A}{2\pi }\left(4\frac{1}{3}\right)\mathrm{log}(\mu r)\right].$$ (30) From this we may read off the universal one-loop beta function $$\beta (g)=\frac{g^3}{(4\pi )^2}\left[4\frac{1}{3}\right].$$ (31) We have decomposed it in this way because it has been shown by a number of authors that the dominant anti-screening contribution (the 4) which is responsible for asymptotic freedom comes from longitudinal glue and the screening part (the $`\frac{1}{3}`$) from gauge invariant glue We now recognise that the result (28) for the interaction energy of the minimally dressed quark/antiquark system is just the anti-screening contribution to the interquark potential. We thus make the important identification that the dominant part of the glue in a $`Q\overline{Q}`$ system which is responsible for anti-screening actually factorises into two individually gauge invariant constituents. The success of constituent models can, to the extent that these low order calculations have been carried though, be explained by our work. We postulate that the introduction of gauge invariant glue, via the incorporation of the phase dressings, will produce the screening effect. Topology and Confinement Having seen the perturbative efficiency and relevance of our variables, we now want to study the non-perturbative sector. Experimentally of course we do not see free quarks, and it can be easily shown that this is in fact predicted by our method. This follows from an intimate link between dressings and gauge-fixing. In the limited space available this can be best explained pictorially. Essentially the existence of a gauge fixing, $`\chi `$, can be shown to imply the existence of a function $`h`$ which transforms as a minimal dressing must to make a quark gauge invariant Similarly we can show that having a dressing implies that we can construct a gauge fixing which slices every gauge orbit once: But Gribov and Singer have showed that, with certain boundary conditions which we have seen above are needed if colour is to be a good quantum number, there is no such good gauge fixing in QCD! Hence any constituent picture must break down! We conclude that there is no non-perturbative, gauge invariant description of a single quark or gluon. Therefore outside of certain dynamical domains we cannot expect to see individual quarks or gluons. Conclusions There are strong phenomenological reasons in both QED and QCD for wanting to be able to describe charged particles. We saw that if we want to describe coloured substructure then this requires both gauge invariance and some restrictions on the allowed gauge transformations. We have seen a method of constructing dressed charges which was designed to describe physical charges with a well defined velocity. It had two inputs: gauge invariance and a kinematical requirement to single out which of the potentially many gauge invariant constructions involving a single matter field corresponds to the ground state. This approach, we have shown, generates structured dressings around charged particles. It has long been argued that it is impossible to describe charged particles in gauge theories . Essentially this is because asymptotic matter fields are not free fields due to the long range nature of the interactions transmitted by massless gauge bosons. However, we have demonstrated that asymptotic dressed fields are indeed free fields (the dressing takes such effects into account). We thus have obtained for the very first time a particle description of charges. It will be shown elsewhere in this meeting that this removes the IR divergences in QED at the level of on-shell Green’s functions. We have calculated the potential between two minimally dressed quark fields. The anti-screening contribution to the interquark potential was shown to be generated by the minimal dressing. We have thus determined the dominant glue configuration around quarks. We have seen that in a ‘meson’ – at least to the level at which we calculate – two separate, gauge invariant, coloured objects are visible. Furthermore there is a topological obstruction to the construction of coloured charges. This means that we can directly demonstrate the non-observability of individual quarks or gluons: outside the domain of perturbation theory and some non-perturbative effects, the Gribov ambiguity will show itself and there will be no locally gauge invariant description of such objects. This shows a new way to calculate the scale of confinement: we need to find out at which stage it becomes impossible to factorise the dressing of, say, a $`Q\overline{Q}`$ state into two individual charges. Beyond this breakdown of the factorisation quarks do not exist in QCD. Acknowledgements: This research was partly supported by the British Council/Spanish Education Ministry Acciones Integradas grant no. Integradas grant 1801 /HB1997-0141. It is a pleasure to thank both the local organisers for their hospitality and also PPARC for a travel grant.
warning/0004/hep-ph0004051.html	ar5iv	text	# Primordial Galactic Magnetic Fields from Domain Walls at the QCD Phase Transition. ## Introduction. The source of cosmic magnetic fields with large scale correlations has remained somewhat of a mystery . There are two possible origins for these fields: primordial sources and galactic sources. Primordial fields are produced in the earlier universe and then evolve and are thought to provide seeds which gravitational dynamos later amplify. Galactic sources would produce the fields as well as amplify them. Many mechanisms have been proposed , however, most fail to convincingly generate fields with large enough correlation lengths to match the observed microgauss fields with $`100`$ kpc correlations. We present here a mechanism which, although probably requiring a dynamo to produce microgauss fields, generates fields with hundred kiloparsec correlations. We present this mechanism as an application of our recent understanding of QCD domain walls, which will be described in detail elsewhere . 1. Sometime near the QCD phase transition, $`T_{\mathrm{QCD}}1`$ GeV, QCD domain walls form. 2. These domain walls rapidly coalesce until there remains, on average, one domain wall per Hubble volume with Hubble-scale correlations. 3. Baryons interact with the domain walls and align their spins along the domain walls. 4. The magnetic and electric dipole moments of the baryons induce helical magnetic fields correlated with the domain wall. 5. The domain walls decay, leaving a magnetic field. 6. As the universe expands, an “inverse cascade” mechanism transfers energy from small to large scale modes, effectively increasing the resulting correlation length of the observed large scale fields. We shall start by discussing the “inverse cascade” mechanism which seems to be the most efficient mechanism for increasing the correlation length of magnetic turbulence. After presenting some estimates to show that this mechanism can indeed generate fields of the observed scales, we shall discuss the domain wall mechanism for generating the initial fields. ## Evolution of Magnetic Fields. As suggested by Cornwall , discussed by Son and confirmed by Field and Carroll , energy in magnetic fields can undergo an apparent “inverse cascade” and be transfered from high frequency modes to low frequency modes, thus increasing the overall correlation length of the field faster than the naïve scaling by the universe’s scale parameter $`R(T)`$. There are two important conditions: turbulence must be supported as indicated by a large Reynolds number $`\mathrm{Re}`$, and magnetic helicity (Abelian Chern-Simons number) $`H=\stackrel{}{𝐀}\stackrel{}{𝐁}\mathrm{d}^3x`$ is approximately conserved. The importance of helicity was originally demonstrated by Pouquet and collaborators . The mechanism is thus: the small scale modes dissipate, but the conservation of helicity requires that the helicity be transfered to larger scale modes. Some energy is transfered along with the helicity and hence energy is transported from the small to large scale modes. This is the inverse cascade. The reader is referred to for a more complete discussion. In the early universe, Re is very large and supports turbulence. This drops to $`\mathrm{Re}1`$ at the $`e^+e^{}`$ annihilation epoch, $`T_0100`$ eV . After this point (and throughout the matter dominated phase) we assume that the fields are “frozen in” and that the correlation length expands as $`R`$ while the field strength decays as $`R^2`$. Note that the inverse cascade is only supported during the radiation dominated phase of the universe. Under the assumption that the field is maximally helical, these conditions imply the following relationships between the initial field $`B_{\mathrm{rms}}(T_i)`$ with initial correlation $`l(T_i)`$ and present fields today ($`T_{\mathrm{now}}2\times 10^4`$ eV) $`B_{\mathrm{rms}}(T_{\mathrm{now}})`$ with correlation $`l(T_{\mathrm{now}})`$ : $`B_{\mathrm{rms}}(T_{\mathrm{now}})`$ $`=`$ $`\left({\displaystyle \frac{T_0}{T_{\mathrm{now}}}}\right)^2\left({\displaystyle \frac{T_i}{T_0}}\right)^{7/3}B_{\mathrm{rms}}(T_i)`$ (1) $`l(T_{\mathrm{now}})`$ $`=`$ $`\left({\displaystyle \frac{T_0}{T_{\mathrm{now}}}}\right)\left({\displaystyle \frac{T_i}{T_0}}\right)^{5/3}l(T_i).`$ (2) As pointed out in , the only way to generate turbulence is either by a phase transition $`T_i`$ or by gravitational instabilities. We consider the former source. As we shall show, our mechanism generates Hubble size correlations $`l_i`$ at a phase transition $`T_i`$. In the radiation dominated epoch, the Hubble size scales as $`T_i^2`$. Combining this with (2), we see that $`l_{\mathrm{now}}T_i^{1/3}`$; thus, the earlier the phase transition, the smaller the possible correlations. The last phase transition is the QCD transition, $`T_i=T_{\mathrm{QCD}}0.2`$ GeV with Hubble size $`l(T_{\mathrm{QCD}})30`$ km. We calculate (9) the initial magnetic field strength to be $`B_{\mathrm{rms}}(T_i)e\mathrm{\Lambda }_{QCD}^2/(\xi \mathrm{\Lambda }_{\mathrm{QCD}})(10^{17}\mathrm{G})/(\xi \mathrm{\Lambda }_{\mathrm{QCD}})`$ where $`\xi `$ is a correlation length that depends on the dynamics of the system as discussed below and $`\mathrm{\Lambda }_{\mathrm{QCD}}0.2`$ GeV. With these estimates, we see that $$B_{\mathrm{rms}}\frac{10^9\mathrm{G}}{\xi \mathrm{\Lambda }_{\mathrm{QCD}}},l100\mathrm{kpc}$$ (3) today. One could consider the electroweak transition which might produce $`100`$ pc correlations today, but this presupposes a mechanism for generating fields with Hubble-scale correlations. Such a mechanism does not appear to be possible in the Standard Model. Instead, the fields produced are correlated at the scale $`T_i^1`$ which can produce only $`1`$ km correlations today. These are crude estimates, and galactic dynamos likely amplify these fields. The important point is that we can generate easily the $`100`$ kpc correlations observed today provided that the fields were initially of Hubble size correlation. Unless another mechanism for amplifying the correlations of magnetic fields is discovered, we suggest that, in order to obtain microgauss fields with $`100`$ kpc correlation lengths, helical fields must be generated with Hubble-scale correlations near or slightly after the QCD phase transition $`T_{\mathrm{QCD}}`$. The same conclusion regarding the relevance of the QCD scale for this problem was also reached in . The rest of this work presents a mechanism that can provide the desired Hubble size fields, justifying the estimate (3). We shall explain the mechanism and give simple estimates here. See for details. ## Magnetic field generation mechanism. The key players in our mechanism are domain walls formed at the QCD phase transition that possess an internal structure of QCD scale. We shall present a full exposition of these walls in but to be concrete, we shall discuss an axion wall similar to that described by Huang and Sikivie . We start with a similar effective Lagrangian to that used by Huang and Sikivie except we included the effects of the $`\eta ^{}`$ singlet field which they neglected: $$_{\mathrm{eff}}=\frac{f_a^2}{2}\left\|_\mu e^{i\stackrel{~}{a}}\right\|^2+\frac{f_\pi ^2}{4}\mathrm{Tr}\left\|_\mu 𝐔\right\|^2V(𝐔,\stackrel{~}{a})$$ (4) where $`\stackrel{~}{a}=f_a^1a`$ is the dimensionless axion field and the matrix $`𝐔=\mathrm{exp}(i\stackrel{~}{\eta }^{}+i\stackrel{~}{\pi }^f\lambda ^f)`$ contains the pion and $`\eta ^{}`$ fields (to simplify the calculations, we consider only the $`SU(2)`$ flavor group). Although the $`\eta ^{}`$ field is not light, it couples to the anomaly and is the dominant player in aligning the magnetic fields. The potential $$V=\frac{1}{2}\mathrm{Tr}(\mathrm{𝐌𝐔}e^{i\stackrel{~}{a}}+\mathrm{h}.\mathrm{c}.)E\mathrm{cos}\left(\frac{i\mathrm{ln}(det(𝐔))}{N_c}\right)$$ (5) was first introduced in . It should be realized that $`i\mathrm{ln}(det(𝐔))i\mathrm{ln}(det(𝐔))+2\pi n`$ is a multivalued function and we must choose the minimum valued branch. Details about this potential are discussed in but several points will be made here. All dimensionful parameters are expressed in terms of the QCD chiral and gluon vacuum condensates, and are well known numerically: $`𝐌=\mathrm{diag}(m_q^i\|\overline{q}^iq^i\|)`$ and $`E=b\alpha _s/(32\pi )G^2`$. This potential correctly reproduces the Veneziano-Witten effective chiral Lagrangian in the large $`N_c`$ limit ; it reproduces the anomalous conformal and chiral Ward identities of QCD; and it reproduces the known dependence in $`\theta `$ for small angles . We should also remark that the qualitative results do not depend on the exact form of the potential: domain walls form naturally because of the discrete nature of the symmetries . The result is that two different types of axion domain walls form . One is almost identical to the one discussed in with small corrections due to the $`\eta ^{}`$. We shall call this the axion/pion ($`a_\pi `$) domain wall. The second type, which we shall call the axion/eta’ ($`a_\eta ^{}`$) domain wall is a new solution characterized by a transition in both the axion and $`\eta ^{}`$ fields. The boundary conditions (vacuum states) for this wall are $`\stackrel{~}{a}(\mathrm{})=\stackrel{~}{\eta }^{}(\mathrm{})=0`$ and $`\stackrel{~}{a}(\mathrm{})=\stackrel{~}{\eta }^{}(\mathrm{})=\pm \pi `$ with $`\pi ^0=0`$ at both boundaries. The main difference between the structures of the two walls is that, whereas the $`a_\pi `$ domain wall has structure only on the huge scale of $`m_a^1`$, the $`\eta ^{}`$ transition in $`a_\eta ^{}`$ has a scale of $`m_\eta ^{}^1\mathrm{\Lambda }_{\mathrm{QCD}}^1`$. The reason is that, in the presence of the non-zero axion ($`\theta `$) field, the pion becomes effectively massless due to its Goldstone nature. The $`\eta ^{}`$ is not sensitive to $`\theta `$ and so its mass never becomes zero. It is crucial that the walls have a structure of scale $`\mathrm{\Lambda }_{\mathrm{QCD}}^1`$: there is no way for the $`a_\pi `$ wall to trap nucleons because of the huge difference in scales but the $`a_\eta ^{}`$ wall has exactly this structure and can therefore efficiently align the nucleons. The model we propose is this: Immediately after the phase transition, the universe is filled with domain walls of scale $`T_{\mathrm{QCD}}^1`$. As the temperature drops, these domain walls coalesce, resulting in an average of one domain wall per Hubble volume with Hubble-scale correlations . It is these $`a_\eta ^{}`$ domain walls which align the dipole moments of the nucleons producing the seed fields. The following steps are crucial for this phenomenon: 1) The coalescing of QCD domain wall gives the fields $`\pi ^f`$, $`\eta ^{}`$ Hubble-scale correlations. 2) These fields interact with the nucleons producing Hubble-scale correlations of nucleon spins residing in the vicinity of the domain wall. (The spins align perpendicular to the wall surface.) 3) Finally, the nucleons, which carry electric and magnetic moments (due to strong CP violation), induce Hubble-scale correlated magnetic and electric fields. 4) These magnetic and electric fields eventually induce a nonzero helicity which has the same correlation. This helicity enables the inverse cascade. ## Quantitative Estimates. As outlined below, we have estimated the strengths of the induced fields in terms of the QCD parameters . We consider two types of interactions. First, the nucleons align with the domain wall. Here we assume that the fluctuations in the nucleon field $`\mathrm{\Psi }`$ are rapid and that these effects cancel leaving the classical domain wall background unaltered. Thus, we are able to estimate many mean values correlated on a large scale on the domain walls such as $`\overline{\mathrm{\Psi }}\gamma _5\sigma _{xy}\mathrm{\Psi }`$ and $`\overline{\mathrm{\Psi }}\gamma _z\gamma _5\mathrm{\Psi }`$ through the interaction $`\overline{\mathrm{\Psi }}(i\overline{)}m_Ne^{i\stackrel{~}{\eta }^{}(z)\gamma _5})\mathrm{\Psi }`$. To estimate the magnetization of the domain wall, we make the approximation that the wall is flat compared to the lengths scales of the nucleon interactions. By assuming that momentum is conserved in the wall, we reduce our problem to an effective $`1+1`$ dimensional theory (in $`z`$ and $`t`$) which allows us to compute easily various mean value using a bosonization trick . The result for the mean value $`\overline{\mathrm{\Psi }}\gamma _5\sigma _{xy}\mathrm{\Psi }`$ for example is : $$\overline{\mathrm{\Psi }}\sigma _{xy}\gamma _5\mathrm{\Psi }\frac{\mu }{\pi }\mathrm{\Lambda }_{\mathrm{QCD}}^2,$$ (6) where $`\mu m_N`$ is a dimensional parameter originating from the bosonization procedure of the corresponding 2D system and the parameter $`\mathrm{\Lambda }_{\mathrm{QCD}}^2dk_xdk_y`$ comes from counting the nucleon degeneracy in the $`x`$$`y`$ plane of a Fermi gas at temperature $`T_c\mathrm{\Lambda }_{\mathrm{QCD}}`$. These mean values are only nonzero within a distance $`\mathrm{\Lambda }_{\mathrm{QCD}}^1`$ of the domain wall and are correlated on the same Hubble-scale as the domain wall. From now on we treat the expectation value (6) as a background classical field correlated on the Hubble-scale. Once these sources are known, one could calculate the generated electromagnetic field by solving Maxwell’s equations with the interaction $$_{\mathrm{int}}=\frac{1}{2}(d_\mathrm{\Psi }\overline{\mathrm{\Psi }}\sigma _{\mu \nu }\gamma _5\mathrm{\Psi }+\mu _\mathrm{\Psi }\overline{\mathrm{\Psi }}i\sigma _{\mu \nu }\mathrm{\Psi })F_{\mu \nu }+\overline{\mathrm{\Psi }}(i\mathrm{D})^2\mathrm{\Psi }$$ (7) where $`d_\mathrm{\Psi }`$ ($`\mu _\mathrm{\Psi }`$) is effective electric (magnetic) dipole moments of the field $`\mathrm{\Psi }`$. Due to the CP violation (nonzero $`\theta `$) along the axion domain wall, the anomalous nucleon dipole moment in (7) $`d_\mathrm{\Psi }\mu _\mathrm{\Psi }\frac{e}{m_N}`$ is also nonzero . This is an important point: if no anomalous moments were induced, then only charged particles could generate the magnetic field: the walls would be diamagnetic not ferromagnetic as argued in and Landau levels would exactly cancel the field generated by the dipoles. Solving the complete set of Maxwell’s equations, however, is extremely difficult. Instead, we use simple dimensional arguments. For a small planar region of area $`\xi ^2`$ filled with aligned dipoles with constant density, we know that the net magnetic field is proportional to $`\xi ^1`$ since the dipole fields tend to cancel, thus for a flat section of our domain wall, the field would be suppressed by a factor of $`(\xi \mathrm{\Lambda }_{\mathrm{QCD}})^1`$. For a perfectly flat, infinite domain wall ($`\xi \mathrm{}`$), there would be no net field as pointed out in . However, our domain walls are far from flat. Indeed, they have many wiggles and high frequency modes, thus, the size of the flat regions where the fields are suppressed is governed by a correlation $`\xi `$ which describes the curvature of the wall. Thus, the average electric and magnetic fields produced by the domain wall are of the order $$F_{\mu \nu }\frac{1}{\xi \mathrm{\Lambda }_{\mathrm{QCD}}}\left[d_\mathrm{\Psi }\overline{\mathrm{\Psi }}\sigma _{\mu \nu }\gamma _5\mathrm{\Psi }+\mu _\mathrm{\Psi }\overline{\mathrm{\Psi }}i\sigma _{\mu \nu }\mathrm{\Psi }\right]$$ (8) where $`\xi `$ is an effective correlation length related to the size of the dominant high frequency modes. To estimate what effective scale $`\xi `$ has, however, requires an understanding of the dynamics of the domain walls. Initially, the domain walls are correlated with a scale of $`\mathrm{\Lambda }_{\mathrm{QCD}}^1`$. As the temperature cools, the walls smooth out and the lower bound $`\xi _{}(t)`$ for the scale of the walls correlations increases from $`\xi _{}(0)\mathrm{\Lambda }_{\mathrm{QCD}}^1`$. This increase is a dynamical feature, however, and is thus slow. In addition, the walls coalesce and become correlated on the Hubble-scale generating large scale correlations. Thus the wall has correlations from $`\xi _{}(t)`$ up to the upper limit set by the Hubble-scale. We expect that $`\xi `$ Hubble size at the time that the fields are aligned and that the suppression is not nearly as great as implied in . Note that, even though the effects are confined to the region close to the wall, the domain walls are moving and twisted so that the effects occur throughout the entire Hubble volume. The picture is thus that fields of strength $$E_zB_z\frac{1}{\xi \mathrm{\Lambda }_{\mathrm{QCD}}}\frac{e}{m_N}\frac{m_N\mathrm{\Lambda }_{\mathrm{QCD}}^2}{\pi }\frac{e\mathrm{\Lambda }_{\mathrm{QCD}}}{\xi \pi }$$ (9) are generated with short correlations $`\xi `$, but then domains are correlated on a large scale by the Hubble-scale modes of the coalescing domain walls. Thus, strong turbulence is generated with correlations that run from $`\mathrm{\Lambda }_{\mathrm{QCD}}`$ up to the Hubble-scale. Finally, we note that this turbulence should be highly helical. This helicity arises from the fact that both electric and magnetic fields are correlated together along the entire domain wall, $`\stackrel{}{𝐄}\stackrel{}{𝐀}/\tau `$ where $`\stackrel{}{𝐀}`$ is the vector potential and $`\tau `$ is a relevant timescale for the electric field to be screened (we expect $`\tau \mathrm{\Lambda }_{\mathrm{QCD}}^1`$ as we discuss below). The magnetic helicity density is thus $$h\stackrel{}{𝐀}\stackrel{}{𝐁}\tau E_zB_z\tau \frac{e^2}{\pi ^2}\frac{\mathrm{\Lambda }_{QCD}^2}{\xi ^2}.$$ (10) Note carefully what happens here: The total helicity was zero in the quark-gluon-plasma phase and remains zero in the whole universe, but the helicity is separated so that in one Hubble volume, the helicity has the same sign. The reason for this is that, as the domain walls coalesce, initial perturbations cause either a soliton or an antisoliton to dominate and fill one Hubble volume. In the neighboring volume, there will be other solitons and antisolitons so that there is an equal number of both, but they are spatially separated which prevents them from annihilating. This is similar to how a particle and antiparticle may be created and then separated so they do not annihilate. In any case, the helicity is a pseudoscalar and thus has the same sign along the domain wall: The entire Hubble volume has helicity of the same sign. This is the origin of the Hubble-scale correlations in the helicity and in $`B^2`$. The correlation parameter $`\xi `$ which affects the magnitude of the fields plays no role in disturbing this correlation. Eventually, the electric field will be screened. The time scale for this is set by the plasma frequency for the electrons (protons will screen much more slowly) $`\omega _p\mathrm{\Lambda }_{\mathrm{QCD}}`$. The nucleons, however, also align on a similar timescale $`\mathrm{\Lambda }_{QCD}^1`$, and the helicity is generated on this scale too, so the electric screening will not qualitatively affect the mechanism. Finally, we note that the turbulence requires a seed which remains in a local region for a timescale set by the conductivity $`\sigma cT/e^2\mathrm{\Lambda }_{\mathrm{QCD}}`$ where for $`T=100`$ MeV, $`c0.07`$ and is smaller for higher $`T`$. Thus, even if the domain walls move at the speed of light (due to vibrations), there is still time to generate turbulence. For this mechanism to work and not violate current observations, it seems that the domain walls must eventually decay. Several mechanisms have been discussed for the decay of axion domain walls and the timescales for these decays are much larger than $`\mathrm{\Lambda }_{\mathrm{QCD}}^1`$, ie. long enough to generate these fields but short enough to avoid cosmological problems. QCD domain walls are quasistable and may nicely solve this problem. We assume that some mechanism exists to resolve the domain wall problem in an appropriate timescale. Thus, all the relevant timescales are of the order $`\mathrm{\Lambda }_{\mathrm{QCD}}^1`$ except for the lifetime of the walls and thus, although the discussed interactions will affect the quantitative results, they will not affect the mechanism or substantially change the order of the effects. ## Conclusion. We have shown that this mechanism can generate the magnetic fields (3) with large correlations, though galactic dynamos should still play an important amplification role. It seems that the crucial conditions for the dynamo to take place are fields $`B>10^{20}`$ G with large ($`100`$ kpc) correlations. From (3) we see that we have a huge interval $`1\xi \mathrm{\Lambda }_{\mathrm{QCD}}10^{10}`$ of $`\xi `$ to seed these dynamos. Also, if $`\xi `$ is small, then this mechanism may generate measurable extra-galactic fields. We mention two new points that distinguish this mechanism from previous proposals . First, the key nucleon is the neutron which generates the fields due to an anomalous dipole moment induced by the CP violating domain walls. The nucleons thus make the wall ferromagnetic, not diamagnetic as discussed in . Second, the interaction between the domain walls and nucleons are substantial because of the QCD scale of the $`\eta ^{}`$ transition. There is no way that axion domain walls with scales $`m_a^1`$ can efficiently align nucleons at a temperature $`T_{\mathrm{QCD}}`$. We should also note that the magnitudes of the fields generated by this mechanism are small enough to satisfy the constraints placed by nucleosynthesis and CMB distortions. Thus, domain walls at the QCD phase transition, in particular those described in , provide a nice method of generating magnetic fields on $`100`$ kpc correlations today (3). This work was supported by the NSERC of Canada. We would like to thank R. Brandenberger for many useful discussions. AZ wishes to thank: M. Shaposhnikov and I. Tkachev for valuable discussions which motivated this study; Larry McLerran and D. Son for discussions on Silk damping; and M. Voloshin and A. Vainshtein for discussions on the magnetic properties of domain walls.
warning/0004/hep-ph0004148.html	ar5iv	text	# References CP nonconservation is a fundamental issue in particle physics. We know that CP is not conserved in $`K\overline{K}`$ mixing (i.e. $`ϵ0`$) and in $`K`$ decay (i.e. $`ϵ^{}0`$). However, only an upper limit ($`0.63\times 10^{25}e`$cm) exists for the electric dipole moment ($`edm`$) of the neutron. This may not be so bothersome until we realize that the currently accepted theory of strong interactions, i.e. quantum chromodynamics (QCD), actually violates CP through the instanton-induced term $$_\theta =\theta _{QCD}\frac{g_s^2}{64\pi ^2}ϵ_{\mu \nu \alpha \beta }G_a^{\mu \nu }G_a^{\alpha \beta },$$ (1) where $`g_s`$ is the strong coupling constant, and $$G_a^{\mu \nu }=^\mu G_a^\nu ^\nu G_a^\mu +g_sf_{abc}G_b^\mu G_c^\nu $$ (2) is the gluonic field strength. The value of $`\theta _{QCD}`$ must then be less than $`10^{10}`$ in magnitude (instead of the expected order of unity) to account for the nonobservation of the neutron $`edm`$. This is known as the strong CP problem. Another fundamental issue in particle physics is supersymmetry. It allows us to solve the hierarchy problem so that our effective theory at the electroweak energy scale ($`M_W`$) is protected against large radiative corrections. However, this requires the scale of soft supersymmetry breaking ($`M_{SUSY}`$) to be not much higher than $`M_W`$. There is no theoretical understanding of why the two scales must be related in this way. In the following we address the question of how the strong CP problem is to be solved in the context of supersymmetry. We find that it can be achieved with a new kind of axion which couples to the gluino rather than to quarks. In a natural implementation of this idea, we find that the breaking of supersymmetry must have the same origin as the axion. The scale of electroweak symmetry breaking is also related. This works because the so-called $`\mu `$ problem is being solved along the way. In our framework, the $`\theta _{QCD}`$ contribution to any quark $`edm`$ is canceled exactly by the minimization of the dynamical gluino phase. Hence the calculation of $`edm`$’s in the MSSM (Minimal Supersymmetric Standard Model) becomes unambiguous. The possibility of cancellation among other different CP nonconserving contributions to $`edm`$’s can now be pursued without fear of contradiction. With the addition of colored fermions, the parameter $`\theta _{QCD}`$ of Eq. (1) is replaced by $$\overline{\theta }=\theta _{QCD}ArgDetM_uM_d3ArgM_{\stackrel{~}{g}},$$ (3) where $`M_u`$ and $`M_d`$ are the respective mass matrices of the charge = 2/3 and charge = –1/3 quarks, and $`M_{\stackrel{~}{g}}`$ is the mass of the gluino. The famous Peccei-Quinn solution is to introduce a dynamical phase to the quark masses which then relaxes $`\overline{\theta }`$ to zero. The specific realization of this requires an axion which is ruled out experimentally. Two other axionic solutions have been proposed which are at present consistent with all observations. The DFSZ solution introduces a singlet scalar field as the source of the axion but its mixing with the doublet scalar fields which couple to the quarks is very much suppressed. The KSVZ solution introduces new heavy quarks so that the axion does not even couple directly to the usual quarks. Neither scheme requires supersymmetry. In the context of supersymmetry however, it is clear that the simplest and most natural thing to do is to attach the axion to the gluino rather than to the quarks in Eq. (3). Because of the structure of supersymmetry, all other superparticles will be similarly affected. This is then a very strong hint that it may have something to do with the breaking of supersymmetry. As shown below with our proposed singlet complex scalar field $`S`$, whose phase contains the axion, all soft supersymmetry breaking parameters are of order $`\|S\|^2/M_{Pl}`$, where $`M_{Pl}10^{19}`$ GeV is the Planck mass. Hence a value of $`10^{11}`$ GeV for $`\|S\|`$, which is allowed by astrophysical and cosmological constraints, would imply $`M_{SUSY}1`$ TeV. Since the electroweak symmetry breaking terms are also among this group, it does not require any stretch of the imagination to find $`M_W`$ and $`M_{SUSY}`$ to be only an order of magnitude apart. It is known that a continuous global $`U(1)_R`$ symmetry can be defined for the MSSM. The quark ($`\widehat{Q},\widehat{u}^c,\widehat{d}^c`$) and lepton ($`\widehat{L},\widehat{e}^c`$) chiral superfields have $`R=+1`$ whereas the Higgs ($`\widehat{H}_u,\widehat{H}_d`$) chiral superfields and the vector superfields have $`R=0`$. The superpotential $`\widehat{W}=\mu \widehat{H}_u\widehat{H}_d+h_u\widehat{Q}\widehat{H}_u\widehat{u}^c+h_d\widehat{Q}\widehat{H}_d\widehat{d}^c+h_e\widehat{L}\widehat{H}_d\widehat{e}^c`$ (4) has $`R=+2`$ except for the $`\mu `$ term (which has $`R=0`$). In the above, the Yukawa couplings $`h_{u,d,e}`$ are nonhermitian matrices in flavor space. The resulting Lagrangian is then invariant only with respect to the usual discrete $`R`$ parity, i.e. $$R(1)^{3B+L+2J},$$ (5) where $`B`$ is baryon number, $`L`$ lepton number, and $`J`$ spin angular momentum, hence $`R`$ is even for particles and odd for superparticles. We now propose to make $`U(1)_R`$ an exact global symmetry of the supersymmetric Lagrangian, as well as that of all the supersymmetric breaking terms. We introduce the composite operator $`\mu (\widehat{S}){\displaystyle \frac{1}{M_{Pl}}}\left(\widehat{S}\right)^2,`$ (6) where the singlet superfield $`\widehat{S}`$ has $`R=+1`$. Our model is then defined by the new superpotential $`\widehat{W}`$ $`=`$ $`\mu (\widehat{S})\widehat{H}_u\widehat{H}_d+m_s^2\mu (\widehat{S})`$ (7) $`+`$ $`h_u\widehat{Q}\widehat{H}_u\widehat{u}^c+h_d\widehat{Q}\widehat{H}_d\widehat{d}^c+h_e\widehat{L}\widehat{H}_d\widehat{e}^c,`$ which has $`R=+2`$, thus yielding a supersymmetric Lagrangian which is invariant under $`U(1)_R`$, together with the following set of supersymmetry breaking terms which are also invariant under $`U(1)_R`$: $`\mathrm{\Delta }`$ $`=`$ $`\|\mu (S)\|^2\left[\stackrel{~}{Q}^{}Y_Q\stackrel{~}{Q}+\stackrel{~}{u^c}^{}Y_u\stackrel{~}{u^c}+\stackrel{~}{d^c}^{}Y_d\stackrel{~}{d^c}+\stackrel{~}{L}^{}Y_L\stackrel{~}{L}+\stackrel{~}{e^c}^{}Y_e\stackrel{~}{e^c}\right]`$ (8) $`+`$ $`\{\mu (S)^{}[k_u\stackrel{~}{Q}H_u\stackrel{~}{u^c}+k_d\stackrel{~}{Q}H_d\stackrel{~}{d^c}+k_e\stackrel{~}{L}H_d\stackrel{~}{e^c}]+h.c.\}`$ $`+`$ $`\|\mu (S)\|^2[y_u\|H_u\|^2+y_d\|H_d\|^2+(k_\mu H_uH_d+h.c.)]`$ $`+`$ $`\{\mu (S)^{}[k_3\stackrel{~}{\lambda }_3^a\stackrel{~}{\lambda }_3^a+k_2\stackrel{~}{\lambda }_2^i\stackrel{~}{\lambda }_2^i+k_1\stackrel{~}{\lambda }_1\stackrel{~}{\lambda }_1]+h.c.\},`$ where $`\stackrel{~}{\lambda }_3^a`$ is the gluino octet, $`\stackrel{~}{\lambda }_2^i`$ the $`SU(2)_L`$ gaugino triplet, and $`\stackrel{~}{\lambda }_1`$ the $`U(1)_Y`$ gaugino singlet. The parameters $`k_{1,2,3}`$ and $`k_\mu `$ are complex, whereas $`y_{u,d}`$ are real. The matrices $`Y_{Q,L}`$ and $`Y_{u,d,e}`$ are hermitian, whereas $`k_{u,d,e}`$ are nonhermitian. Obviously, we have assumed in the above that the source of all supersymmetry breaking terms is $`\mu (S)`$. Together with $`U(1)_R`$, this solves the so-called $`\mu `$ problem in the MSSM, because the scale of $`\mu (S)`$ is $`\|S\|^2/M_{Pl}`$ which is of order 1 TeV for $`\|S\|10^{11}`$ GeV, instead of the typical grand unification scale of $`10^{16}`$ GeV. The Lagrangian $`\mathrm{\Delta }`$ describes the interaction of $`\mu (S)`$ with sfermions, Higgs doublets and gauginos only. However, a complete description of our model requires the self interactions of the singlet to be specified as well. The pure singlet contribution $`m_s^2\mu (\widehat{S})`$ in $`\widehat{W}`$ is allowed by the symmetries of the model and the mass parameter $`m_s^2`$ is a priori arbitrary. Through the $`F`$-term contributions, this induces a positive mass-squared parameter for the singlet: $`m_F^24m_s^4/M_{Pl}^2`$. However, interactions at higher energies at or near the Planck scale can provide an additional mass-squared parameter $`m_0^2`$ as well as a quartic coupling $`\lambda _s`$. Hence the effective potential for the singlet takes the form $`V_s=M_s^2\|S\|^2+\lambda _s\|S\|^4`$ with $`M_s^2m_0^2+m_F^2`$. Since the Higgs doublets have vanishing $`R`$ charges, the electroweak breaking cannot have any effect on the fate of $`U(1)_R`$. The only way to break it is to allow the singlet to develop a nonvanishing vacuum expectation value. This can happen only when $`M_s^2<0`$ so that $`v_s^2=M_s^2/2\lambda _s`$. Since $`m_F^2`$ is positive, $`m_0^2`$ should be negative enough to induce a negative $`M_s^2`$. This impies that $`m_s^2`$ cannot be as large as $`M_{Pl}^2`$ as it would leave $`U(1)_R`$ unbroken; hence $`\|m_s^2\|\|m_0^2\|v_s^2`$ is a natural choice. The singlet field could then be expanded around $`v_s`$ as $`S(x)={\displaystyle \frac{1}{\sqrt{2}}}[v_s+s(x)]e^{i\phi (x)},`$ (9) where $`\phi (x)`$ is the corresponding Nambu–Goldstone boson which has a strictly flat potential, and $`s(x)`$ is a real scalar field with a mass of order $`v_s`$. It is clear from the above that our $`U(1)_R`$ plays the role of what is usually called $`U(1)_{PQ}`$. Whereas the conventional $`U(1)_{PQ}`$ applies to the usual quarks and leptons, our $`U(1)_{PQ}`$ applies only to the superparticles, and the gluino is the only colored fermion having a nonvanishing $`U(1)_{PQ}`$ current: $`J_\mu ^{5,\stackrel{~}{g}}=\overline{\lambda ^a}\gamma _\mu \gamma _5\lambda ^a,`$ (10) where we have used the four-component notation: $`\lambda ^a=(\stackrel{~}{\lambda }_3^a,\overline{\stackrel{~}{\lambda }_3^a})`$. Now the gluino also contributes to the color current with respect to which $`J_\mu ^{5,\stackrel{~}{g}}`$ has a nonvanishing quantum anomaly: $`^\mu J_\mu ^{5,\stackrel{~}{g}}={\displaystyle \frac{6g_s^2}{64\pi ^2}}ϵ_{\mu \nu \alpha \beta }G_a^{\mu \nu }G_a^{\alpha \beta }.`$ (11) Since $`J_\mu ^5`$ couples to $`\phi `$ as $`^\mu \phi J_\mu ^5`$, the effective QCD vacuum angle takes the form $`\overline{\theta }=\theta _{QCD}+6\phi (x),`$ (12) where the nondynamical phases in the quark mass matrices and the phase of the complex constant $`k_3`$ can be included in $`\theta _{QCD}`$ by a chiral rotation. In close analogy with the KSVZ scenario, our $`\phi (x)`$ also receives a potential from the instanton background so as to develop a vacuum expectation value which enforces $`\overline{\theta }0`$ \[i.e. $`\phi =\theta _{QCD}/6`$\], to all orders in perturbation theory. Rather than the quarks, it is thus the gluino which realizes the Peccei-Quinn mechanism of solving the strong CP problem. The axion, $`av_s[\phi (x)\phi ]`$, has a mass and lifetime given by $`m_am_\pi {\displaystyle \frac{f_\pi }{f_a}},\tau (a2\gamma )\left({\displaystyle \frac{m_\pi }{m_a}}\right)^5\tau (\pi 2\gamma ),`$ (13) where its decay constant $`f_a`$ is equal to $`v_s/6`$. Our axion is not a DFSZ axion as it does not couple to quarks and leptons; it is also not a KSVZ axion as it does not couple to unknown colored fermion multiplets beyond the MSSM spectrum. We may call it the $`gluino`$ $`axion`$ as it is induced by promoting the masses of the gauginos to local operators. Let us choose $`v_s/\sqrt{2}10^{11}`$ GeV, which is in the middle of the range of $`10^9`$ to $`10^{12}`$ GeV allowed by astrophysical and cosmological bounds on $`f_a`$. The effective theory below $`v_s`$ is then a replica of the MSSM with the effective $`\mu `$ parameter $`\mu _{eff}={\displaystyle \frac{v_s^2}{2M_{Pl}}}e^{i\theta _{QCD}/3}10^3\mathrm{GeV}\times e^{i\theta _{QCD}/3},`$ (14) which is the right scale for supersymmetry breaking. This seesaw mechanism for the $`\mu `$ parameter results from the introduction of the composite operator given by Eq. (6) into the theory, the dynamics of which are presumably dictated by physics at or near the Planck scale. On the other hand, the scale of the $`\mu `$ parameter is fixed by the astrophysical and cosmological bounds on the axion decay constant $`f_a`$ which gives (or receives) a meaning to (from) the intermediate scale $`v_s`$ <sup>1</sup><sup>1</sup>1If we abandon the composite operator idea (which correlates the axion scale with $`M_{SUSY}`$), we can still get a phenomenologically acceptable model, though less economical than the present one, as follows: Let $`S_2`$ have $`R=+2`$ which couples as $`\mu `$. Introduce $`S_1`$ with $`R=+1`$ which contains the axion. Let $`S_21`$ TeV, but $`S_110^{11}`$ GeV. The mixing between $`S_1`$ and $`S_2`$ is assumed small, so the axion coupling to the gluino is suppressed, i.e. a kind of DFSZ model applied to gluinos.. The low-energy effective theory is the softly broken MSSM with $`R`$ parity conservation. Indeed, after replacing the effective $`\mu `$ parameter \[Eq. (14)\] and $`\phi =\theta _{QCD}/6`$ into the effective Lagrangian \[Eq. (8)\], we obtain $`_{MSSM}^{soft}`$ $`=`$ $`\stackrel{~}{Q}^{}M_Q^2\stackrel{~}{Q}+\stackrel{~}{u^c}^{}M_{u^c}^2\stackrel{~}{u^c}+\stackrel{~}{d^c}^{}M_{d^c}^2\stackrel{~}{d^c}+\stackrel{~}{L}^{}M_L^2\stackrel{~}{L}+\stackrel{~}{e^c}^{}M_{e^c}^2\stackrel{~}{e^c}`$ (15) $`+`$ $`\{A_u\stackrel{~}{Q}H_u\stackrel{~}{u^c}+A_d\stackrel{~}{Q}H_d\stackrel{~}{d^c}+A_e\stackrel{~}{L}H_d\stackrel{~}{e^c}]+h.c.\}`$ $`+`$ $`M_{H_u}^2\|H_u\|^2+M_{H_d}^2\|H_d\|^2+(\mu _{eff}BH_uH_d+h.c.)`$ $`+`$ $`\{M_3\stackrel{~}{\lambda }_3^a\stackrel{~}{\lambda }_3^a+M_2\stackrel{~}{\lambda }_2^i\stackrel{~}{\lambda }_2^i+M_1\stackrel{~}{\lambda }_1\stackrel{~}{\lambda }_1+h.c.\},`$ which is nothing but the soft supersymmetry-breaking part of the MSSM Lagrangian. It is in fact this part of the Lagrangian that possesses all sources of CP violation through the complex $`A`$ parameters, the gaugino masses, and $`\mu _{eff}`$ itself. The explicit expressions for the mass parameters in Eq. (15) read as follows. The gaugino masses are given by $`M_3=\|k_3\|\mu _{eff}^{},M_2=k_2\mu _{eff}^{},M_1=k_1\mu _{eff}^{},`$ (16) which are not necessarily universal. The soft masses for the Higgs sector are given by $`M_{H_u}^2=y_u\|\mu _{eff}^2\|,M_{H_d}^2=y_d\|\mu _{eff}^2\|,\mu _{eff}B=\|\mu _{eff}\|^2(8{\displaystyle \frac{m_s^2}{v_s^2}}+k_\mu ),`$ (17) and are responsible for electroweak symmetry breaking, with similar expressions for the mass-squared matrices of the sfermion fields. In particular, since $`m_s^2/v_s^2`$ is of order unity, the $`B`$ parameter is also of the same scale. Finally the $`A`$ parameters are given by $`A_u=\mu _{eff}^{}k_u,A_d=\mu _{eff}^{}k_d,A_e=\mu _{eff}^{}k_e,`$ (18) which do not have to be proportional to $`h_u`$, $`h_d`$, and $`h_e`$ of Eq. (7) as in the constrained MSSM. As noted before, all mass scales of the MSSM Lagrangian \[Eq. (15)\] are fixed in terms of $`\|\mu _{eff}\|`$. More than this, the phase of $`\mu _{eff}`$, i.e. $`\theta _{QCD}/3`$, contributes universally to all mass parameters which are complex. However, the phases of the gaugino masses as well as those of the $`A`$ and $`B`$ terms also depend on the $`k`$ parameters. Hence if the flavor structure of these matrices is not the same as those of the usual quarks and leptons, then the CP violation in flavor-changing processes is a powerful probe into this sector of the effective theory. In the calculation of electric dipole moments due to supersymmetry, these CP phases can be considered as they are without worrying about whether there is an additional contribution from $`\overline{\theta }`$. In conclusion, we have presented in the above a simultaneous solution to two hierarchy problems, i.e. why $`\overline{\theta }`$ is so small (the strong CP problem) and why $`\mu `$ is 1 TeV and not $`10^{16}`$ GeV (the $`\mu `$ problem), as well as the related issue of why $`M_W`$ and $`M_{SUSY}`$ are only one order of magnitude apart. The primary difference between our approach and previous other attempts lies in the fact that the gaugino masses are promoted here to local operators given by $`\mu (\widehat{S})`$. Indeed, finite bare mass terms for the gauginos would automatically break the $`U(1)_{PQ}`$ symmetry, making it impossible for the relaxation of $`\overline{\theta }`$ to zero. As it is, $`S`$ serves two important purposes: its magnitude determines the scale of supersymmetry breaking and its phase solves the strong CP problem. Let us summarize our proposal. ($`i`$) We work in the framework of supersymmetry and identify $`U(1)_{PQ}`$ as $`U(1)_R`$ which contains the usual $`R`$ parity as a discrete subgroup. ($`ii`$) We require the supersymmetric Lagrangian and all supersymmetry breaking terms to be invariant under $`U(1)_R`$. ($`iii`$) We implement this with the composite operator $`\mu (\widehat{S})(\widehat{S})^2/M_{Pl}`$ where $`\widehat{S}`$ is a singlet superfield having $`R=+1`$. ($`iv`$) The spontaneous breaking of $`U(1)_R`$ generates an axion and relaxes the effective QCD vacuum angle $`\overline{\theta }`$ to zero, using the dynamical gluino phase, thus solving the strong CP problem. ($`v`$) The existing astrophysical and cosmological bounds on the axion decay constant implies a supersymmetry breaking scale of 1 TeV. ($`vi`$) The effective Lagrangian at low energy is that of the MSSM with $`R`$ parity conservation. All mass scales are of order 1 TeV, thus solving the $`\mu `$ problem and the related issue of why $`M_W`$ and $`M_{SUSY}`$ are only an order of magnitude apart. ($`vii`$) Since $`\overline{\theta }=0`$ in this consistent supersymmetric theory, electric dipole moments can be calculated unambiguously from the other explicit CP violating terms of the MSSM. D.A.D. acknowledges the hospitality of the UCR Physics Department, where this work was initiated. The research of E.M. was supported in part by the U. S. Department of Energy under Grant No. DE-FG03-94ER40837.
warning/0004/hep-ph0004186.html	ar5iv	text	# The QCD Coupling Constant 11footnote 1This work was supported in part by the Director, Office of Science, Office of High Energy and Nuclear Physics, Division of High Energy Physics of the U.S. Department of Energy under Contracts DE-AC03-76SF00098. ## 1 QCD AND ITS COUPLING Quantum chromodynamics (QCD) is a gauge field theory that describes the strong interactions of quarks and gluons . All experimental results to date are consistent with QCD predictions to within the experimental and theoretical errors. In this review, we discuss the current status of the extraction of the strong interaction coupling constant $`\alpha _s`$ from the experimental data. The QCD Lagrangian describing the interactions of quarks and gluons is $$L=\frac{1}{4}F_{\mu \nu }^aF^{a\mu \nu }+\underset{k}{}\overline{\psi }_k\left(i\text{/}Dm_k\right)\psi _k,$$ (1) where $$F_{\mu \nu }^a=_\mu A_\nu ^a_\nu A_\mu ^a+gf^{abc}A_\mu ^bA_\nu ^c$$ (2) is the gluon field strength tensor, $$D_\mu =_\mu igA_\mu ^aT^a$$ (3) is the gauge covariant derivative, and $`T^a`$ are the $`SU(3)`$ representation matrices normalized so that $`\mathrm{tr}T^aT^b=\delta ^{ab}/2`$, and the sum on $`k`$ is over the six different flavors ($`u,d,s,c,b,t`$) of quarks. At the classical level, the QCD Lagrangian depends on the six quark masses $`m_k`$, and the strong interaction coupling constant $`g`$, or equivalently, the strong fine-structure constant $`\alpha _s=g^2/4\pi `$. The quantum theory contains an additional parameter, the $`\theta `$-angle, that violates CP. The experimental limit on this parameter is $`\theta <10^9`$ , so we will set it to zero for the purposes of this article. One can evaluate QCD scattering amplitudes in powers of $`\alpha _s`$ using a Feynman diagram expansion. As is typical in a quantum field theory, loop graphs are divergent and need to be treated using a renormalization scheme. The most commonly used scheme is modified minimal subtraction ($`\overline{\mathrm{MS}}`$) , and we will use this scheme throughout. An important consequence of renormalization is that the parameters $`\alpha _s`$ and $`m_k`$ of the QCD Lagrangian depend in a calculable manner on the $`\overline{\mathrm{MS}}`$ subtraction-scale $`\mu `$. The $`\mu `$ dependence of $`\alpha _s`$ is described by the $`\beta `$-function, $$\mu \frac{d\alpha _s}{d\mu }=\beta \left(\alpha _s\left(\mu \right)\right).$$ (4) In perturbation theory, $$\beta \left(\alpha _s\right)=\beta _0\frac{\alpha _s^2}{2\pi }\beta _1\frac{\alpha _s^3}{\left(2\pi \right)^2}\beta _2\frac{\alpha _s^4}{\left(2\pi \right)^3}\mathrm{},$$ (5) where (for $`n_f`$ flavors of quarks) $`\beta _0=11{\displaystyle \frac{2}{3}}n_f,`$ (6) $`\beta _1=51{\displaystyle \frac{19}{3}}n_f,`$ (7) and the next two terms are also known . If $`\alpha _s`$ is small, the renormalization group equation Eq. (4) can be integrated using only the $`\beta _0`$ term to give $$\frac{1}{\alpha _s(\mu _1)}=\frac{1}{\alpha _s(\mu _1)}+\frac{\beta _0}{2\pi }\mathrm{ln}\frac{\mu _1}{\mu _2}.$$ (8) Since $`\beta _0>0`$ for $`n_f<16.5`$, $`\alpha _s(\mu )0`$ as $`\mu \mathrm{}`$. The vanishing of the QCD coupling for large values of $`\mu `$ is referred to as asymptotic freedom. One important consequence of asymptotic freedom, is that QCD processes at high energies can be reliably computed in a perturbation expansion in $`\alpha _s`$. A measurable quantity, such as the total cross-section for $`e^+e^{}\mathrm{hadrons}`$ at high energies can be computed as a function of the QCD coupling constant $`\alpha _s(\mu )`$ and the center of mass energy $`E_{\mathrm{CM}}`$, $$\sigma \left(e^+e^{}\mathrm{hadrons}\right)=\frac{1}{E_{\mathrm{CM}}^2}f(\alpha _s\left(\mu \right),\mathrm{ln}\frac{E_{\mathrm{CM}}}{\mu }).$$ (9) where $`f`$ is a dimensionless function of its arguments. The form of the cross-section given in Eq. (9) follows from dimensional analysis: $`\sigma `$ has dimensions of $`\mathrm{energy}^2`$, and $`\mu `$ has dimensions of energy. In the $`\overline{\mathrm{MS}}`$ sscheme, any dependence on $`\mu `$ is logarithmic, so $`f`$ can only depend on $`\mathrm{ln}E_{\mathrm{CM}}/\mu `$. The cross section $`\sigma \left(e^+e^{}\mathrm{hadrons}\right)`$ is a measurable quantity and cannot depend on the subtraction scale $`\mu `$, so the $`\mu `$ dependence on the right hand side of Eq. (9) must cancel, $$\mu \frac{d}{d\mu }f(\alpha _s\left(\mu \right),\frac{E_{\mathrm{CM}}}{\mu })=0,$$ (10) and any value of $`\mu `$ can be used on the right hand side of Eq. (9). In practice, one can only compute the right hand side of Eq. (9) at some finite order in perturbation theory, and the approximate value of $`f`$ can depend on $`\mu `$ at higher order in perturbation theory. Typically, one finds that the perturbation expansion has terms of the form $$\left[\alpha _s(\mu )\mathrm{ln}\frac{E_{\mathrm{CM}}}{\mu }\right]^n,$$ (11) which are referred to as “leading logarithms.” Even if $`\alpha _s(\mu )`$ is small, the perturbation expansion can break down if $`\mathrm{ln}E_{\mathrm{CM}}/\mu `$ is large. For this reason, it is conventional to choose the subtraction scale of order the center of mass energy $`E_{\mathrm{CM}}`$. The exact choice of scale (for example, whether $`\mu =2E_{\mathrm{CM}}`$ or $`E_{\mathrm{CM}}`$ or $`E_{\mathrm{CM}}/2`$) is arbitrary, and differences in choice of scale are formally of higher order in $`\alpha _s`$. Many methods have been proposed to determine the optimum scale to use for a given calculation . The only way to determine the “best” scale at a given order is to compute the cross-section at next order. \[Of course, in this case, one might as well use the more accurate formula to determine the cross-section.\] The scale dependence of a given quantity can also be used to estimate the size of neglected higher order corrections. Scale dependence is a dominant source of error in many of the quantities that will be used to determine $`\alpha _s`$. Perturbation theory is valid if one chooses $`\mu `$ to be of order $`E_{\mathrm{CM}}`$, so that the expansion parameter is $`\alpha _s(E_{\mathrm{CM}})`$, with no large logarithms. This shows that at high-energies, the coupling constant is small because of asymptotic freedom, and QCD cross-sections are approximately those of free quarks and gluons. At low, energies, non-perturbative effects become important. The value of $`\alpha _s`$ is determined by computing a quantity in terms of $`\alpha _s`$, and comparing with its measured value. One might think that it is better to use high-energy processes to determine $`\alpha _s`$, since perturbation theory is more reliable. This is not necessarily the case. High energy processes can be computed more reliably precisely because they do not depend very much on $`\alpha _s`$. This means that errors in the experimental measurement or theoretical calculation get amplified when they are converted to an error on $`\alpha _s`$. We will see in this article that low-energy extractions of $`\alpha _s`$ have comparable errors to those at high energy. In addition to scale dependence, the coupling constant $`\alpha _s`$ is also subtraction scheme dependent. The scheme dependence of the coupling constant is compensated for by the scheme dependence in the functional form for a measurable quantity, so that the value of an observable is scheme independent. The $`\overline{\mathrm{MS}}`$ scheme will be used in this article, but there is still some residual scheme dependence we need to consider. In the $`\overline{\mathrm{MS}}`$ scheme, heavy quarks do not decouple in loop graphs at low-energy. For example, the $`\beta `$-function coefficient $`\beta _0=112n_f/3`$, where $`n_f`$ is the number of quark flavors. This expression is true for all energies, irrespective of the mass of the quark. One might expect that at energies much smaller than the mass $`m_Q`$ of a heavy quark, the quark does not contribute to the $`\beta `$-function. This cannot happen in the $`\overline{\mathrm{MS}}`$ scheme, since $`\overline{\mathrm{MS}}`$ is a mass-independent subtraction scheme, which results in a mass-independent $`\beta `$-function. What happens is that there are large logarithms of the form $`\mathrm{ln}m_Q/\mu `$ that compensate for the “incorrect” $`\beta `$-function at low-energies. In practice, one deals with this problem using an effective field theory. At energies smaller than $`m_Q`$, one switches from a QCD Lagrangian with $`n_f`$ flavors to a QCD Lagrangian with $`n_f1`$ flavors by integrating out the heavy quark flavor. The effect of the heavy quark is taken into account by higher dimension operators in the QCD Lagrangian, and by shifts in the parameters $`\alpha _s`$ and $`m_k`$. Thus it is necessary to specify the effective theory when quoting the value of $`\alpha _s`$. We will use the notation $`\alpha _s^{(n_f)}`$ to denote the value of $`\alpha _s`$ in the $`n_f`$ flavor theory. One can compute the relation between $`\alpha _s^{(n_f)}`$ and $`\alpha _s^{(n_f1)}`$ at the scale $`\mu =m_Q`$ of the heavy quark. This relation is known to three loops , $`\alpha _s^{(n_f1)}\left(m_Q\right)`$ $`=`$ $`\alpha _s^{(n_f)}\left(m_Q\right)[1+0.1528\left({\displaystyle \frac{\alpha _s^{(n_f)}\left(m_Q\right)}{\pi }}\right)^2`$ (12) $`+(0.97210.0847n_f)\left({\displaystyle \frac{\alpha _s^{(n_f)}\left(m_Q\right)}{\pi }}\right)^3+\mathrm{}].`$ When quoting the value of $`\alpha _s(\mu )`$, it is also necessary to specify the number of flavors in the effective theory. In most of the $`\alpha _s`$ determinations we consider, the appropriate effective theory to use is one with $`n_f=5`$, and $`\alpha _s`$ will refer to $`\alpha _s^{(5)}`$ unless otherwise specified. Classical QCD is a scale invariant theory, but this scale invariance is broken at the quantum level. The quantum theory has a dimensionful parameter $`\mathrm{\Lambda }`$ that characterizes the scale of the strong interactions. The $`\mathrm{\Lambda }`$ parameter is determined in terms of $`\alpha _s(\mu )`$. The solution of the renormalization group equation Eq. (4) including the first three terms in the $`\beta `$-function is $`\alpha _s(\mu )`$ $`=`$ $`{\displaystyle \frac{4\pi }{\beta _0\mathrm{ln}\left(\mu ^2/\mathrm{\Lambda }^2\right)}}[1{\displaystyle \frac{2\beta _1}{\beta _0^2}}{\displaystyle \frac{\mathrm{ln}\left[\mathrm{ln}\left(\mu ^2/\mathrm{\Lambda }^2\right)\right]}{\mathrm{ln}\left(\mu ^2/\mathrm{\Lambda }^2\right)}}`$ (13) $`+{\displaystyle \frac{4\beta _1^2}{\beta _0^4\mathrm{ln}^2\left(\mu ^2/\mathrm{\Lambda }^2\right)}}((\mathrm{ln}\left[\mathrm{ln}(\mu ^2/\mathrm{\Lambda }^2)\right]{\displaystyle \frac{1}{2}})^2+{\displaystyle \frac{\beta _2\beta _0}{\beta _1^2}}{\displaystyle \frac{5}{4}})].`$ This equation can be used to determine $`\mathrm{\Lambda }`$ if $`\alpha _s`$ is known at some scale $`\mu `$. The value of $`\mathrm{\Lambda }`$ depends on the number of terms retained in Eq. (13). The expansion parameter in Eq. (13) is $$\frac{\mathrm{ln}\left[\mathrm{ln}\left(\mu ^2/\mathrm{\Lambda }^2\right)\right]}{\mathrm{ln}\left(\mu ^2/\mathrm{\Lambda }^2\right)},$$ (14) which is small as long as $`\mu \mathrm{\Lambda }`$. In QCD, $`\mathrm{\Lambda }`$ is of order 200 MeV. The last term in Eq. (13) is often dropped in the definition of $`\mathrm{\Lambda }`$. For a fixed value of $`\alpha _s(M_Z)`$, the shift in $`\mathrm{\Lambda }`$ is approximately 15 MeV if the last term in Eq. (13) is dropped. The QCD $`\beta `$-function depends on $`n_f`$, and so changes across quark thresholds. This in turn implies that $`\mathrm{\Lambda }`$ changes across quark thresholds, so that $`\mathrm{\Lambda }^{(n_f)}`$ is the value of $`\mathrm{\Lambda }`$ with $`n_f`$ dynamical quark flavors. The matching conditions for $`\mathrm{\Lambda }^{(n_f)}\mathrm{\Lambda }^{(n_f1)}`$ can be computed using the matching condition Eq. (12) for $`\alpha _s`$. The differences between $`\mathrm{\Lambda }^{(3)}`$, $`\mathrm{\Lambda }^{(4)}`$, and $`\mathrm{\Lambda }^{(5)}`$ are numerically very significant. In addition to the perturbative effects discussed so far, non-perturbative effects play an important role in strong interaction processes. The size of non-perturbative effects is governed by the ratio of the strong interaction scale $`\mathrm{\Lambda }`$ to the typical energy $`E_{\mathrm{CM}}`$ of a given process. In many cases, non-perturbative effects are estimated using a model analysis, or by a phenomenological fit to the experimental data. A few processes can be analyzed rigorously using the operator product expansion (OPE) ; in this case, one can calculate the size of non-perturbative effects in terms of matrix elements of gauge invariant local operators. Two classic examples of this type are deep inelastic scattering, and the total cross-section for $`e^+e^{}\mathrm{hadrons}`$. It is convenient to calculate $`R(s)`$, the ratio of the total cross-sections for $`e^+e^{}\mathrm{hadrons}`$ and $`e^+e^{}\mu ^+\mu ^{}`$ at center of mass-energy $`E_{\mathrm{CM}}=\sqrt{s}`$. One can show that $$R(s)=f_0\left(\alpha _s(s)\right)+f_1\left(\alpha _s(s)\right)\frac{F_{\mu \nu }F^{\mu \nu }}{s^2}+\mathrm{},$$ (15) where the first non-perturbative correction depends on the vacuum expectation value of the square of the gluon field-strength tensor. By dimensional analysis, this quantity is of order $`\mathrm{\Lambda }^4`$, so the non-perturbative corrections are of order $`\mathrm{\Lambda }^4/s^2`$. The existence of an OPE provides some crucial information on the size of non-perturbative corrections. For $`R(s)`$, we know that the corrections vanish at least as fast as $`\mathrm{\Lambda }^4/s^2`$ for large values of $`s`$, because $`F_{\mu \nu }F^{\mu \nu }`$ is the lowest dimension operator that can contribute in the OPE. Similarly, it is known in deep inelastic scattering that the first non-perturbative corrections arise from twist-four operators, and are of order $`\mathrm{\Lambda }^2/Q^2`$, where $`Q`$ is the momentum transfer. One can estimate the size of non-perturbative corrections for processes with an OPE, by estimating the value of operator matrix elements. In some cases, one is fortunate enough that the relevant matrix element can actually be determined from some other measurement, or computed from first principles. The size of non-perturbative corrections is much less certain if the process does not have an OPE. Non-perturbative effects can fall off like a fractional power of $`\mathrm{\Lambda }/s`$, or could have some more complicated dependence on $`s`$. Typically, one uses some model estimate of the non-perturbative corrections. Perturbative and non-perturbative corrections to scattering cross-sections are interrelated, because the QCD perturbation series is an asymptotic expansion, rather than a convergent expansion. A dimensionless quantity $`f(\alpha _s)`$ has an expansion of the form $$f(\alpha _s)=c_0+c_1\alpha _s+c_2\alpha _s^2+\mathrm{}.$$ (16) Typically, the coefficients $`c_n`$ grow as $`n!`$, so that the series has zero radius of convergence. The large-order behavior of the perturbation series can be computed in certain limiting cases . If one studies QCD in the limit of a large number of flavors, $`n_f\mathrm{}`$, with $`\alpha _sn_f`$ fixed, one can sum all terms of the form $`(\alpha _sn_f)^n`$. This is sometimes referred to as the “bubble chain” approximation, because the graphs one sums are of the type show in Fig. 1. QCD is not asymptotically free as $`n_f\mathrm{}`$. Nevertheless, one can try and apply the bubble chain results to QCD. The bubble chain graphs contribute to the QCD $`\beta `$-function. In the large $`n_f`$ limit, the coefficient $`\beta _0=112/3n_f2/3n_f`$. One therefore computes the bubble chain sum, makes the replacement $`n_f3\beta _0/2`$, and uses the resultant expression for QCD with $`\beta _0>0`$. This seemingly unjustified procedure has provided some useful insights into the nature of the QCD perturbation series. A detailed discussion of this method is beyond the scope of the present article. It is typically found that the coefficients $`c_n`$ in the perturbation expansion have a factorial divergence in the bubble chain approximation. One can try and sum a series of this type using a Borel transformation. One defines the Borel transform of $`f`$ by $$f_B(t)=c_0\delta (t)+c_1+c_2t+\mathrm{}+\frac{c_{n+1}}{n!}t^n+\mathrm{}.$$ (17) Then the original function can be obtained by the inverse Borel transform, $$f(\alpha _s)=_0^{\mathrm{}}𝑑tf_B(t)e^{t/\alpha _s}.$$ (18) Suppose the coefficients of $`f(\alpha _s)`$ have the form $$c_{n+1}=a^nn!,n>0.$$ (19) Then $$f_B(t)=c_0\delta (t)+\underset{n=0}{\overset{\mathrm{}}{}}a^nt^n=c_0\delta (t)+\frac{1}{1at},$$ (20) and the inverse Borel transform gives $$f(\alpha _s)=c_0+_0^{\mathrm{}}𝑑te^{t/\alpha _s}\frac{1}{1at}.$$ (21) The behavior of the integral is governed by the singularities in the complex $`t`$ plane, which are referred to as renormalons. If $`a<0`$, the integral is well-defined, since the singularity at $`t=1/a`$ is not along the path of integration. If $`a>0`$, the singularity is along the contour of integration. One can regulate the integral by deforming the contour around the singularity. The integral depends on the precise prescription used. The prescription dependence is related to the pole at $`t=1/a`$, which gives a contribution to the integral of the form $$\mathrm{exp}\left[\frac{1}{a\alpha _s}\right]=\mathrm{exp}\left[\frac{\beta _0\mathrm{ln}(\mu /\mathrm{\Lambda })}{2\pi a}\right]=\left(\frac{\mathrm{\Lambda }}{\mu }\right)^{\beta _0/2\pi a}.$$ (22) The value of $`\mu `$ is chosen to be the typical energy in the process, such as the momentum transfer $`Q`$ for deep inelastic scattering. The perturbative series then has the same structure as a typical non-perturbative correction, a power law correction of the form $$\left(\frac{\mathrm{\Lambda }^2}{Q^2}\right)^{u_0}$$ (23) where $`u_0=\beta _0/(4\pi a)`$ is the renormalon singularity in the variable $`u=\beta _0t/(4\pi )`$. \[It is conventional to refer to the location of the renormalon singularity in $`u`$ rather than in $`t`$.\] Contributions of the form Eq. (23) are called renormalon ambiguities, since their value depends on the way in which one performs the inverse Borel transform. Renormalon ambiguities have the same structure as non-perturbative corrections. It has been suggested in the literature that renormalons can be used as a guide to the size of non-perturbative corrections. There is one non-trivial check to this idea. A renormalon singularity at $`u_0`$ corresponds to a non-perturbative ambiguity of the form Eq. (23). In processes that have an OPE, all non-perturbative effects should be given by the matrix elements of gauge invariant local operators. To every renormalon ambiguity in the perturbation expansion, there should be a corresponding ambiguity in the operator matrix element, such that the sum is well-defined . This can only happen provided that there is a gauge invariant local operator corresponding to every renormalon singularity. For example, the first gauge invariant operator corrections to deep inelastic scattering are of the form $`\mathrm{\Lambda }^2/Q^2`$, and it is known that the first renormalon ambiguity is at $`u=1`$. Similarly, for $`R`$, the first non-perturbative corrections are of the form $`\mathrm{\Lambda }^4/Q^4`$, and the first renormalon ambiguity is at $`u=2`$. The matching between renormalon singularities and the OPE occurs in all examples that have been computed so far. For this reason, renormalon singularities have also been taken as an indication of the size of non-perturbative effects in processes without an OPE. Non-perturbative effects are expected to fall off faster if the renormalons are at larger values of $`u`$. Infrared sensitivity is also used to estimate the size of non-perturbative corrections to a measurable quantity . One computes the quantity in the presence of an infrared cutoff momentum $`\lambda `$. For example, one can imagine working in a box of size $`1/\lambda `$, or using a gluon mass of order $`\lambda `$. Cross-sections can have infrared divergences of the form $`\mathrm{ln}\lambda `$. If one computes measurable cross-sections for color singlet states to scatter into color singlet states, one finds that the $`\mathrm{ln}\lambda `$ terms cancel, and the cross-section is infrared finite, a result known as the KLN theorem. It is important to include finite detector resolution to get a finite cross-section, as for QED. For example, the Bhabha scattering cross-section for $`e^+e^{}e^+e^{}`$ has an infrared divergence at one-loop order. However, it is impossible to distinguish $`e^+e^{}e^+e^{}`$ from $`e^+e^{}e^+e^{}\gamma `$ if the photon energy $`E_\gamma `$ is smaller than the detector resolution $`\delta `$. The measurable quantity is the sum of the $`e^+e^{}e^+e^{}`$ cross-section and the $`e^+e^{}e^+e^{}\gamma `$ cross-section for $`E_\gamma <\delta `$, which is free of infrared singularities. While the $`\mathrm{ln}\lambda `$ term must cancel, terms of order $`\lambda `$, $`\lambda ^2\mathrm{ln}\lambda `$, etc. which vanish as $`\lambda 0`$ need not cancel. The first non-vanishing term is an indication of the infrared-sensitivity of a given quantity. In QCD, one can imagine that the scale $`\lambda `$ represents the confinement scale $`\mathrm{\Lambda }`$. A process that is infrared sensitive at order $`\lambda ^n`$ would then be expected to have non-perturbative corrections of order $`\mathrm{\Lambda }^n/Q^n`$, where $`Q`$ is the typical momentum transfer. In a few cases, one can analyze the problem using renormalon methods, and by using the criterion of infrared sensitivity. It is found in these cases that both methods give the same estimate for the size of non-perturbative corrections. ## 2 $`\alpha _s`$ FROM $`Z`$ DECAYS AND $`e^+e^{}`$ TOTAL RATES The total cross section for $`e^+e^{}\text{hadrons}`$ is obtained (at low values of $`\sqrt{s}`$) by multiplying the muon-pair cross section by the factor $`R`$. At lowest order in QCD perturbation theory $`R=R^0=3\mathrm{\Sigma }_qe_q^2`$ where $`e_q`$ is the electric charge of the quark of flavor $`q`$. The higher-order QCD corrections to this are known, and the results can be expressed in terms of the factor: $$R=R^{(0)}\left[1+\frac{\alpha _s}{\pi }+C_2\left(\frac{\alpha _s}{\pi }\right)^2+C_3\left(\frac{\alpha _s}{\pi }\right)^3+\mathrm{}\right]$$ (24) where $`C_2=1.411`$ and $`C_3=12.8`$ . This result is only correct in the zero-quark-mass limit. The $`𝒪(\alpha _s`$) corrections are also known for massive quarks . The principal advantages of determining $`\alpha _s`$ from $`R`$ in $`e^+e^{}`$ annihilation are that the measurement is inclusive, that there is no dependence on the details of the hadronic final state and that non-perturbative corrections are suppressed by $`1/s^2`$. A measurement by CLEO at $`\sqrt{s}=10.52`$ GeV yields $`\alpha _s(10.52\text{ GeV})=0.20\pm 0.01\pm 0.06`$ which corresponds to $`\alpha _s(M_Z)=0.13\pm 0.005\pm 0.03.`$ A comparison of the theoretical prediction of Eqn. 24 (corrected for the $`b`$-quark mass), with all the available data at values of $`\sqrt{s}`$ between 20 and 65 GeV, gives $`\alpha _s(35\text{ GeV})=0.146\pm 0.030.`$ It should be noted that the size of the order $`\alpha _s^3`$ term is of order 40% of that of the order $`\alpha _s^2`$ and 3% of the order $`\alpha _s`$. If the order $`\alpha _s^3`$ term is not included, the extracted value decreases to $`\alpha _s(35\text{ GeV})=0.142\pm 0.03`$, a difference smaller than the experimental error. Measurements of the ratio of the hadronic to leptonic width of the $`Z`$ at LEP and SLC, $`\mathrm{\Gamma }_h/\mathrm{\Gamma }_\mu `$ probe the same quantity as $`R`$. Using the average of $`\mathrm{\Gamma }_h/\mathrm{\Gamma }_\mu =20.783\pm 0.029`$ gives $`\alpha _s(M_Z)=0.123\pm 0.004`$ . The prediction depends upon the couplings of the quarks and leptons to the $`Z`$. The precision is such that higher order electroweak corrections to these couplings must be included. There are theoretical errors arising from the values of top-quark and Higgs masses which enter in these radiative corrections. Hence, while this method has small theoretical uncertainties from QCD itself, it relies sensitively on the electroweak couplings of the $`Z`$ to quarks and on the ability of the Standard Model of electroweak interactions to predict these correctly. The presence of new physics which changes these couplings via electroweak radiative corrections would invalidate the extracted value of $`\alpha _s(M_Z)`$. Since the Standard Model fits the measured $`Z`$ properties well, this concern is ameliorated and more precise value of $`\alpha _s`$ can be obtained by using a global fit to the many precisely measured properties of the $`Z`$ boson and the measured $`W`$ and top masses. This gives $$\alpha _s(M_Z)=0.1192\pm 0.0028$$ This error is larger than the shift in the value of $`\alpha _s(M_Z)`$ ($`0.002`$) that would result if the order $`\alpha _s(M_Z)^3`$ term were omitted and hence one can conclude that it is very unlikely that the uncertainty due to the unknown $`\alpha _s(M_Z)^4`$ terms will dominate over the experimental uncertainty. ## 3 DETERMINATION OF $`\alpha _s`$ FROM DEEP INELASTIC SCATTERING The original and still one of the most powerful quantitative tests of perturbative QCD is the breaking of Bjorken scaling in deep-inelastic lepton-hadron scattering. Consider the case of electron-proton scattering ($`epeX`$), where the cross-section can be written as $$\frac{d\sigma }{dxdy}=\frac{4\pi \alpha _{em}^2s}{Q^4}\left[\frac{1+(1y)^2}{2}2xF_1(x,Q^2)+(1y)(F_2(x,Q^2)2xF_1(x,Q^2))\right]$$ (25) The variables are defined as follows (see Figure 2): $`q`$ is the momentum of the exchanged photon, $`P`$ is the momentum of the target proton, $`k`$ is that of the incoming electron, and $`Q^2`$ $`=q^2`$ $`\nu `$ $`={\displaystyle \frac{qP}{m_p}}`$ $`x`$ $`={\displaystyle \frac{Q^2}{2m_p\nu }}`$ $`y`$ $`={\displaystyle \frac{qp}{kp}}`$ $`s`$ $`=2pk+m_p^2`$ For charged current scattering, which proceeds via the exchange of a virtual $`W`$ boson between the lepton and target nucleus, there is an additional parity violating structure function $`F_3`$ $$\begin{array}{cc}\frac{d\sigma ^{\nu N}}{dxdy}=\hfill & \frac{G_F^2M_W^4s}{2\pi (Q^2+M_W^2)^2}(xy^2F_1^{\nu N}(x,Q^2)\hfill \\ & +(1yx^2y^2M^2/Q^2)F_2^{\nu N}(x,Q^2)\hfill \\ & \frac{1}{2}x((1y)^21)F_3^{\nu N}(x,Q^2))\hfill \end{array}$$ (27) For $`\overline{\nu }N`$ scattering the sign of the last ($`xF_3`$) term is reversed. In the leading-logarithm approximation, the measured structure functions $`F_i(x,Q^2`$) are related to the quark distribution functions $`q_i(x,Q^2)`$ according to the naive parton model, for example $$F_2(x,Q^2)=\underset{i}{}e_i^2q_i(x,Q^2)$$ (28) Here $`q_i(x,Q^2)`$ is the probability for a parton of type $`i`$ to carry a fraction $`x`$ of the nucleon’s momentum. The $`Q^2`$ dependence of the parton distribution functions is predicted by perturbative QCD, hence a measurement of the $`Q^2`$ dependence (“scaling violation”) can by used to measure $`\alpha _s`$. In describing the way in which scaling is broken in QCD, it is convenient to define nonsinglet and singlet quark distributions: $$F^{NS}=q_iq_jF^S=\underset{i}{}(q_i+\overline{q}_i)$$ The nonsinglet structure functions have nonzero values of flavor quantum numbers such as isospin or baryon number. The variation with $`Q^2`$ of these is described by the so-called DGLAP equations : $$\begin{array}{cc}Q^2\frac{F^{NS}}{Q^2}\hfill & =\frac{\alpha _s(\|Q\|)}{2\pi }P^{qq}F^{NS}\hfill \\ & \\ Q^2\frac{}{Q^2}\left(\genfrac{}{}{0pt}{}{F^S}{G}\right)\hfill & =\frac{\alpha _s(\|Q\|)}{2\pi }\left(\genfrac{}{}{0pt}{}{P^{qq}}{P^{gq}}\genfrac{}{}{0pt}{}{2n_fP^{qg}}{P^{gg}}\right)\left(\genfrac{}{}{0pt}{}{F^S}{G}\right)\hfill \end{array}$$ where $``$ denotes a convolution integral: $$fg=_x^1\frac{dy}{y}f(y)g\left(\frac{x}{y}\right)$$ The leading-order Altarelli-Parisi splitting functions are $$\begin{array}{cc}P^{qq}\hfill & =\frac{4}{3}\left[\frac{1+x^2}{(1x)_+}\right]+2\delta (1x)\hfill \\ P^{qg}\hfill & =\frac{1}{2}\left[x^2+(1x)^2\right]\hfill \\ P^{gq}\hfill & =\frac{4}{3}\left[\frac{1+(1x)^2}{x}\right]\hfill \\ P^{gg}\hfill & =6\left[\frac{1x}{x}+x(1x)+\frac{x}{(1x)_+}+\frac{11}{12}\delta (1x)\right]\hfill \\ & \frac{n_f}{3}\delta (1x)\hfill \end{array}$$ Here the gluon distribution $`G(x,Q^2)`$ has been introduced and $`1/(1x)_+`$ means $$_0^1𝑑x\frac{f(x)}{(1x)_+}=_0^1𝑑x\frac{f(x)f(1)}{(1x)}$$ Measurement of the structure functions over a large range of $`x`$ and $`Q^2`$ allows both $`\alpha _s`$ and the parton distributions to be determined. Notice that $`\alpha _s`$ and the gluon distribution can only be obtained by measuring the $`Q^2`$ dependence. The precision of contemporary experimental data demands that higher-order corrections also be included . The above results are for massless quarks. Algorithms exist for the inclusion of nonzero quark masses . These are particularly important for neutrino scattering near the charm threshold. At low $`Q^2`$ values, there are also important “higher-twist” (HT) contributions of the form: $$F_i(x,Q^2)=F_i^{(LT)}(x,Q^2)+\frac{F_i^{(HT)}(x,Q^2)}{Q^2}+\mathrm{}$$ Leading twist (LT) terms are those whose behavior can be predicted using the parton model, and are related to the parton distribution functions. Higher-twist corrections depend on matrix elements of higher dimension operators. These corrections are numerically important only for $`Q^2<𝒪`$(few GeV<sup>2</sup>) except for $`x`$ very close to 1. At very large values of $`x`$ corrections proportional to $`\mathrm{log}(1x)`$ can become important . From Eqn. LABEL:AltarelliParisi;a, it is clear that a nonsinglet structure function offers in principle the most precise test of the theory, since the $`Q^2`$ evolution is independent of the unmeasured gluon distribution. The CCFR collaboration fit to the Gross-Llewellyn Smith sum rule which is known to order $`\alpha _s^3`$ (Estimates of the order $`\alpha _s^4`$ term are available ) $$\begin{array}{cc}& _0^1𝑑x\left\{F_3^{\overline{\nu }p}(x,Q^2)+F_3^{\nu p}(x,Q^2)\right\}=\hfill \\ & 3\left[1\frac{\alpha _s}{\pi }(1+3.58\frac{\alpha _s}{\pi }+19.0(\frac{\alpha _s}{\pi })^2)\right]\mathrm{\Delta }HT\hfill \end{array}$$ where the higher-twist contribution $`\mathrm{\Delta }HT`$ is estimated to be $`(0.09\pm 0.045)/Q^2`$ in and to be somewhat smaller by . The CCFR collaboration , combines their data with that from other experiments and gives $`\alpha _s(\sqrt{3}\text{GeV})=0.28\pm 0.035(\text{expt.})\pm 0.05(\text{sys})_{0.03}^{+0.035}(\text{theory})`$. The error from higher-twist terms (assumed to be $`\mathrm{\Delta }HT=0.05\pm 0.05`$) dominates the theoretical error. If the higher twist result of is used, the central value increases to 0.31 in agreement with the fit of . This value extrapolates to $`\alpha _s(M_Z)=0.118\pm 0.011`$. Measurements involving singlet-dominated structure functions, such as $`F_2`$, result in correlated measurements of $`\alpha _s`$ and the gluon structure function. A full next to leading order fit combining data from SLAC , BCDMS , E665 and HERA has been performed . These authors extend the analysis to next to next to leading order (NNLO). In this case the full theoretical calculation is not available as not all the three-loop anomalous dimensions are known; their analysis uses moments of structure functions<sup>2</sup><sup>2</sup>2The moments are defined by $`M_n=_0^1x^nF(x,Q^2)𝑑x`$. and is restricted to those moments where the full calculation is available . The NNLO result is $`\alpha _s(M_Z)=0.1172\pm 0.0017(\text{expt.})\pm 0.0017(\text{sys})`$. Here the first error is a combination of statistical and systematic experimental errors, and the second error is due to the uncertainties, quark masses, higher twist and target mass corrections, and errors from the gluon distribution. If only a next to leading order fit is performed then the value decreases to $`\alpha _s(M_Z)=0.116`$ indicating that the theoretical results are stable. No error is included from the choice of $`\mu `$; $`\mu =Q`$ is assumed. We use a total error of $`\pm 0.0045`$ to take into account an estimate of the scale uncertainty. This result is consistent with earlier determinations , , and . The spin-dependent structure functions, measured in polarized lepton nucleon scattering, can also be used to determine $`\alpha _s`$. The spin structure functions $`G_1`$ and $`G_2`$ are defined in terms of the asymmetry in polarized lepton nucleon scattering $$a(x,y)=\frac{d\sigma _p^{eN}}{dxdy}\frac{d\sigma _{ap}^{eN}}{dxdy}$$ (29) where the subscript $`p`$ ($`ap`$) refers to the state where the nucleon spin is parallel (anti-parallel) to its direction of motion in the center of mass frame of the lepton-nucleon system. In both cases the lepton has its spin aligned along its direction of motion. $$\begin{array}{cc}a(x,y)=\hfill & \frac{8\pi \alpha _{em}^2y}{MQ^2}((12/y^2+2x^2y^2M^2/Q^2)G_1(x,Q^2)\hfill \\ & +4x^2M^2G_2(x,Q^2)/Q^2)\hfill \end{array}$$ (30) The $`Q^2`$ evolution of the spin structure functions $`G_1(x,Q^2)`$ and $`G_2(x,Q^2)`$ is similar to that of the unpolarized ones and is known at next to leading order .Here the values of $`Q^22.5\text{ GeV}^2`$ are small particularly for the E143 data and higher-twist corrections are important. A fit using the measured spin dependent structure functions measured by themselves and by other experiments gives $`\alpha _s(M_Z)=0.121\pm 0.002(\text{expt.})\pm 0.006(\text{theory and syst.})`$. Data from HERMES are not included in this fit; they are consistent with the older data. $`\alpha _s`$ can also be determined from the Bjorken sum rule . $$S_{\mathrm{Bj}}=_0^1𝑑x\left(G_1^pG_1^n\right)=\frac{1}{6}a_3$$ (31) At lowest order in QCD $`a_3=g_A=\frac{G_V}{G_A}=1.2573\pm 0.0028`$. A fit gives $`\alpha _s(M_Z)=0.118_{0.024}^{+0.010}`$; a significant contribution to the error being due to the extrapolation into the (unmeasured) small $`x`$ region. Theoretically, the sum rule is preferable as the perturbative QCD result is known to higher order, and these terms are important at the low $`Q^2`$ involved. It has been shown that the theoretical errors associated with the choice of scale are considerably reduced by the use of Pade approximants which results in $`\alpha _s(1.7\text{GeV})=0.328\pm 0.03(\text{expt.})\pm 0.025(\text{theory})`$ corresponding to $`\alpha _s(M_Z)=0.116_{0.005}^{+0.003}(\text{expt.})\pm 0.003(\text{theory})`$. No error is included from the extrapolation into the region of $`x`$ that is unmeasured. Should data become available at smaller values of $`x`$ so that this extrapolation could be more tightly constrained, the sum rule method could provide a better determination of $`\alpha _s`$ than that from the spin structure functions themselves. At very small values of $`x`$ and $`Q^2`$, both the $`x`$ and $`Q^2`$ dependence of the structure functions is predicted by perturbative QCD . Here terms to all orders in $`\alpha _s\mathrm{ln}(1/x)`$ are summed. The data from HERA on $`F_2^{ep}(x,Q^2)`$ can be fitted to this form , including the NLO terms which are required to fix the $`Q^2`$ scale. The data are dominated by $`4\text{ GeV}^2<Q^2<100\text{ GeV}^2`$. The fit using H1 data gives $`\alpha _s(M_Z)=0.122\pm 0.004(\text{expt.})\pm 0.009(\text{theory})`$. (The theoretical error is taken from .) The dominant part of the theoretical error is from the scale dependence; errors from terms that are suppressed by $`1/\mathrm{log}(1/x)`$ in the quark sector are included while those from the gluon sector are not. Typically, $`\mathrm{\Lambda }`$ is extracted from the deep inelastic scattering data by parameterizing the parton densities in a simple analytic way at some $`Q_0^2`$, evolving to higher $`Q^2`$ using the next-to-leading-order evolution equations, and fitting globally to the measured structure functions. Thus, an important by-product of such studies is the extraction of parton densities at a fixed-reference value of $`Q_0^2`$. These can then be evolved in $`Q^2`$ and used as input for phenomenological studies in hadron-hadron collisions (see below). These densities will have errors associated with the that value of $`\alpha _s`$. A next-to-leading order fit must be used if the process being calculated is known to next-to-leading order in QCD perturbation theory. In such a case, there is an additional scheme dependence; this scheme dependence is reflected in the $`𝒪(\alpha _s)`$ corrections that appear in the relations between the structure functions and the quark distribution functions. There are two common schemes: a deep-inelastic scheme where there are no order $`\alpha _s`$ corrections in the formula for $`F_2(x,Q^2)`$ and the minimal subtraction scheme. It is important when these next-to-leading order fits are used in other processes (see below), that the same scheme is used in the calculation of the partonic rates. Most current sets of parton distributions are obtained using fits to all relevant data . In particular, data from purely hadronic initial states are used as they can provide important constraints on the gluon distributions. ### 3.0.1 PHOTON STRUCTURE FUNCTIONS Experiments in $`e^+e^{}`$ collisions can be used to study photon-photon interactions and to measure the structure function of a photon , by selecting events of the type $`e^+e^{}e^+e^{}+\mathrm{hadrons}`$ which proceeds via two photon scattering. If events are selected where one of the photons is almost on mass shell and the other has a large invariant mass $`Q`$, then the latter probes the photon structure function at scale $`Q`$; the process is analogous to deep inelastic scattering where a highly virtual photon is used to probe the proton structure. The $`Q^2`$ variation of this structure function follows that shown above (see Eq LABEL:AltarelliParisi;a). A review of the data can be found in . Data have become available from LEP and from TRISTAN which extend the range of $`Q^2`$ to of order 300 GeV<sup>2</sup> and $`x`$ as low as $`2\times 10^3`$and show $`Q^2`$ dependence of the structure function that is consistent with QCD expectations. Experiments at HERA can also probe the photon structure function by looking at jet production in $`\gamma p`$ collisions; this is analogous to the jet production in hadron-hadron collisions which is sensitive to hadron structure functions. The data are consistent with theoretical models . ## 4 $`\alpha _s`$ FROM FRAGMENTATION FUNCTIONS Measurements of the fragmentation function $`d_i(z,E)`$, the probability that a hadron of type $`i`$ be produced with energy $`zE`$ in $`e^+e^{}`$ collisions at $`\sqrt{s}=2E`$, can be used to determine $`\alpha _s`$. As in the case of scaling violations in structure functions, QCD predicts only the $`E`$ dependence in a form similar to the $`Q^2`$ dependence of Eq LABEL:AltarelliParisi;a. Hence, measurements at different energies are needed to extract a value of $`\alpha _s`$. Because the QCD evolution mixes the fragmentation functions for each quark flavor with the gluon fragmentation function, it is necessary to determine each of these before $`\alpha _s`$ can be extracted. The ALEPH collaboration has used data in the energy range $`\sqrt{s}=22`$ GeV to $`\sqrt{s}=91`$ GeV. A flavor tag is used to discriminate between different quark species, and the longitudinal and transverse cross sections are used to extract the gluon fragmentation function . The result obtained is $`\alpha _s(M_Z)=0.126\pm 0.007(\text{expt.})\pm 0.006(\text{theory})`$ . The theory error is due mainly to the choice of scale at which $`\alpha _s`$ is evaluated. The OPAL collaboration has also extracted the separate fragmentation functions. DELPHI has performed a similar analysis using data from other experiments at center of mass energies between 14 and 91 GeV with the result $`\alpha _s(M_Z)=0.124\pm 0.007\pm 0.009(\text{theory})`$. The larger theoretical error is because the value of $`\mu `$ was allowed to vary between $`0.5\sqrt{s}`$ and $`2\sqrt{s}`$. These results can be combined to give $`\alpha _s(M_Z)=0.125\pm 0.005\pm 0.008(\text{theory})`$. ## 5 $`\alpha _s`$ FROM EVENT SHAPES AND JET COUNTING An alternative method of determining $`\alpha _s`$ in $`e^+e^{}`$ annihilation involves measuring the the topology of the hadronic final states. There are many possible choices of inclusive event shape variables: thrust , energy-energy correlations , average jet mass, etc.. These quantities must be infrared safe, which means that they are insensitive to the low energy properties of QCD and can therefore be reliably calculated in perturbation theory. For example, the thrust distribution is defined by $$T=max(\underset{i}{}\left\|\overline{p_i}\overline{n}\right\|/\underset{i}{}\left\|\overline{p_i}\right\|),$$ (32) where the sum runs over all hadrons in the final state and the unit vector $`\overline{n}`$ is varied. At lowest order in QCD the process $`e^+e^{}q\overline{q}`$ results in a final state with back to back quarks i.e. “pencil-like” event with $`T=1`$. Alternatively, the event can be divided by a plane normal to the thrust axis and the invariant mass of the particles in the two hemispheres is computed, the larger (smaller) of these is $`M_h`$ ($`M_l`$). At lowest order in QCD $`M_h=M_l=0`$. The observed final state consists of hadrons rather than the quarks and gluons of perturbation theory. The hadronization of the partonic final state has an energy scale of order $`\mathrm{\Lambda }`$. The resulting hadrons acquire momentum components perpendicular to the original quark direction of order $`\mathrm{\Lambda }`$. This effect induces corrections to the shape variables of order $`\mathrm{\Lambda }/\sqrt{s}`$. A model is needed to describe the detailed evolution of a partonic final state into one involving hadrons, so that detector corrections can be applied. Furthermore if the QCD matrix elements are combined with a parton-fragmentation model, this model can then be used to correct the data for a direct comparison with the perturbative QCD calculation. The different hadronization models that are used model the dynamics that are controlled by non-perturbative QCD effects which we cannot yet calculate. The fragmentation parameters of these Monte Carlo simulations are tuned to get agreement with the observed data. The differences between these models can be used to estimate systematic errors. In addition to using a shape variable, one can perform a jet counting experiment. At order $`\alpha _s`$ the partonic final state $`q\overline{q}g`$ appears which can manifest itself as a three-jet final state after hadronization. Every higher order produce a higher jet multiplicity and measuring quantities that are sensitive to the relative rates of two-, three-, and four-jet events can lead to a determination of $`\alpha _s`$. There are theoretical ambiguities in the way that particles are combined to form jets. Quarks and gluons are massless, whereas the observed hadrons are not, so that the massive jets that result from combining them cannot be compared directly to the massless jets of perturbative QCD. The jet-counting algorithm, originally introduced by the JADE collaboration , has been used by many other groups. Here, particles of momenta $`p_i`$ and $`p_j`$ are combined into a pseudo-particle of momentum $`p_i+p_j`$ if the invariant mass of the pair is less than $`y_0\sqrt{s}`$. The process is then iterated until no more pairs of particles or pseudo-particles remain. The remaining number of pseudo-particles is then defined to be the number of jets in the event, and can be compared to the perturbative QCD prediction which depends on $`y_0`$. The Durham algorithm is slightly different: in computing the mass of a pair of partons, it uses $`M^2=2\text{min}(E_1^2,E_2^2)(1\mathrm{cos}\theta _{ij})`$ for partons of energies $`E_i`$ and $`E_j`$ separated by angle $`\theta _{ij}`$ . Different recombination schemes have been tried, for example combining 3-momenta and then rescaling the energy of the cluster so that it remains massless. These varying schemes result in the same data giving slightly different values of $`\alpha _s`$. These differences can be used to estimate a systematic error. However, such an error may be conservative as it is not based on a systematic approximation. The starting point for all these quantities is the multijet cross section. For example, at order $`\alpha _s`$, for the process $`e^+e^{}qqg`$ $$\frac{1}{\sigma }\frac{d^2\sigma }{dx_1dx_2}=\frac{2\alpha _s}{3\pi }\frac{x_1^2+x_2^2}{(1x_1)(1x_2)}$$ where $`x_i=\frac{2E_i}{\sqrt{s}}`$are the center-of-mass energy fractions of the final-state (massless) quarks. The order $`\alpha _s^2`$ corrections to this process have been computed, as well as the 4-jet final states such as $`e^+e^{}qqgg`$ . A distribution in a “three-jet” variable, such as those listed above, is obtained by integrating this differential cross section over an appropriate phase space region for a fixed value of the variable. Thus $`<1T>\alpha _s`$, $`<M_h^2>/s\alpha _s`$ and $`<M_l^2>/s\alpha _s^2`$. The result of this integration depends explicitly on $`\alpha _s`$ but scale $`\mu `$ at which $`\alpha _s(\mu )`$ is to be evaluated is not clear. In the case of jet counting, the invariant mass of a typical jet (or $`\sqrt{sy_0}`$) is probably a more appropriate choice than the $`e^+e^{}`$ center-of-mass energy. While there is no justification for doing so, if the value of $`\mu `$ is allowed to float in the fit to the data, the fit improves and the data tend to prefer values of order $`\sqrt{s}/10`$ GeV for some variables ; the exact value depends on the variable that is fitted. Typically experiments assign a systematic error from the choice of $`\mu `$ by varying it by a factor of 2 around the value determined by the fit. The choice of this factor is arbitrary Estimates for the non-perturbative corrections to $`<1T>`$ have been made using an operator product expansion. $$<1T>=A\frac{\alpha _s(\mu )}{2\pi }+B(\frac{\alpha _s(\mu )}{2\pi })^2+C\frac{\alpha _0}{\sqrt{s}}$$ (33) where A and B known quantities , $`\mu `$ is the renormalization scale and $`\alpha _0`$ is the non-perturbative parameter (the matrix element of an appropriate operator) to be determined from experiment. Note that the corrections are only suppressed by $`\sqrt{s}`$. This provides an alternative to the use of hadronization models for estimating these non-perturbative corrections. The DELPHI collaboration uses data below the $`Z`$ mass from many experiments and Eq. 33 to determine $`\alpha _s(M_Z)=0.119\pm 0.006`$, the error being dominated by the choice of scale. The values of $`\alpha _s`$ and the non-perturbative parameter $`\alpha _0`$ are also determined by a fit to using the variable $`<M_h^2>/s`$. While the extracted values of $`\alpha _s(M_Z)`$ are consistent with each other, the values of $`\alpha _0`$ are not. The analysis is useful as one can directly determine the size of the $`1/E`$ corrections; they are approximately 20% (50%) of the perturbative result at $`\sqrt{s}=91(11)`$ GeV. Even at $`\sqrt{s}=91`$ GeV the omission of these perturbative terms will cause a shift on the extracted value of $`\alpha _s`$ of $`0.05`$ which is much larger than the quoted experimental errors. The perturbative QCD formulae can break down in special kinematical configurations. For example, the first term in Eq. 33 contains a term of the type $`\alpha _s\mathrm{ln}^2(1T)`$. The higher orders in the perturbation expansion contain terms of order $`\alpha _s^n\mathrm{ln}^m(1T)`$. For $`T1`$ (the region populated by 2-jet events), the perturbation expansion in $`\alpha _s`$ is unreliable. The terms with $`nm`$ can be summed to all orders in $`\alpha _s`$ . If the jet recombination methods are used, higher-order terms involve $`\alpha _s^n\mathrm{ln}^m(y_0)`$, these too can be resummed . The resummed results give better agreement with the data at large values of $`T`$. Some caution should be exercised in using these resummed results because of the possibility of overcounting; the showering Monte Carlos that are used for the fragmentation corrections also generate some of these leading-log corrections. Different schemes for combining the order $`\alpha _s^2`$ and the resummations are available . These different schemes result in shifts in $`\alpha _s(M_Z)`$ of order $`\pm 0.002`$. The use of the resummed results improves the agreement between the data and the theory. Studies on event shapes have been undertaken at lower energies at TRISTAN, PEP/PETRA, and CLEO. A combined result from various shape parameters by the TOPAZ collaboration gives $`\alpha _s(58\text{GeV})=0.125\pm 0.009`$, using the fixed order QCD result, and $`\alpha _s(58\text{GeV})=0.132\pm 0.008`$ (corresponding to $`\alpha _s(M_Z)=0.123\pm 0.007`$) where the error is dominated by scale and fragmentation uncertainties. The CLEO collaboration fits to the order $`\alpha _s^2`$ results for the two jet fraction at $`\sqrt{s}=10.53`$ GeV, and obtains $`\alpha _s(10.93)=0.164\pm 0.004(\text{expt.})\pm 0.014(\text{theory})`$ . The dominant systematic error arises from the choice of scale ($`\mu `$), and is determined from the range of $`\alpha _s`$ that results from fit with $`\mu =10.53`$ GeV, and a fit where $`\mu `$ is allowed to vary to get the lowest $`\chi ^2`$. The latter results in $`\mu =1.2`$ GeV. Since the quoted result corresponds to $`\alpha _s(1.2)=0.35`$, it is by no means clear that the perturbative QCD expression is reliable and the resulting error should, therefore, be treated with caution. A fit to many different variables as is done in the LEP/SLC analyses would give added confidence to the quoted error. Recently studies have been carried out at energies between $``$130 GeV and $``$189 GeV . These can be combined to give $`\alpha _s(130\text{GeV})=0.114\pm 0.008`$ and $`\alpha _s(189\text{GeV})=0.1104\pm 0.005`$. The dominant errors are theoretical and systematic and, as most of these are in common at the different energies, these data, those at the $`Z`$ resonance and lower energy provide very clear confirmation of the expected decrease in $`\alpha _s`$ as the energy is increased. A combined analysis of the data between 35 and 189 GeV using data from OPAL and JADE using a large set of shape variables shows excellent agreement with $`\alpha _s(M_Z)=0.1187_{0.0019}^{+0.0034}`$. A comparison of this result with those at the $`Z`$ resonance from SLD , OPAL , L3 , ALEPH , and DELPHI , indicates that they are all consistent with this value. The experimental errors are smaller than the theoretical ones arising from choice of scale $`\mu `$ and modeling of non-perturbative effects, which are common to all of the experiments. The SLD collaboration determines the allowed range of $`\mu `$ by allowing any value that is consistent with the fit. This leads to a larger error ($`0.0056`$) than that obtained by DELPHI who vary $`\mu `$ by a factor of 2 around the best fit value and obtain $`\pm 0.0008`$. We elect to use a more conservative average of $`\alpha _s(M_Z)=0.119\pm 0.005`$. At lowest order in $`\alpha _s`$, the $`epeX`$ scattering process produces a final state of (1+1) jets, one from the proton fragment and the other from the quark knocked out by the underlying process $`e+\text{quark}e+\text{quark}`$. At next order in $`\alpha _s`$, a gluon can be radiated, and hence a (2+1) jet final state produced. By comparing the rates for these (1+1) and (2+1) jet processes, a value of $`\alpha _s`$ can be obtained. A NLO QCD calculation is available . The basic methodology is similar to that used in the jet counting experiments in $`e^+e^{}`$ annihilation discussed above. Unlike those measurements, the ones in $`ep`$ scattering are not at a fixed value of $`Q^2`$. In addition to the systematic errors associated with the jet definitions, there are additional ones since the structure functions enter into the rate calculations. Results from H1 and ZEUS can be combined to give $`\alpha _s(M_Z)=0.118\pm 0.0015(\text{stat.})\pm 0.009(\text{syst.})`$. The contributions to the systematic errors from experimental effects (mainly the hadronic energy scale of the calorimeter) are comparable to the theoretical ones arising from scale choice, structure functions, and jet definitions. The theoretical errors are common to the two measurements; therefore, we have not reduced the systematic error after forming the average. ## 6 $`\alpha _s`$ FROM $`\tau `$ DECAY The coupling constant $`\alpha _s`$ can be determined from an analysis of hadronic $`\tau `$ decays . The quantity that will be used is the ratio $$R_\tau =\frac{\mathrm{\Gamma }(\tau \nu _\tau +\mathrm{hadrons}+(\gamma ))}{\mathrm{\Gamma }(\tau \nu _\tau e\overline{\nu }_e+(\gamma ))},$$ (34) where $`(\gamma )`$ represents possible electromagnetic radiation, or lepton pairs. In the absence of radiative corrections, the ratio $`R_\tau `$ is $$R_\tau =3\left(\left\|V_{ud}\right\|^2+\left\|V_{us}\right\|^2\right)3,$$ (35) where $`3`$ is the number of colors. The experimental value $`R_\tau =3.61\pm 0.05`$ is close to three, which is experimental evidence for the existence of three colors in QCD. The deviation of $`R_\tau `$ from three is used to extract $`\alpha _s`$. The weak decay Lagrangian for non-leptonic $`\tau `$ decay is $`L={\displaystyle \frac{4G_F}{\sqrt{2}}}C_\tau (\mu )\left[V_{ud}^{}\overline{\nu }_\tau \gamma ^\mu P_L\tau \overline{d}\gamma _\mu P_Lu+V_{us}^{}\overline{\nu }_\tau \gamma ^\mu P_L\tau \overline{s}\gamma _\mu P_Lu\right],`$ (36) where $`V_{us}`$ and $`V_{ud}`$ are the CKM mixing angles. The Lagrangian Eq. (36) is obtained at the scale $`\mu =M_W`$ by integrating out the $`W`$ boson to generate a local four-Fermion operator in the effective theory below $`M_W`$, and $`C_\tau =1`$ at $`\mu =M_W`$. The typical momentum transfer in $`\tau `$ decays is of order $`m_\tau `$, so it is necessary to scale the Lagrangian Eq. (36) from $`\mu =M_W`$ to $`\mu =m_\tau `$. Electromagnetic interactions renormalize the Lagrangian. At one-loop, the renormalization from graphs shown in Fig. (3) produce a multiplicative renormalization of the Lagrangian, and give $$C_\tau (m_\tau )=1+\frac{\alpha _{\mathrm{em}}}{\pi }\mathrm{ln}\frac{M_W}{m_\tau }1.009.$$ (37) The $`\tau `$ decay amplitude $`\tau \nu _\tau X`$, where $`X`$ is the final hadronic state, can be written as $$A=i\frac{4G_F}{\sqrt{2}}C_\tau (\mu )\overline{u}(p_\nu )\gamma ^\mu P_Lu(p_\tau )\left[V_{ud}^{}X\left\|\overline{d}\gamma _\mu P_Lu\right\|0+V_{us}^{}X\left\|\overline{s}\gamma _\mu P_Lu\right\|0\right].$$ (38) Squaring the amplitude, and computing the decay rate gives $`\mathrm{\Gamma }`$ $`=`$ $`4{\displaystyle \frac{G_F^2\left\|C_\tau \right\|^2\left\|V_{ud}\right\|^2}{m_\tau }}{\displaystyle \underset{X}{}}{\displaystyle \frac{d^3p_\nu }{(2\pi )^32E_\nu }\left[p_\nu ^\mu p_\tau ^\nu +p_\nu ^\nu p_\tau ^\mu p_\nu p_\tau g^{\mu \nu }iϵ^{\mu \nu }{}_{\alpha \beta }{}^{}p_{\nu }^{\alpha }p_\tau ^\beta \right]}`$ (39) $`\times (2\pi )^4\delta ^4(p_\tau p_\nu p_X)0\left\|\overline{u}\gamma _\nu P_Ld\right\|XX\left\|\overline{d}\gamma _\mu P_Lu\right\|0,`$ where we have retained only the $`V_{ud}`$ term for simplicity. The sum on $`X`$ is symbolic for the sum over all final states, including phase space factors. The $`\delta `$ function can be written as $$\delta ^4(p_\tau p_\nu p_X)=d^4q\delta ^4(p_\tau p_\nu q)\delta ^4(qp_X),$$ (40) and the sum on $`X`$ can be written as $$\underset{X}{}(2\pi )^4\delta ^4(qp_X)0\left\|\overline{u}\gamma _\nu P_Ld\right\|XX\left\|\overline{d}\gamma _\mu P_Lu\right\|0=W_{\nu \mu }^{(ud)}(q).$$ (41) The tensor $`W_{\mu \nu }`$ is related to another quantity $`\mathrm{\Pi }_{\mu \nu }`$, defined by $$\mathrm{\Pi }_{\mu \nu }^{(ud)}(q)=id^4qe^{iqx}0\|T(j_\nu ^{}(x)j_\mu (0)\|0,$$ (42) where $`j^\mu =\overline{d}\gamma _\mu P_Lu`$. Inserting a complete set of states in the time-ordered product, one finds that $$W_{\mu \nu }^{(ud)}=2\mathrm{Im}\mathrm{\Pi }_{\mu \nu }^{(ud)}$$ (43) The tensor $`\mathrm{\Pi }_{\mu \nu }`$ depends on the only variable, $`q`$, and must have the form $$\mathrm{\Pi }_{\mu \nu }^{(ud)}(q)=\left(q^2g_{\mu \nu }+q_\mu q_\nu \right)\mathrm{\Pi }_T^{(ud)}(q^2)+q_\mu q_\nu \mathrm{\Pi }_L^{(ud)}(q^2),$$ (44) by Lorentz invariance. The tensor $`W_{\mu \nu }`$ is then given by $$W_{\mu \nu }^{(ud)}(q)=\left(q^2g_{\mu \nu }+q_\mu q_\nu \right)\mathrm{\Omega }_T^{(ud)}(q^2)+q_\mu q_\nu \mathrm{\Omega }_L^{(ud)}(q^2),$$ (45) where $$\mathrm{\Omega }_L^{(ud)}=2\mathrm{Im}\mathrm{\Pi }_L^{(ud)},\mathrm{\Omega }_T^{(ud)}=2\mathrm{Im}\mathrm{\Pi }_T^{(ud)}.$$ (46) If the light quark mass difference $`m_dm_u`$ and $`m_sm_u`$ are neglected, the hadronic currents $`\overline{d}\gamma _uP_Lu`$ and $`\overline{s}\gamma _uP_Lu`$ are conserved. This implies that $`q^\mu \mathrm{\Pi }_{\mu \nu }=0`$, so that $`\mathrm{\Pi }_L(q^2)=0`$. Inserting Eq. (40) and Eq. (45) into Eq. (47) gives $`\mathrm{\Gamma }`$ $`=`$ $`2{\displaystyle \frac{G_F^2\left\|C_\tau \right\|^2}{m_\tau }}{\displaystyle d^4q\delta (p_\tau p_\nu q)\frac{d^3p_\nu }{(2\pi )^32E_\nu }}`$ (47) $`(m_\tau ^2q^2)\left[\mathrm{\Omega }_T(m_\tau ^2+2q^2)+\mathrm{\Omega }_Lm_\tau ^2\right],`$ where $$\mathrm{\Omega }_{T,L}=\left\|V_{ud}\right\|^2\mathrm{\Omega }_{T,L}^{(ud)}+\left\|V_{us}\right\|^2\mathrm{\Omega }_{T,L}^{(us)}$$ (48) and we added back the $`V_{us}`$ contribution. There is no interference term (at lowest order in the weak interactions), because the $`ud`$ and $`us`$ currents lead to final states with different flavor quantum numbers, and the strong interactions conserve flavor. The hadronic invariant mass distribution can then be written as $`{\displaystyle \frac{d\mathrm{\Gamma }}{ds}}`$ $`=`$ $`{\displaystyle \frac{G_F^2\left\|C_\tau \right\|^2}{8\pi ^2m_\tau ^3}}(m_\tau ^2s)^2\left[\mathrm{\Omega }_T(s)(m_\tau ^2+2s)+\mathrm{\Omega }_L(s)m_\tau ^2\right]`$ (49) The ratio of the hadronic to leptonic decay rate of the $`\tau `$ is given by $$R_\tau =6\pi \left\|C_\tau \right\|^2_0^{m_\tau ^2}\frac{ds}{m_\tau ^2}\left(1\frac{s}{m_\tau ^2}\right)^2\left[\mathrm{\Omega }_T(s)\left(1+\frac{2s}{m_\tau ^2}\right)+\mathrm{\Omega }_L(s)\right]$$ (50) The hadronic tensors $`\mathrm{\Pi }_{L,T}(s)`$ are analytic in the complex $`s`$ plane, except for a branch cut along the positive real axis. The discontinuity across the cut is $`\mathrm{\Omega }_{L,T}(s)`$, and is the cross-section for the currents to create hadrons. Clearly, the hadron production rate is sensitive to non-perturbative effects, and can not be computed reliably. Far away from the physical cut, there are no infrared singularities, and QCD perturbation theory is valid. One can rewrite the integral Eq. (50) as $$R_\tau =6\pi i_C\frac{ds}{m_\tau ^2}\left(1\frac{s}{m_\tau ^2}\right)^2\left[\mathrm{\Pi }_T(s)\left(1+\frac{2s}{m_\tau ^2}\right)+\mathrm{\Pi }_L(s)\right],$$ (51) where $`C=C_1`$ is the contour shown in Fig. (4). The difference of $`\mathrm{\Pi }_{L,T}`$ above and below the cut is $`\mathrm{\Omega }_{L,T}`$, so this gives back Eq. (50). Since $`\mathrm{\Pi }_{L,T}`$ have no singularities in the complex $`s`$ plane other than the branch cut, one can deform the contour $`C_1`$ to the contour $`C_2`$, and use Eq. (51) with $`C=C_2`$. The advantage of using Eq. (51) with $`C=C_2`$ rather than $`C=C_1`$ is that one needs to know $`\mathrm{\Pi }_{L,T}`$ far away from the cut for most of the integration contour. The contour approaches the cut at $`s=m_\tau ^2`$, but at this point, the integrand vanishes as $`(sm_\tau )^2`$, so the contribution of the region near $`s=m_\tau ^2`$ to the total integral is suppressed . $`\mathrm{\Pi }(s)`$ can be computed in perturbation theory using an operator product expansion, which is valid away from the physical cut. The perturbation theory result for $`\mathrm{\Pi }(s)`$ is then substituted in Eq. (51). In practice, the calculation can be simplified by using the perturbation theory value for $`\mathrm{\Omega }(s)`$ in Eq. (50). We have argued above that perturbation theory is not valid for $`\mathrm{\Omega }(s)`$. Nevertheless, using the perturbation theory value for $`\mathrm{\Omega }(s)`$ in Eq. (50) is justified because $`\mathrm{\Pi }(s)`$ in perturbation theory has the same analytic structure in QCD. Thus using the perturbation theory value of $`\mathrm{\Pi }(s)`$ in Eq. (51) is equivalent to using the perturbation theory value for $`\mathrm{\Omega }(s)`$ in Eq. (50), even though the perturbative computation of $`\mathrm{\Omega }(s)`$ is not valid. The OPE for $`\mathrm{\Pi }(s)`$ is closely related to that for $`e^+e^{}\mathrm{hadrons}`$, which depends on the time-ordered product of two electromagnetic currents. $$\mathrm{\Pi }(s)=c_i(\mu ,s,\alpha _s(\mu ))O_i(\mu ),$$ (52) where $`c_i`$ are the coefficient functions, and $`O_i`$ are the local operators. Since the contour $`C_2`$ is a circle of radius $`m_\tau ^2`$ in the complex $`s`$ plane, one expects that logarithms in the uncalculated higher order corrections are minimized if one chooses $`\mu =m_\tau `$. The leading order operator is the unit operator. In the limit that the light quark masses are neglected, $`\mathrm{\Pi }_L`$ vanishes, and we only need to compute $`\mathrm{\Pi }_T`$, giving $`2\pi \mathrm{\Omega }_T(s)`$ $`=`$ $`\left\|C_\tau \right\|^2\left(\left\|V_{ud}\right\|^2+\left\|V_{us}\right\|^2\right)`$ (53) $`\times \left[1+{\displaystyle \frac{\alpha _s(\sqrt{s})}{\pi }}+F_3\left({\displaystyle \frac{\alpha _s(\sqrt{s})}{\pi }}\right)^2+F_4\left({\displaystyle \frac{\alpha _s(\sqrt{s})}{\pi }}\right)^3+\mathrm{}\right].`$ The first coefficient is the well-known result that the ratio $`\sigma (e^+e^{}\mathrm{hadrons})/\sigma (e^+e^{}q\overline{q})`$ is $`3(1+\alpha _s(\sqrt{s})/\pi )`$. The next two coefficients are $$F_3=1.98570.1153n_f,F_4=6.63681.2001n_f0.0052n_f^2.$$ (54) Using the $`\beta `$-function to write $`\alpha _s(\sqrt{s})`$ in terms of $`\alpha _s(m_\tau )`$, evaluating the $`s`$ integral, and setting $`n_f=3`$ gives $`R_\tau `$ $`=`$ $`3\left\|C_\tau \right\|^2\left(\left\|V_{ud}\right\|^2+\left\|V_{us}\right\|^2\right)`$ (55) $`\times \left[1+{\displaystyle \frac{\alpha _s(m_\tau )}{\pi }}+5.2023\left({\displaystyle \frac{\alpha _s(m_\tau )}{\pi }}\right)^2+26.366\left({\displaystyle \frac{\alpha _s(m_\tau )}{\pi }}\right)^3+\mathrm{}\right]`$ Assuming $`\alpha _s(m_\tau )0.35`$, the series in brackets is $`1+0.111+0.065+0.036+\mathrm{}`$, so the terms are still numerically decreasing till order $`\alpha ^3`$. One can use the value of the last term as an estimate of the theoretical uncertainty in the perturbative value for the coefficient of the unit operator. This also tells us that we can neglect corrections from higher dimension operator that are smaller than about $`3\%`$. The $`1/m_\tau ^2`$ corrections to $`R_\tau `$ arise from the quark mass corrections to the coefficient of the unit operator in the OPE. At order $`m^2/m_\tau ^2`$, the currents are no longer conserved, so one needs to compute both $`\mathrm{\Pi }_L`$ and $`\mathrm{\Pi }_T`$. The only light quark mass contribution of any significance is the $`s`$-quark mass correction , $$\delta R_\tau =3\left\|V_{us}\right\|^2\left[8\frac{m_s^2}{m_\tau ^2}\left(1+\frac{16\alpha _s(m_\tau )}{3\pi }\right)\right].$$ (56) Using $`m_s150`$ MeV as an estimate for the $`s`$-quark mass gives $`\delta R_\tau 0.008`$, which is smaller than the error in the perturbation series. The $`1/m_\tau ^4`$ and corrections in the OPE arise from the dimension four operators $`F_{\mu \nu }F^{\mu \nu }`$ and $`m\overline{\psi }\psi `$, and the $`1/m_\tau ^6`$ corrections from the four-quark operators $`\overline{\psi }\mathrm{\Gamma }\psi \overline{\psi }\mathrm{\Gamma }\psi `$, where $`\mathrm{\Gamma }`$ is some combination of $`\gamma `$ matrices. An analysis of these corrections, based on model estimates of the operator matrix elements indicates that these corrections are smaller than the uncertainty in the perturbation series . The size of non-perturbative corrections can be determined directly from the experimental data. Instead of considering the integral Eq. (50) that gives the total hadronic width, one compares the integral of $`d\mathrm{\Gamma }/ds`$ weighted with $`(1s/m_\tau ^2)^k(s/m_\tau ^2)^l`$ with the corresponding moment of the experimental data. By studying the moments for different values of $`k`$ and $`l`$, one finds that the non-perturbative corrections are about 3% , and so are comparable in size to the uncertainty in the perturbation series. The experimentally measured quantity is $`R_\tau ^{(ud)}=3.484\pm 0.024`$ , the ratio for $`\tau `$ to decay into non-strange hadrons to the leptonic decay rate. This is given by Eq. (55), dropping $`V_{us}`$, and gives $$\alpha _s(m_\tau )=0.34\pm 0.03,$$ (57) where we have assumed a theoretical uncertainty of 100% in the $`\alpha ^3`$ term. This value corresponds to $$\alpha _s(m_Z)=0.119\pm 0.003,$$ (58) ## 7 $`\alpha _s`$ FROM LATTICE GAUGE THEORY COMPUTATIONS The strong coupling constant $`\alpha _s`$ can be determined from lattice gauge theory calculations of the hadronic spectrum. The basic procedure used is to choose a definition of $`\alpha _s`$, and measure its value on the lattice. One then has to set the scale at which $`\alpha _s`$ takes on the measured value. The lattice scale can be normalized using the hadronic spectrum measured on the same lattice. Finally, one has to convert the lattice definition of $`\alpha _s`$ to the value defined in the continuum in a scheme such as $`\overline{\mathrm{MS}}`$. There are several sources of systematic errors that limit the current accuracy in determining $`\alpha _s`$. Typically, $`\alpha _s`$ is determined by determining the spectrum of heavy quark bound states on the lattice. There are corrections due to the finite volume and finite lattice spacing $`a`$. The finite lattice spacing errors can be reduced by using improved actions, that are accurate to higher order in $`a`$. To some extent, one can estimate the error due to finite volume and finite lattice spacing by repeating the simulation on a larger lattice. The dominant systematic uncertainty is due to the quenched approximation, in which light quark loops are neglected. It is difficult to reliably estimate the systematic errors due to this approximation without doing a full simulation including dynamical fermions. Simulations with dynamical fermions are just starting to be done, and in a few years one should have more reliable estimates of $`\alpha _s`$. There is one important advantage to using a heavy quark system such as the $`\mathrm{{\rm Y}}`$ to determine $`\alpha _s`$. The leading correction to the $`\mathrm{{\rm Y}}`$ energy levels due to the light quark masses is linear in the quark masses, and can only depend on the flavor singlet combination $`m_u+m_d+m_s`$. Thus the light quark mass corrections can be computed to a good approximation using three light quarks of mass $`(m_u+m_d+m_s)/3`$. This avoids having to simulate almost massless dynamical quarks, which is very difficult. Lattice calculations can also be used to test theoretical calculations, and determine the regime in which perturbation theory is applicable. In the quenched approximation, one can study the scale dependence of the coupling constant on the lattice. This provides a check on the perturbation theory calculation with $`n_f=0`$. The result is in remarkable agreement with the perturbation theory result in the regime where the coupling constant is weak . The Fermilab and SCRI groups use the $`SP`$ and $`1S2S`$ splittings in the $`\mathrm{{\rm Y}}`$ system to determine $`\alpha _s`$. There are some systematic deviations of the calculated numbers from their experimental values in the quenched approximation ($`n_f=0`$), which are dramatically reduced if one includes $`n_f=2`$ dynamical flavors. The value of $`\alpha _s(M_Z)`$ in the $`\overline{\mathrm{MS}}`$ scheme is $$\alpha _s(M_Z)=0.1159\pm 0.0019\pm 0.0013\pm 0.0019$$ (59) where the first error is due to discretization effects, relativistic corrections, and statistical errors, the second is due to dynamical fermions, and the third is from conversion uncertainties. More recent computations give (in the $`n_f=5`$ scheme) $$\alpha _s(M_Z)=0.1174\pm 0.0024$$ (60) and $$\alpha _s(M_Z)=0.1118\pm 0.0017.$$ (61) An average of these newer values gives $`\alpha _s(M_Z)=0.115\pm 0.004`$, where we have included the difference between the two central values as an estimated additional systematic error. ## 8 $`\alpha _s`$ FROM HEAVY QUARK SYSTEMS Heavy quark bound states such as the $`\mathrm{{\rm Y}}`$ can also be used to extract a value for $`\alpha _s`$. If the bound state is treated using non-relativistic quantum mechanics, the annihilation decays $`\mathrm{{\rm Y}}\mu ^+\mu ^{}`$, $`\mathrm{{\rm Y}}ggg`$ and $`\mathrm{{\rm Y}}\gamma gg`$ can be computed as the product of the probability to find the quark-antiquark pair at the origin, times the annihilation rate for $`Q\overline{Q}`$ at rest to decay to the final state. The relevant Feynman graphs are shown in Fig. 5. The decay rate $`\mathrm{{\rm Y}}ggg`$ is the inclusive decay rate for $`\mathrm{{\rm Y}}\mathrm{hadrons}`$, and the decay rate for $`\mathrm{{\rm Y}}\gamma gg`$ is that for $`\mathrm{{\rm Y}}\gamma +\mathrm{hadrons}`$. The probability to find $`Q\overline{Q}`$ at the origin, $`\left\|\psi (0)\right\|^2`$, is sensitive to the detailed dynamics of the $`Q\overline{Q}`$ bound state. If one takes the ratio of decay rates, $`\left\|\psi (0)\right\|^2`$ drops out, and the ratio of decay rates can be used to determine $`\alpha _s`$. The above qualitative discussion can be made precise using NRQCD (non-relativistic QCD) to calculate the properties of the $`\mathrm{{\rm Y}}`$ . NRQCD has an expansion in powers of $`v`$, the velocity of quarks in the bound state. The $`\alpha _s`$ expansion is coupled to the $`v`$ expansion, since $`v\alpha _s`$ in a Coulombic system. The NRQCD approach allows one to systematically factor the decay rate into short distance coefficients that are calculable in perturbation theory, and non-perturbative hadronic matrix elements that generalize the notion of the wavefunction at the origin. The $`\mathrm{{\rm Y}}`$ wavefunction in NRQCD has different Fock components. The lowest order (in $`v`$) component is $`\|\overline{Q}Q`$ and the first correction contains a gluon, $`\|\overline{Q}Qg`$. The $`\|\overline{Q}Qg`$ is referred to as the color-octet component, because the two quarks are in a color octet state. The NRQCD velocity counting rules show that the probability to find the $`\mathrm{{\rm Y}}`$ in the octet component is of order $`v^2`$. The decay rate for $`\mathrm{{\rm Y}}\mu ^+\mu ^{}`$ in NRQCD is $$\mathrm{\Gamma }(\mathrm{{\rm Y}}\mu ^+\mu ^{})=\frac{2\mathrm{I}\mathrm{m}f_{ee}({}_{}{}^{3}S_{1}^{})}{M_b^2}\mathrm{{\rm Y}}\left\|𝒪_1({}_{}{}^{3}S_{1}^{})\right\|\mathrm{{\rm Y}}+\frac{2\mathrm{I}\mathrm{m}g_{ee}({}_{}{}^{3}S_{1}^{})}{M_b^4}\mathrm{{\rm Y}}\left\|𝒫_1({}_{}{}^{3}S_{1}^{})\right\|\mathrm{{\rm Y}}$$ (62) where $`M_b`$ is the $`b`$-quark pole mass, $`𝒪_1({}_{}{}^{3}S_{1}^{})=\psi ^{}𝝈\chi \chi ^{}𝝈\psi `$ and $`𝒫_1({}_{}{}^{3}S_{1}^{})=i/4(\psi ^{}𝝈\chi \chi ^{}𝝈(𝐃)^\mathrm{𝟐}\psi +\mathrm{h}.\mathrm{c}.)`$, and $`\psi `$ and $`\chi ^{}`$ annihilate quarks and antiquarks, respectively. The first term is the leading order contribution, and the second term is the $`v^2`$ correction. The coefficients $`f_{ee}({}_{}{}^{3}S_{1}^{})`$ and $`g_{ee}({}_{}{}^{3}S_{1}^{})`$ can be computed from the first graph in Fig. 5 at lowest order in perturbation theory. The values are $$\mathrm{Im}f_{ee}({}_{}{}^{3}S_{1}^{})=\frac{\pi \alpha ^2e_b^2}{3}\left[116\frac{\alpha _s}{\pi }\right],\mathrm{Im}g_{ee}({}_{}{}^{3}S_{1}^{})=\frac{4\pi \alpha ^2e_b^2}{9},$$ (63) where $`e_b=1/3`$ is the charge of the $`b`$-quark, and the radiative correction to $`f_{ee}({}_{}{}^{3}S_{1}^{})`$ has also been included. The decay rate for $`\mathrm{{\rm Y}}\mathrm{hadrons}`$ is $`\mathrm{\Gamma }(\mathrm{{\rm Y}}\mathrm{hadrons})`$ $`=`$ $`{\displaystyle \frac{2\mathrm{I}\mathrm{m}f_1({}_{}{}^{3}S_{1}^{})}{M_b^2}}\mathrm{{\rm Y}}\left\|𝒪_1({}_{}{}^{3}S_{1}^{})\right\|\mathrm{{\rm Y}}+{\displaystyle \frac{2\mathrm{I}\mathrm{m}g_1({}_{}{}^{3}S_{1}^{})}{M_b^4}}\mathrm{{\rm Y}}\left\|𝒫_1({}_{}{}^{3}S_{1}^{})\right\|\mathrm{{\rm Y}}`$ (64) $`+\mathrm{\Gamma }^{(8)}(\mathrm{{\rm Y}}\mathrm{hadrons})`$ where $`\mathrm{\Gamma }^{(8)}`$ is the contribution to the decay rate from the color-octet component of the $`\mathrm{{\rm Y}}`$. In NRQCD, the color octet decay rate is $`v^4`$ suppressed relative to the color singlet decay rate. However, the color octet component can decay into two gluons, rather than three, so the color octet decay rate of order $`\alpha _s^2v^4`$ can compete with the relativistic correction to the color singlet decay rate of order $`\alpha _s^3v^2`$. The coefficient functions are $`\mathrm{Im}f_1({}_{}{}^{3}S_{1}^{})`$ $`=`$ $`{\displaystyle \frac{10}{243}}(\pi ^29)\alpha _s^3(M_b)`$ $`\times [1+(9.46(2)C_F+4.13(17)C_A1.161(2)n_f){\displaystyle \frac{\alpha _s}{\pi }})]`$ $`\mathrm{Im}g_1({}_{}{}^{3}S_{1}^{})`$ $`=`$ $`{\displaystyle \frac{5}{1458}}(19\pi ^2132)\alpha _s^3(M_b)`$ (65) The $`\alpha _s^3`$ term for $`\mathrm{Im}f_1({}_{}{}^{3}S_{1}^{})`$ has the given value when the scale of the $`\alpha _s^2`$ term is the $`b`$-quark pole mass. The decay rate $`\mathrm{\Gamma }(\mathrm{{\rm Y}}\gamma +\mathrm{hadrons})`$ has the form $$\mathrm{\Gamma }(\mathrm{{\rm Y}}\gamma +\mathrm{hadrons})=\frac{2\mathrm{I}\mathrm{m}f_\gamma ({}_{}{}^{3}S_{1}^{})}{M_b^2}\mathrm{{\rm Y}}\left\|𝒪_1({}_{}{}^{3}S_{1}^{})\right\|\mathrm{{\rm Y}}+\frac{2\mathrm{I}\mathrm{m}g_\gamma ({}_{}{}^{3}S_{1}^{})}{M_b^2}\mathrm{{\rm Y}}\left\|𝒫_1({}_{}{}^{3}S_{1}^{})\right\|\mathrm{{\rm Y}},$$ (66) where the coefficient function is $`\mathrm{Im}f_\gamma ({}_{}{}^{3}S_{1}^{})`$ $`=`$ $`{\displaystyle \frac{8}{27}}(\pi ^29)\alpha _s^2(M_b)\alpha e_b^2`$ (67) $`\times [1+(9.46(2)C_F+2.75(11)C_A0.774(1)n_f){\displaystyle \frac{\alpha _s}{\pi }})]`$ The equations of motion can be used to relate the matrix elements of the $`S`$-wave and $`P`$-wave operators , $$\mathrm{{\rm Y}}\left\|𝒫_1({}_{}{}^{3}S_{1}^{})\right\|\mathrm{{\rm Y}}=\frac{M_\mathrm{{\rm Y}}M_b}{M_b}\mathrm{{\rm Y}}\left\|𝒪_1({}_{}{}^{3}S_{1}^{})\right\|\mathrm{{\rm Y}}+O(v^2).$$ (68) The matrix element of the $`S`$-wave operator is the NRQCD analog of the wavefunction at the origin in a potential model calculation. The matrix element is non-perturbative, and can be eliminated by considering ratios of decay rates. The matrix element of $`𝒫_1({}_{}{}^{3}S_{1}^{})`$ can be determined from Eq. (68) using estimates of the $`b`$-quark pole mass $`M_b`$, and the measured $`\mathrm{{\rm Y}}`$ mass, as was done in Ref. . However, one can instead replace $`(M_\mathrm{{\rm Y}}2M_b)/M_b`$ by $`4/9\alpha _s^2`$, the lowest-order result for the binding energy for a Coulomb bound state. This reduces somewhat the uncertainty in the extraction of $`\alpha _s`$, since it eliminates any uncertainty from the pole mass. The experimental value of the ratio $`\mathrm{\Gamma }(\mathrm{{\rm Y}}\mathrm{hadrons})/\mathrm{\Gamma }(\mathrm{{\rm Y}}\mathrm{}^+\mathrm{}^{})=39.11\pm 0.4`$ , where $`\mathrm{}=e,\mu ,\tau `$ gives $`\alpha _s(M_b)=0.177\pm 0.01`$, using the ratio of the theoretical formula for the decay widths. The unknown octet decay rate has been estimated to be less than 9% in Ref. , and this has been included as a theoretical uncertainty. The decay rates Eq. (62)–(66) have been written in terms of $`\alpha _s(M_b)`$. One can instead rewrite them in terms of $`\alpha _s(\mu )`$ using Eq. (8), extract $`\alpha _s(\mu )`$, and convert this into $`\alpha _s(M_b)`$ using Eq. (8). We have include the uncertainty of a scale change by a factor of two in the theoretical estimate. The octet and scale uncertainties are comparable in size. The experimental value of the ratio $`\mathrm{\Gamma }(\mathrm{{\rm Y}}\gamma +\mathrm{hadrons})/\mathrm{\Gamma }(\mathrm{{\rm Y}}\mathrm{hadrons})=2.75\pm 0.04\pm 0.15`$ can also be used to extract $`\alpha _s`$. It is convenient to use the experimental value of $`\mathrm{\Gamma }(\mathrm{{\rm Y}}\mathrm{hadrons})/\mathrm{\Gamma }(\mathrm{{\rm Y}}\mathrm{}^+\mathrm{}^{})`$ to convert this to $`\mathrm{\Gamma }(\mathrm{{\rm Y}}\gamma +\mathrm{hadrons})/\mathrm{\Gamma }(\mathrm{{\rm Y}}\mathrm{}^+\mathrm{}^{})=1.075\pm 0.06`$, before comparing with the theoretical results. This eliminates the theoretical uncertainty due to the octet component in the hadronic decay rate. The extracted value of $`\alpha _s(M_b)`$ is $`0.189\pm 0.01`$, where we have included a scale uncertainty as above. Averaging the two extractions gives $`\alpha _s(M_b)=0.183\pm 0.01`$ which corresponds to $`\alpha _s(M_Z)=0.108\pm 0.004`$ ## 9 $`\alpha _s`$ FROM HADRON-HADRON SCATTERING There are many process at high-energy hadron-hadron colliders which can constrain the value of $`\alpha _s`$. All rely on the QCD improved parton model, and on the factorization theorems of QCD . The rate for any process is expressed as a convolution of the partonic scattering amplitude $`\sigma _{i,j}^p`$ and parton distribution functions discussed in section 3; see Eq 28 (note that here we use $`f`$ rather than $`q`$ as the sum on $`i,j`$ runs over quarks and gluons. $$\sigma =\underset{i,j}{}𝑑x_1𝑑x_2f_i(x,M^2)f_j(x,M^2)\sigma _{ij}^p(M)$$ (69) The factorization scale $`M`$ is arbitrary. As in the case of the scale $`\mu `$ used in $`\alpha _s(\mu )`$ (see Eq 10 and the surrounding discussion), the exact result cannot depend on its choice. However as the processes $`\sigma _{i,j}^p`$ is only calculated to some finite order in perturbation theory, some residual $`M`$ dependence will remain. As in the case of $`\mu `$ the sensitivity to $`M`$ will be small if it is chosen to be a characteristic scale of the process; for example, in the case of the production of a pair of jets of momentum, $`p_T`$, transverse to the direction defined by the incoming hadrons, $`M=p_T`$ is a reasonable choice. The quantitative tests of QCD and the consequent extraction of $`\alpha _s`$ which appears in $`\sigma _{i,j}^p`$ are possible only if the process in question has been calculated beyond leading order in QCD perturbation theory. The production of hadrons with large transverse momentum in hadron-hadron collisions provides a direct probe of the scattering of quarks and gluons: $`qqqq`$, $`qgqg`$, $`gggg`$, etc.. Here the leading order term in $`\sigma _{i,j}^p`$ is of order $`\alpha _s^2`$ so the rates are sensitive to its value. Higher–order QCD calculations of the jet rates and shapes are in impressive agreement with data . This agreement has led to the proposal that these data could be used to provide a determination of $`\alpha _s`$ . A set of structure functions is assumed and Tevatron collider data are fitted over a very large range of transverse momenta, to the QCD prediction for the underlying scattering process that depends on $`\alpha _s`$. The evolution of the coupling over this energy range (40 to 250 GeV) is therefore tested in the analysis. CDF obtains $`\alpha _s(M_Z)=0.1129\pm 0.0001(\text{stat.})\pm 0.0085(\text{syst.})`$ . Estimation of the theoretical errors is not straightforward. The structure functions used depend implicitly on $`\alpha _s`$ and an iteration procedure must be used to obtain a consistent result; different sets of structure functions yield different correlations between the two values of $`\alpha _s`$. We estimate an uncertainty of $`\pm 0.005`$ from examining the fits. Ref. estimates the error from unknown higher order QCD corrections to be $`\pm 0.005`$. Combining these then gives: $`\alpha _s(M_Z)=0.1129\pm 0.011`$ QCD corrections to Drell-Yan type cross sections (the production in hadron collisions by quark-antiquark annihilation of lepton pairs of invariant mass $`Q`$ from virtual photons or of real $`W`$ or $`Z`$ bosons), are known . These processes are not very sensitive to $`\alpha _s`$ as the leading piece in $`\sigma _{i,j}^p`$ is of order $`\alpha _s^0`$. The production of $`W`$ and $`Z`$ bosons and photons at large transverse momentum begins at order $`\alpha _s^0`$. The leading-order QCD subprocesses are $`q\overline{q}\gamma g`$ and $`qg\gamma q`$. The next-to-leading-order QCD corrections are known for photons, and for $`W/Z`$ production , and so an extraction of $`\alpha _s`$ is possible in principle. Data exist on photon production from the CDF and DØ collaborations and from fixed target experiments . Detailed comparisons with QCD predictions may indicate an excess of the data over the theoretical prediction at low value of transverse momenta, although other authors find smaller excesses. These differences indicate that while the process may be understood, no meaningful extraction of $`\alpha _s`$ is possible. The UA2 collaboration has extracted a value of $`\alpha _s(M_W)=0.123\pm 0.018(\text{stat.})\pm 0.017(\text{syst.})`$ from the measured ratio $`R_W=\sigma (W+1\text{jet})/\sigma (W+0\text{jet})`$. The result depends on the algorithm used to define a jet, and the dominant systematic errors due to fragmentation and corrections for underlying events (the former causes jet energy to be lost, the latter causes it to be increased) are connected to the algorithm. The scale at which $`\alpha _s(M)`$ is to be evaluated is not clear. A change from $`\mu =M_W`$ to $`\mu =M_W/2`$ causes a shift of 0.01 in the extracted $`\alpha _s`$, and the quoted error should be increased to take this into account. There is also dependence on the parton distribution functions, and hence, $`\alpha _s`$ appears explicitly in the formula for $`R_W`$, and implicitly in the distribution functions. Data from CDF and DØ on the W $`p_T`$ distribution are in agreement with QCD but are not able to determine $`\alpha _s`$ with sufficient precision to have any weight in a global average. The production rates of $`b`$ quarks in $`p\overline{p}`$ have been used to determine $`\alpha _s`$ . The next to leading order QCD production processes have been used. At order $`\alpha _s`$ the production processes are $`ggb\overline{b}`$ and $`q\overline{q}b\overline{b}`$ result in b-hadrons that are back to back in azimuth. By selecting events in this region the next-to leading order calculation can be used to compare rates to the measured value and a value of $`\alpha _s`$ extracted. The errors are dominated by the measurement errors, the choice of $`\mu `$ and $`M`$, and uncertainties in the choice of structure functions. The last were estimated by varying the structure functions used. The result is $`\alpha _s(M_Z)=0.113_{0.013}^{+0.009}`$. ## 10 CONCLUSION The previous sections have illustrated the large number of processes where quantitative tests of QCD can be made and a value of $`\alpha _s`$ extracted. Figure 6 shows the values of $`\alpha _s(M_Z)`$ deduced from the various processes shown above. The consistency and precision of these results is remarkable. Figure 7 shows the values of $`\alpha _s(\mu )`$ and the values of $`\mu `$ where they are measured. This figure clearly shows the experimental evidence for the variation of $`\alpha _s(\mu )`$ with $`\mu `$ predicted by Eq.4. An average of the values in Figure 6 and in Table 1 gives $`\alpha _s(M_z)=0.1173`$, with a total $`\chi ^2`$ of 9 for twelve fitted points, showing good consistency among the data. The value from heavy quark systems contributes slightly more that one half of the total $`\chi ^2`$. If this result is omitted the average increases to 0.1185. All of the other results are within $`1.1\sigma `$ of the average value. The error on the average, assuming that all of the errors in the contributing results are uncorrelated, is $`\pm 0.0014`$, and may be an underestimate. We have seen that in almost all of the cases discussed, the errors are dominated by systematic, usually theoretical errors. Only some of these, notably from the choice of scale, are correlated. it is important to note that the average is not dominated by a single measurement; there are many results with comparable small errors: from $`\tau `$ decay, lattice gauge, theory deep inelastic scattering and the $`Z^0`$ width. We quote our average value as $`\alpha _s(M_Z)=0.1173\pm 0.002`$, which corresponds to $`\mathrm{\Lambda }^{(5)}=200_{23}^{+24}`$ MeV using Eq. 13. The reader may wish to consult other recent articles for different opinions . Significant improvements in the precision in the near future are not likely. The accuracy of data from LEP will not improve. It is possible that a better understanding of the jet rates in hadron-hadron colliders and a systematic treatment of the errors from the structure functions will lead to and improvement in the precision of the value of $`\alpha _s`$ derived. In many cases where the data are quite precise, such as heavy quark system, theoretical uncertainties limit the precision. In the very long term precision at the 1% level may be achievable . ## Acknowledgments This work was supported by the Director, Office of Science, Office of Basic Energy Services, of the U.S. Department of Energy under Contract DE-AC03-76SF0098, and by DOE grant DOE-FG03-97ER40546.
warning/0004/cond-mat0004191.html	ar5iv	text	# Spectral partitions on infinite graphs ## I Introduction The study of model systems without translation invariance is an interesting and complex subject of modern statistical mechanics. A very general description of this situation is in terms of statistical models on graphs, that is on generic networks formed by sites, where dynamical variables reside, and links connecting pairwise sites whose variables are coupled. This is the direct extension of the typical setup valid for crystalline lattices, which are indeed very special, homogeneous graphs. On the other hand, graphs are not in general homogeneous and the main question is how these inhomogeneities affect physical properties and give rise to relevant changes with respect to lattices. While small scale inhomogeneities will affect local properties, one expects that only large scale inhomogeneities are relevant for bulk thermodynamic properties. Most likely, the latter properties are those that show universal features which depend only on a few global parameters, just as in the case of lattices. The study of such universality requires consideration of infinite graphs (with certain natural restrictions given below), where the thermodynamic limit is taken. The main relevant geometrical parameter affecting universal properties is the spectral dimension $`\overline{d}`$ of an infinite graph $`𝒢`$ . It generalizes the Euclidean dimension of lattices to arbitrary real values and is naturally defined from the infrared behaviour of the spectral density of the Laplacian operator on $`𝒢`$ . An equivalent definition, the one adopted in this work, is in terms of average properties of random walks on $`𝒢`$ at large times, that is to say of the singularities of the Gaussian model on the same graph . On the other hand, the spectral dimension of the whole graph $`𝒢`$, by itself turns out not to be sensitive to macroscopic inhomogeneities strong enough to give rise to true thermodynamic inhomogeneities. Indeed it may happen that distinct macroscopic parts of an infinite graph exhibit distinct thermodynamic behaviour. We shall show below that such parts can be characterized in terms of their own spectral dimension, possibly plus a spectral weight, resulting in an effective spectral partition of $`𝒢`$. The crucial point is that these parts form subgraphs which are thermodynamically independent, that is to say completely uncoupled as far as thermodynamic properties are concerned. In other words, inhomogeneous thermodynamic behaviour on the same infinite graph necessarily imply effective decoupling. ## II Infinite graphs: basic definitions, measure and averages A (unoriented) graph $`𝒢`$ is the ordered couple $`(G,G_L)`$ formed by a countable set $`G`$ of vertices (or sites, or nodes), that we shall generically indicate with small-case Latin letters, $`i`$, $`j`$, $`k`$, …, and a set $`G_L`$ of unoriented links (or bonds) which connect pairwise the sites and are therefore naturally denoted by couples $`(i,j)=(j,i)`$. When the set $`G`$ is finite, $`𝒢`$ is a finite graph and we shall denote $`N`$ the number of vertices of $`𝒢`$. A subgraph $`𝒢^{}`$ of $`𝒢`$ is a graph such that $`G^{}G`$ and $`G_L^{}G_L`$. A subgraph is said to be complete if its has all the available links, that is if, given the subset of nodes $`G^{}`$, the subset of links $`G_L^{}`$ is the largest possible one. A path in $`𝒢`$ is a sequence of consecutive links $`\{(i,k)(k,h)\mathrm{}(n,m)(m,j)\}`$. A graph is said to be connected, if for any two points $`i,jG`$ there is always a path joining them. In the following we will consider only connected graphs. The graph topology can be algebraically described by its adjacency matrix $`𝑨`$ with elements $$A_{ij}=\{\begin{array}{cc}1& \mathrm{if}(i,j)G_L\hfill \\ 0& \mathrm{if}(i,j)G_L\hfill \end{array}$$ (1) The Laplacian matrix $`𝑳`$ on the graph $`𝒢`$ has elements: $$L_{ij}=z_i\delta _{ij}A_{ij}$$ (2) where $`z_i=_jA_{ij}`$, the number of nearest neighbours of $`i`$, is called the coordination number (or degree) of site $`i`$. Here we will consider graphs with $`z_{max}=sup_iz_i<\mathrm{}`$. One can also consider a generalization of the adjacency matrix, which corresponds to the ferromagnetic and uniformly bounded coupling $`J_{ij}`$, with $`J_{ij}0A_{ij}=1`$ and $`supJ_{ij}<\mathrm{}`$, $`infJ_{ij}>0`$. The elements of the generalized Laplacian matrix then read: $$_{ij}=J_i\delta _{ij}J_{ij}$$ (3) where $`J_i=_jJ_{ij}`$. Every connected graph $`𝒢`$ is endowed with an intrinsic metric generated by the chemical distance $`r_{i,j}`$ which is defined as the number of links in the shortest path(s) connecting vertices $`i`$ and $`j`$. Let us now consider thermodynamic averages on infinite graphs . The Van Hove sphere $`S_{o,r}𝒢`$ of centre $`oG`$ and radius $`r`$ is the complete subgraph of $`𝒢`$ containing all $`iG`$ whose distance from $`o`$ is $`r`$ and all the links of $`𝒢`$ joining them. We will call $`N_{o,r}`$ the number of vertices contained in $`S_{o,r}`$. In the thermodynamic limit the average $`[f]_G`$ of a real–valued function $`f`$ on $`G`$ is: $$[f]_G\underset{r\mathrm{}}{lim}\frac{1}{N_{o,r}}\underset{iS_{o,r}}{}f_i$$ (4) This average does not depend on the choice of the origin $`oG`$ provided $`f`$ is bounded from below and $$\underset{r\mathrm{}}{lim}\frac{\|S_{o,r}\|}{N_{o,r}}=0$$ (5) where $`\|S_{o,r}\|`$ is the number of the vertices of the sphere $`S_{o,r}`$ connected with the rest of the graph . Here we shall restrict our attention to graphs with this property. The measure $`\|A\|`$ of a subset $`AG`$ is the average value $`[\chi (A)]_G`$ of its characteristic function $`\chi _i(A)`$ defined by $`\chi _i(A)=1`$ if $`iA`$ and $`\chi _i(A)=0`$ if $`iA`$. The measure of a subset of links $`G_L^{}G_L`$ is similarly given by: $$\|G_L^{}\|\underset{r\mathrm{}}{lim}\frac{N_{L,r}^{}}{N_{o,r}}$$ (6) where $`N_{L,r}^{}`$ is the number of links of $`G_L^{}`$ contained in the sphere $`S_{o,r}`$. Any two nonzero–measure subsets $`A`$ and $`B`$ of $`G`$ are said to be equivalent if their symmetric difference has zero measure, that is $`\|A\|=\|B\|=\|AB\|`$. For any given nonzero–measure subsets $`AG`$ we shall denote $`\{A\}`$ its equivalence class. Then $`A`$ is said to be a representative of $`\{A\}`$. With the subgraph $`𝒢^{}`$ defined by the ordered double $`(G^{},G_L^{})`$, we identify the measure of the subgraph as the measure $`\|G^{}\|`$ of its points. Given a (nonzero–measure) subset $`AG`$, we define the average on $`A`$ of any real–valued function $`f`$ on $`G`$ as: $$[f]_A=[\chi (A)f]_G$$ (7) By definition $`[f]_A`$ is a function only of the equivalence classes, that is $`[f]_A=[f]_{\{A\}}`$. Moreover, quite evidently $`[f]_C=[f]_A+[f]_B`$ whenever $`C=AB`$ and $`\|AB\|=0`$. Given a complete subgraph $`=(M,M_L)`$, we denote $`\overline{}`$ its complement in $`𝒢`$. This is formed by all points that do not belong to $`M`$ and by all links of $`𝒢_L`$ which connect them. $`\overline{}`$ is therefore a complete subgraph. We call the pair $`(,\overline{})`$ a partition of order two of $`𝒢`$ whenever both $`M`$ and its complement $`\overline{M}`$ are nonzero–measure subsets of $`G`$. We introduce now the important concept of minimal distance $`\underset{¯}{D}(𝒜,)`$ between any pair $`𝒜,`$ of nonzero–measure subgraphs of $`𝒢`$ such that $`\|AB\|=0`$. It is defined as $$\underset{¯}{D}(𝒜,)=\mathrm{min}(n:\left\|A_nB\right\|>0)$$ (8) where $$A_nB=\{iA:\text{dist}(i,B)=n\},\text{dist}(i,B)=\underset{jB}{\mathrm{min}}\text{r}_{(i,j)}$$ (9) For $`n=0`$, $`_n`$ reduces to the usual intersection operator. Notice that, while in general the relation $`A_nB`$ is not symmetric in $`A,B`$, the minimal distance is symmetric: $`\underset{¯}{D}(𝒜,)=\underset{¯}{D}(,𝒜)`$. In fact, from the boundedness of $`z_i`$, it can be shown by induction on $`n`$ that $$\|B_nA\|(z_{max})^n\|A_nB\|$$ (10) so that $$\|A_nB\|>0\|B_nA\|>0$$ (11) implying our assertion. Consider now the minimal distance between the two members of a partition of order two. Suppose $`\underset{¯}{D}(,\overline{})=n>1`$; then $`\|M_n\overline{M}\|>0\|M_{n1}\overline{M}\|>0`$ from the boundedness of $`z_i`$. This implies that if $`\underset{¯}{D}(,\overline{})`$ is finite, then $`\underset{¯}{D}(,\overline{})=1`$. In this case we may say that $``$ and $`\overline{}`$ are densely interlaced, while in the opposite case that they are infinitely separated. From the definition of minimal distance, it follows that if two subgraphs $`𝒜`$ and $``$ of $`𝒢`$ are infinitely separated, their common frontier $`(𝒜,)`$ (i.e. the links $`(i,j)𝒢_L`$ with $`iA`$ and $`jB`$) is a zero–measure set. Then the two subgraphs can be disconnected by cutting such a zero–measure set of links. This relates the property of infinite separability to the simple separability property defined in . Indeed, the two definitions coincide. We shall term separable partition a partition $`(,\overline{})`$ where $``$ and $`\overline{}`$ are infinitely separated. ## III The Gaussian model: infrared behaviour and the spectral dimension The Gaussian model on $`𝒢`$ is defined by assigning a real–valued random variable $`\varphi _i`$ to each node $`iG`$ and then prescribing the following probability measure $$d\mu _r[\varphi ]=\frac{1}{Z_r}\mathrm{exp}\left[\underset{i,jS_{o,r}}{}\varphi _i(𝑳+m^2𝜼)_{ij}\varphi _j\right]\underset{iS_{o,r}}{}d\varphi _i$$ (12) for the collection $`\varphi =\{\varphi _i;iS_{o,r}\}`$. Here $`Z_r`$ is the proper normalization factor, $`m>0`$ is a free parameter and $`𝜼`$ is the diagonal matrix with elements $`\eta _{ij}=\eta _i\delta _{ij}`$ with the real numbers $`\eta _i`$ positive definite and uniformly bounded throughout $`G`$ (that is $`0<\eta _{\mathrm{min}}\eta _i\eta _{\mathrm{max}}`$, $`i𝒢`$). The thermodynamic limit is achieved by letting $`r\mathrm{}`$ and defines a Gaussian measure over the entire $`\varphi =\{\varphi _i;iG\}`$ which does not depend on the centre of the Van Hove sphere $`o`$ . The covariance of this Gaussian process reads $$\varphi _i\varphi _jC_{ij}(m^2)=(𝑳+m^2𝜼)_{ij}^1$$ (13) and hence it satisfies by definition the Schwinger–Dyson (SD) equation $$(J_i+m^2\eta _i)C_{ij}(m^2)\underset{kG}{}J_{ik}C_{kj}(m^2)=\delta _{ij}$$ (14) Setting $$C_{ij}=\frac{(1W)_{ij}^1}{J_i+m^2\eta _i},W_{ij}=\frac{J_{ij}}{J_j+m^2\eta _j}$$ (15) one obtains the standard connection with the random walk (RW) over $`𝒢`$ : $$(1W)_{ij}^1=\underset{t=0}{\overset{\mathrm{}}{}}(W^t)_{ij}=\underset{\gamma :ij}{}W[\gamma ]$$ (16) where the last sum runs over all paths from $`j`$ to $`i`$, each weighted by the product along the path of the one–step probabilities in $`W`$: $$\gamma =(i,k_{t1},\mathrm{},k_2,k_1,j)W[\gamma ]=W_{ik_{t1}}W_{k_{t1}k_{t2}},\mathrm{},W_{k_2k_1}W_{k_1j}$$ (17) Notice that, as long as $`m>0`$, we have $`_i(W^t)_{ij}<1`$ for any $`t`$, namely the walker has a nonzero death probability. This implies that $`C_{ij}`$ is a smooth functions of $`m^2`$ for $`mϵ>0`$. In the limit $`m0`$ the walker never dies and the sum over paths in eq. (16) is dominated by the infinitely long paths which sample the large scale structure of the entire graph (“large scale” refers here to the metric induced by the chemical distance alone). This typically reflects itself into a singularity of $`C_{ij}`$ at $`m=0`$ whose nature does not depend on the detailed form of $`J_{ij}`$ or $`\eta _i`$, as long these stay uniformly positive and bounded. Of particular importance is the leading singular infrared behaviour, as $`m^20`$, of the average $`[C(m^2)]_G`$ of $`C_{ii}(m^2)`$, which is a positive definite quantity, over all points $`i`$ of the graph $`𝒢`$, which we may write in general as $$\mathrm{Sing}[C(m^2)]_Gc(m^2)^{\overline{d}/21}$$ (18) The parameter $`\overline{d}`$ is called the spectral dimension of the graph $`𝒢`$ and on regular lattices it coincides with the usual Euclidean dimension. Henceforth we shall call spectral weight the coefficient $`c`$ in eq. (18). The name spectral dimension is related to the behaviour of the spectral density $`\rho (l)`$ of low-lying eigenvalues of the Laplacian $`𝑳`$; indeed it can be shown that $`\rho (l)`$ scales as a power of $`ł`$ for $`ł0`$, that is $`\rho (l)l^{\overline{d}/21}`$. ## IV Large scale inhomogeneity: homogeneity classes and spectral classes In the study of statistical models one often has to deal with the average $`[C(m^2)]_A`$ of $`C_{ii}(m^2)`$ over a generic positive measure subset $`AG`$ and in particular one has to consider the leading singular behaviour of $`[C(m^2)]_A`$ as $`m^20`$. On regular lattices this singular behaviour is independent of $`A`$ and it actually coincides with that obtained averaging over all points of $`𝒢`$: $$\mathrm{Sing}[C(m^2)]_A=\mathrm{Sing}[C(m^2)]_G,AG,\|A\|>0$$ (19) This property arises from the large scale homogeneity of regular lattices due to translation invariance. On graphs, where translation invariance is lost, this property can still hold if the inhomogeneity is limited to finite scales. More generally it may happen that inhomogeneity extends to large scales and the singular parts of eq. (19) are different on different subsets. However we will prove that such subsets must satisfy very strong topological constraints: a large–scale inhomogeneous graph always consists of homogeneous parts joined together by a zero–measure set of links. Therefore the splitting of infrared behaviour always corresponds to a macroscopically evident inhomogeneity of the graph. In this section we will give a rigorous formulation of these statements through the following steps. * Let us suppose that the graph $`𝒢`$ has indeed a large–scale inhomogeneity that manifests itself through the existence of at least one nonzero–measure subset $`AG`$ such that, as $`m^20`$, $$\mathrm{Sing}[C(m^2)]_Ac_A(m^2)^{\overline{d}_A/21}$$ (20) with $`\overline{d}_A\overline{d}`$. * We then define $`MG`$ to be a maximally homogeneous (or more briefly maximal) subset with respect to $`\overline{d}_A`$ whenever 1. $`\|MA\|>0`$ 2. $`\mathrm{Sing}[C(m^2)]_Mc_M(m^2)^{\overline{d}_M/21}`$, with $`\overline{d}_M=\overline{d}_A`$. 3. For any nonzero–measure subset $`BM`$ we have $`\overline{d}_B=\overline{d}_M`$. 4. There exists no $`BM`$ such that $`\overline{d}_B=\overline{d}_M`$ and $`\|B\|>\|M\|`$. By this definition it follows that the set of all maximal subsets with respect to $`\overline{d}_M`$ coincides with the equivalence class $`\{M\}`$ and we will call it the homogeneity class of $`\overline{d}_M`$. * Next we prove the Theorem 1: The subgraphs $``$ and its complement $`\overline{}`$ are infinitely separated, i.e. their minimal distance $`\underset{¯}{D}(,\overline{})`$ is infinite and they define a separable partition of $`𝒢`$. Since this separability is induced by the spectral properties embodied by the spectral dimension, we call this a spectral partition (of order two) of $`𝒢`$. * Finally we consider a Gaussian model on the graph $``$ showing that, from the infinite separability of $``$ and $`\overline{}`$ the spectral dimension of $``$ is $`\overline{d}_M`$. Therefore, $`\overline{d}_M`$ is a property of the graph $``$ and defines a spectral class. This chain of arguments may now be applied to $`\overline{}`$, splitting off a new spectral class if $`\overline{}`$ has a large scale inhomogeneity of the type given above. The process can be repeated until necessary, yielding a complete spectral partition of the original graph $`𝒢`$ into spectral classes. Proof of Theorem 1: Let us suppose ad absurdum that $`\underset{¯}{D}(,\overline{})=1`$ and therefore that there exists a nonzero–measure subset $`\overline{M}^{}\overline{M}`$ such that $`\underset{¯}{D}(,\overline{}^{})=1`$. From the maximality of $`M`$ it follows that $`\overline{d}_M\overline{d}_{\overline{M}^{}}`$. Let us consider the random walk representation (16) of $`C_{ii}(m^2)`$ with $`i\overline{M}^{}`$: $$C_{ii}(m^2)=\frac{1}{J_i+m^2\eta _i}\underset{\gamma :ii}{}W[\gamma ]$$ (21) Next consider a site $`kM`$ whose distance from $`i`$ is $`1`$. This site exists from the hypothesis $`\underset{¯}{D}(,\overline{}^{})=1`$. Then, from the sum over paths in the left hand side of (21) let us retain only the paths containing $`k`$. Then, from the boundedness and positivity of $`J_{ij}`$ and $`\eta _i`$ one gets: $$C_{ii}(m^2)\frac{C_{kk}(m^2)}{J_{max}+m^2\eta _{max}}$$ (22) Averaging over $`M`$ and then over $`\overline{M}^{}`$ we get: $$[C(m^2)]_{\overline{M}^{}}K[C(m^2)]_M$$ (23) where $`K`$ is a positive constant. Now, taking $`m^20`$ and using the asymptotic expression for $`[C(m^2)]`$ given in (18) we obtain $$(m^2)^{\overline{d}_{\overline{M}^{}}/21}K^{}(m^2)^{\overline{d}_M/21}$$ (24) Since this argument applies equally well with $``$ and $`\overline{}`$ interchanged, one gets: $$(m^2)^{\overline{d}_M/21}K^{\prime \prime }(m^2)^{\overline{d}_{\overline{M}^{}}/21}$$ (25) which gives $`\overline{d}_M=\overline{d}_{\overline{M}^{}}`$ contradicting the hypothesis. Therefore $`\underset{¯}{D}(,\overline{})=\mathrm{}`$ and $``$ and $`\overline{}`$ must be infinitely separated. The infinite separability of $``$ and $`\overline{}`$ implies that the two subgraphs can be disconnected by cutting a zero–measure set of links. This very peculiar property implies thermodynamic independence, that is the decoupling, in the thermodynamic limit, of a model defined on the whole graph $`𝒢`$ into two models defined independently on on $``$ and $`\overline{}`$ . This applies in particular to the Gaussian model, so that the two averages of $`C_{ii}(m^2)`$ on $``$ and $`\overline{}`$ are independent quantities, each satisfying a relation like eq. (18) with two distinct spectral dimensions. Most importantly, to any nonzero–measure subset of $``$ there corresponds by construction the same spectral dimension $`\overline{d}`$ of $``$. We can say then that $`\overline{d}`$ is a universal property of $``$. ## V Spectral weights and Subclasses of spectral classes In the singular behaviour of $`[C(m^2)]`$, inhomogeneity at large scale can appear also in the coefficient of the leading infrared part (18). However, following the same steps as the previous section, we will show that once again a splitting in the value of the coefficient corresponds to a macroscopic inhomogeneity of the graph and that a macroscopically homogeneous graph is indeed characterized by universal $`\overline{d}`$ and $`c`$. Actually in this case the proof is subtler and requires some further mathematical steps. We first define the spectral subclasses of a given spectral class by looking at the spectral weight $`c_A`$, proceeding along steps similar to those followed above. * Let us suppose that, for a given graph $`𝒢`$ belonging to the spectral class characterized by $`\overline{d}`$, there exists at least one nonzero–measure subset $`AG`$ such that, as $`m^20`$, $$\mathrm{Sing}[C(m^2)]_Ac_A(m^2)^{\overline{d}/21}$$ (26) with $`c_Ac`$, with $`c`$ given as in eq. (18). * Then we say that a nonzero–measure subset $`MG`$, which certainly is maximal w.r.t. $`\overline{d}`$, due to its universality, is maximal also w.r.t. $`c_A`$ whenever 1. $`\|MA\|>0`$ 2. $`\mathrm{Sing}[C(m^2)]_Mc_M(m^2)^{\overline{d}/21}`$, with $`c_M=c_A`$. 3. For any nonzero–measure subset $`BM`$ we have $`c_B=c_M`$. 4. There exists no $`BM`$ such that $`c_B=c_M`$ and $`\|B\|>\|M\|`$. By this definition it follows that the set of all maximal subsets with respect to $`c_M`$ coincides with the equivalence class $`\{M\}`$ and we will call it the homogeneity subclass of spectral weight $`c_M`$. * We then prove the Theorem 2: The subgraphs $``$ and its complement $`\overline{}`$ are infinitely separated and define a spectral partition of $`𝒢`$. * Following the same steps as the previous section, we then consider a Gaussian model on the graph $``$ showing that, from the infinite separability of $``$ and $`\overline{}`$, the coefficient of $`\mathrm{Sing}[C(m^2)]_M`$ is $`c_M`$. Therefore we can say that $`c_M`$ is a universal property of the graph $``$ and defines a spectral subclass separated from the rest. Proof of Theorem 2: To prove this theorem we first need the following lemma: Lemma: Within a given spectral subclass, for any subset $`A`$ of the subclass, the asymptotic form of $`[C(m^2)]_A`$ is invariant under pre–averaging over any normalized point distribution with nonzero–measure support. In other words, if we define $$[C(m^2)]_{A,\alpha }=\frac{[\alpha C(m^2)]_A}{[\alpha ]_A}$$ (27) where $`\alpha _i>0`$ on a subset of $`A`$ with nonzero measure, then again $$\mathrm{Sing}[C(m^2)]_{A,\alpha }c_A(m^2)^{\overline{d}/21}$$ (28) with no dependence at all for $`c_A`$ and $`\overline{d}`$ on the distribution $`\alpha =\{\alpha _i;iA\}`$. The proof of this statement is elementary: we define the quantities $$f_i=(m^2)^{\overline{d}/2+1}C_{ii}(m^2)c_A$$ (29) Then, by construction, for any $`ϵ>0`$ there exist a $`\delta >0`$ such that we have $`\left\|[f]_A\right\|<ϵ`$ as soon as $`m^2<\delta `$. Hence we also have $$\left\|[\alpha f]_A\right\|<\left(\underset{iA}{sup}\alpha _i\right)\left\|[f]_A\right\|<\left(\underset{iA}{sup}\alpha _i\right)ϵ$$ (30) which immediately implies our assertion. Now we can prove Theorem 2: Let us suppose ad absurdum that $`\underset{¯}{D}(,\overline{})=1`$ and therefore that it exists a nonzero–measure subset $`\overline{M}^{}\overline{M}`$ such that $`\underset{¯}{D}(,\overline{}^{})=1`$. From the maximality of $`M`$ it follows that $`c_Mc_{\overline{M}^{}}`$. The following proof is given only for $`\overline{d}<4`$, owing to brevity and physical requirements. Indeed a real structure has necessarily a dimension $`\overline{d}3`$; moreover, from a purely theoretical point of view, the class of models we have in mind, with site variables and link interactions, typically have $`4`$ as an upper critical dimension for the scaling behaviour. Let us consider first the case of a spectral class where $`[C(m^2)]_G`$ diverges when $`m^20`$, that is such that $`\overline{d}<2`$. The Schwinger-Dyson equation for $`C_{ii}[m^2]`$ reads: $$(J_i+m^2\eta _i)C_{ii}(m^2)\underset{k𝒢}{}J_{ik}C_{ki}(m^2)=1$$ (31) Averaging equation (31) over $`M`$, we obtain the relation $$[JC]_M+m^2[\eta C]_M[JC]_M=\|M\|$$ (32) where $`(JC)_iJ_iC_{ii}`$, $`(\eta C)_i\eta _iC_{ii}`$ and $`(JC)_i=_kJ_{ik}C_{ki}`$. We then divide by $`[JC]_M`$ and let $`m^20`$. Due to the divergence of $`[JC]_M`$ we have that, for any $`ϵ>0`$ there exists a $`\delta >0`$ such that, as soon as $`m<\delta `$, $$1ϵ\frac{[JC]_M}{[JC]_M}$$ (33) Next we set $$J_{\overline{M}^{},i}=\underset{k\overline{M}^{}}{}J_{ik},(JC)_{\overline{M}^{},i}=\underset{k\overline{M}^{}}{}J_{ik}C_{ki}$$ (34) and use the positivity of $`C_{ii}C_{ik}`$ to push the above inequality to $$1ϵ1\frac{[J_{\overline{M}^{}}C]_M}{[JC]_M}+\frac{[(JC)_{\overline{M}^{}}]_M}{[JC]_M}$$ (35) which yields $$\underset{m^20}{lim}\frac{[(JC)_{\overline{M}^{}}]_M}{[J_{\overline{M}^{}}C]_M}=1$$ (36) Owing to the symmetry of $`\underset{¯}{D}(,\overline{}^{})`$, we may repeat the above steps with $`M`$ and $`\overline{M}^{}`$ interchanged. Since the symmetry of $`J_{ij}`$ and $`C_{ij}`$ implies $`[(JC)_{\overline{M}^{}}]_M=[(JC)_M]_{\overline{M}^{}}`$, we finally obtain $$\underset{m^20}{lim}\frac{[J_{\overline{M}^{}}C]_M}{[J_MC]_{\overline{M}^{}}}=1$$ (37) At this stage we apply the lemma given above with $`\alpha `$ identified with $`J_{\overline{M}^{}}`$ or $`J_M`$, namely $$[J_{\overline{M}^{}}C]_Mc_M[J_{\overline{M}^{}}]_M(m^2)^{\overline{d}/21},[J_MC]_{\overline{M}^{}}c_{\overline{M}^{}}[J_M]_{\overline{M}^{}}(m^2)^{\overline{d}/21}$$ (38) But $`[J_{\overline{M}^{}}]_M=[J_M]_{\overline{M}^{}}`$ so that eq. (37) implies $`c_M=c_{\overline{M}^{}}`$, contradicting our initial hypothesis that $`\underset{¯}{D}(,\overline{})=1`$ with $`M`$ maximal. Hence necessarily $`\underset{¯}{D}(,\overline{})=\mathrm{}`$, proving our assertion. Let us now consider a spectral class where $`C(m^2)_G`$ does not diverge in the limit $`m^20`$ while its first derivative with respect to $`m^2`$, $`C^{}(m^2)_G`$, diverges in the same limit. This is the case of a spectral class characterized by a spectral dimension $`2<\overline{d}<4`$, where: $$[C^{}(m^2)]_{M,\alpha }=\frac{[\alpha C^{}(m^2)]_M}{[\alpha ]_M}(\overline{d}/21)c_M(m^2)^{\overline{d}/22},m^20$$ (39) Taking the first derivative with respect to $`m^2`$ in the Schwinger-Dyson equation (31), we obtain: $$\eta _iC_{ii}(m^2)+m^2\eta _iC_{ii}^{}(m^2)=\underset{k𝒢}{}J_{ik}[C_{ki}^{}(m^2)C_{ii}^{}(m^2)]$$ (40) which can be averaged over $`M`$ giving $$[\eta C]_M+m^2[\eta C^{}]_M=[JC^{}]_M[JC^{}]_M$$ (41) Together with eq. (39), this implies $$\underset{m^20}{lim}(m^2)^{2\overline{d}/2}([JC^{}]_M[JC^{}]_M)=0^+$$ (42) that is, for any $`ϵ>0`$ there exists a $`\delta >0`$ such that, as soon as $`m^2<\delta `$ $$0<\xi ([JC^{}]_M[JC^{}]_M)<ϵ$$ (43) with $`\xi =(m^2)^{2\overline{d}/2}`$. This can be rewritten as $$0<[(JC^{})_M]_M[J_MC^{}]_M+[(JC^{})_{\overline{M}}]_M[J_{\overline{M}}C^{}]_M<\xi ^1ϵ$$ (44) Now, since $`C_{ij}^{}_k\eta _kC_{ik}C_{kj}`$ are the elements of a negative semi–definite matrix, one has that $`[(JC^{})_M]_M[J_MC^{}]_M>0`$. Therefore $$0[(JC^{})_{\overline{M}}]_M[J_{\overline{M}^{}}C^{}]_M<\xi ^1ϵ$$ (45) Again owing to the symmetry of $`\underset{¯}{D}(,\overline{}^{})`$, the previous steps can be repeated with $`M`$ and $`\overline{M}`$ interchanged, leading to: $$0[(JC^{})_M]_{\overline{M}}[J_MC^{}]_{\overline{M}^{}}<\xi ^1ϵ$$ (46) Since $`[(JC^{})_{\overline{M}}]_M=[(JC^{})_M]_{\overline{M}}`$, these two relations imply: $$0\|[J_{\overline{M}^{}}C^{}]_M[J_MC^{}]_{\overline{M}^{}}\|<\xi ^1ϵ$$ (47) Eq. (39) entails in the limit $`m^20`$: $$[J_{\overline{M}^{}}C^{}]_M(\overline{d}/21)c_M[J_{\overline{M}^{}}]_M\xi ^1,[J_MC^{}]_{\overline{M}^{}}(\overline{d}/21)c_{\overline{M}^{}}[J_M]_{\overline{M}^{}}\xi ^1$$ (48) so that, since $`[J_{\overline{M}^{}}]_M=[J_M]_{\overline{M}^{}}`$ from (47) one obtains $`c_M=c_{\overline{M}^{}}`$, which contradicts our hypothesis $`\underset{¯}{D}(,\overline{}^{})=1`$ and therefore proves our assertion $`\underset{¯}{D}(,\overline{})=\mathrm{}`$.
warning/0004/cond-mat0004381.html	ar5iv	text	# Spin-triplet superconductivity in repulsive Hubbard models with disconnected Fermi surfaces: a case study on triangular and honeycomb lattices ## Abstract We propose that spin-fluctuation-mediated spin-triplet superconductivity may be realized in repulsive Hubbard models with disconnected Fermi surfaces. The idea is confirmed for Hubbard models on triangular (dilute band filling) and honeycomb (near half-filling) lattices using fluctuation exchange approximation, where triplet pairing order parameter with f-wave symmetry is obtained. Possible relevance to real superconductors is suggested. A fascination toward spin-triplet superconductivity has a long history, but recent experimental suggestions for triplet pairing in a heavy fermion system UPt<sub>3</sub>, organic conductors (TMTSF)<sub>2</sub>X (X=ClO<sub>4</sub>,PF<sub>6</sub>), and a ruthenate Sr<sub>2</sub>RuO<sub>4</sub>, have renewed our interests in mechanisms of triplet superconductivity. In particular, it is fairly intriguing to investigate whether electron-electron repulsive interactions can lead to triplet superconductivity. Ferromagnetic-spin-fluctuation mechanism has been proposed from the early days, but to our knowledge, realization of triplet superconductivity (at sizable temperatures) has yet to be established theoretically for repulsive electron models with renormalization effects of the quasiparticles taken into account. The lifetime of the quasiparticles is important since this is a factor dominating $`T_c`$. Recently, the present authors with Aoki have investigated the possibility of triplet pairing in the Hubbard model for a variety of lattice structures and band fillings using fluctuation exchange (FLEX) approximation. A naive expectation is that triplet superconductivity may be realized when the band is away from half-filled and the density of states (DOS) at the Fermi level is large, since ferromagnetic fluctuations become strong in such a situation. In ref., however, it has turned out that the transition temperature $`(T_c)`$ of triplet superconductivity, if any, is too low to be detected as far as the cases surveyed there are concerned. A typical case is a square lattice with appreciable next nearest neighbor hoppings and dilute band fillings. First let us briefly review this situation as a reference for the results presented later. We consider the Hubbard model, $`=_{i,j\sigma =,}t_{ij}c_{i\sigma }^{}c_{j\sigma }+U_in_in_i,`$ on a square lattice shown in Fig.1, where $`t(=1)`$ is the nearest and $`t_1^{}(=t_2^{}`$ here) is the next nearest neighbor hopping. In the FLEX calculation, (i) Dyson’s equation is solved to obtain the renormalized Green’s function $`G(k)`$, where $`k(𝐤,iϵ_n)`$ denotes the 2D wave-vectors and the Matsubara frequencies, (ii) the effective electron-electron interaction $`V^{(1)}(q)`$ is calculated by collecting RPA-type bubbles and ladder diagrams consisting of the renormalized Green’s function, namely, by summing up powers of the irreducible susceptibility $`\chi _{\mathrm{irr}}(q)\frac{1}{N}_kG(k+q)G(k)`$ ($`N`$:number of $`k`$-point meshes), (iii) the self energy is obtained as $`\mathrm{\Sigma }(k)\frac{1}{N}_qG(kq)V^{(1)}(q)`$, which is substituted into Dyson’s equation in (i), and the self-consistent loops are repeated until convergence is attained. Throughout the study, we take $`64\times 64`$ $`k`$-point meshes and the Matsubara frequencies $`ϵ_n`$ from $`(2N_c1)\pi T`$ to $`(2N_c1)\pi T`$ with $`N_c`$ up to 16384 in order to ensure convergence at low temperatures. To obtain $`T_c`$, we solve the eigenvalue (Éliashberg) equation for the superconducting order parameter $`\varphi (k)`$, $`\lambda \varphi (k)`$ $`=`$ $`{\displaystyle \frac{T}{N}}{\displaystyle \underset{k^{}}{}}\varphi (k^{})\|G(k^{})\|^2V^{(2)}(kk^{}),`$ (1) where the pairing interaction $`V^{(2)}`$, which mediates pair scattering from $`(𝐤,𝐤)`$ to $`(𝐤^{},𝐤^{})[(𝐤+𝐪,𝐤𝐪)]`$, is given as $`V_s^{(2)}(q)={\displaystyle \frac{3}{2}}\left[{\displaystyle \frac{U^2\chi _{\mathrm{irr}}(q)}{1U\chi _{\mathrm{irr}}(q)}}\right]{\displaystyle \frac{1}{2}}\left[{\displaystyle \frac{U^2\chi _{\mathrm{irr}}(q)}{1+U\chi _{\mathrm{irr}}(q)}}\right]`$ (2) for singlet pairing, and $`V_t^{(2)}(q)={\displaystyle \frac{1}{2}}\left[{\displaystyle \frac{U^2\chi _{\mathrm{irr}}(q)}{1U\chi _{\mathrm{irr}}(q)}}\right]{\displaystyle \frac{1}{2}}\left[{\displaystyle \frac{U^2\chi _{\mathrm{irr}}(q)}{1+U\chi _{\mathrm{irr}}(q)}}\right]`$ (3) for triplet pairing. In either case, the first (second) term represents the contribution from spin (charge) fluctuations. $`T_c`$ is the temperature at which the maximum eigenvalue $`\lambda `$ reaches unity. We denote the eigenvalue and the order parameter for triplet (singlet) pairing as $`\lambda _t`$ ($`\lambda _s`$) and $`\varphi _t`$ $`(\varphi _s)`$, respectively. In ref., $`t^{}`$, $`U`$, and the band filling $`n`$ were varied in search of triplet superconductivity, but $`\lambda _t`$ remained below $`0.2`$ in the tractable temperature range as typically displayed in Fig.1 (dash-dotted line). The main reason why triplet pairing instability is weak is because $`\|V_t^{(2)}\|`$ is only one third of $`\|V_s^{(2)}\|`$ when spin fluctuation is dominant as can be seen from eqs.(2) and (3). In this Letter, we propose that the above difficulty for spin-fluctuation-mediated triplet pairing in the Hubbard model can be overcome under certain conditions. Let us first present our idea. We consider a situation (see Fig.2) where (i) the Fermi surface (FS) is disconnected (preferably well separated) into two pieces which are located point symmetrically about $`𝐤=\mathrm{𝟎}`$, and (ii) the spin structure is pronounced around a wave vector Q in such a way that two electrons with zero total momentum can be scattered within each piece of the FS (this process will be called intra-FS pair scattering hereafter) by exchanging spin fluctuations having momentum $`𝐐`$. Now, in order to have large $`\lambda `$, the quantity $`[_{𝐤,𝐤^{}FS}V^{(2)}(𝐤𝐤^{})\varphi (𝐤)\varphi (𝐤^{})]/[_{𝐤FS}\varphi ^2(𝐤)]`$ (the numerator being $`_{𝐤FS}V^{(2)}(𝐐)\varphi (𝐤)\varphi (𝐤+𝐐)`$) has to be positive and large. Then, pair scatterings from $`(𝐤,𝐤)`$ to $`(𝐤+𝐐,𝐤𝐐)`$ for singlet pairing, mediated by repulsive $`V_s^{(2)}(𝐐)`$, have to accompany a sign change in the order parameter $`\varphi _s(𝐤)`$ (Fig.2(a)). Hence the nodes of $`\varphi _s(𝐤)`$ must intersect the FS. For triplet pairing, by contrast, pairs can be scattered within a region having the same sign in $`\varphi _t(𝐤)`$ because $`V_t^{(2)}(𝐐)`$ is attractive. In this case, since the gap nodes (which exists due to triplet pairing symmetry $`\varphi (𝐤)=\varphi (𝐤)`$) do not intersect the FS (Fig.2(b)), the entire FS can be exploited for pairing, so that triplet pairing may be enhanced. Quite recently, a related proposal has been raised by Kohmoto and Sato for systems with both phonons and spin fluctuations present, as discussed later. The above conditions are not satisfied for the $`t`$-$`t^{}`$ square lattice because it has a connected FS. As an example of a system in which the above conditions are indeed satisfied, we consider the Hubbard model on an isotropic triangular lattice with dilute band fillings. The band dispersion for $`U=0`$ is given by $`\epsilon _{\mathrm{}}(𝐤)=2[\mathrm{cos}k_x+\mathrm{cos}k_y+\mathrm{cos}(k_x+k_y)]`$ when we represent an isotropic triangular lattice by setting $`t=t_1^{}=1`$ and $`t_2^{}=0`$ in Fig.1. Superconductivity on an isotropic triangular lattice has been studied by several authors, but their interest was mainly focused on $`n1`$. In ref, possibility of triplet superconductivity was studied at quarter filling $`(n=0.5)`$, where ferromagnetic fluctuations become strong because the Fermi level lies right at the position where the DOS diverges. However, $`\lambda _t`$ was again found to be small, which, in the present context, is because the FS is connected. If we set $`n<0.5`$, on the other hand, the FS is disconnected into two pieces, which are centered respectively at $`k=(2\pi /3,2\pi /3)`$ and $`k=(4\pi /3,4\pi /3)`$. Here we take $`n=0.15`$, where the two pieces of the FS are well separated. In Fig.3, we plot the FLEX result for $`\|G(𝐤,i\pi k_BT)\|^2`$ (a) and the spin susceptibility $`\chi (𝐤,0)\chi _{\mathrm{irr}}(𝐤,0)/[1U\chi _{\mathrm{irr}}(𝐤,0)]`$ (b) against $`𝐤`$ for $`U=8`$ and $`T=0.01`$. The FS as identified from the ridge in $`\|G(𝐤,i\pi k_BT)\|^2`$ is indeed disconnected into two pieces. $`\chi (𝐤,0)`$ is sharply peaked at $`𝐤=0`$ as seen in Fig.3(b), indicating ferromagnetic fluctuations. This is partially because the FS is small, but it is also because the Fermi level for $`n=0.15`$ is still not too far away from the peak position of the DOS. In this case, $`\lambda _s`$ is shown to be small, which is because $`V_s^{(2)}(𝐐)`$ can only mediate pair scatterings in the vicinity of the nodes when $`𝐐\mathrm{𝟎}`$. If we turn to triplet pairing, the order parameter $`\varphi _t(𝐤,i\pi k_BT)`$, plotted against $`𝐤`$ for $`T=0.01`$ in Fig.3(c), has f-wave (f$`_{x^33xy^2}`$-wave in the notation of the $`C_6`$ symmetry group) symmetry with three sets of nodal lines ($`k_x0(\mathrm{mod2}\pi ),k_y0,`$ and $`k_x+k_y0`$). Comparing Figs.3(a) and (c), we can see that these nodes do not intersect the FS. Accordingly, $`\lambda _t`$ (Fig.1,solid line) is strongly enhanced compared to the case for the $`t`$-$`t^{}`$ square lattice. A spline extrapolation to low temperatures suggests a possible low but finite $`T_c`$. As another example, we next propose that the Hubbard model on a honeycomb lattice (Fig.4) should also be interesting. Since there are two sites (A and B) in a unit cell, this is a two-band system. The noninteracting band dispersion $`\epsilon _{\mathrm{hc}}(𝐤)=\pm \sqrt{\epsilon _{\mathrm{}}(𝐤)+3}`$ has two pairs of vertex-sharing cones at $`k=(2\pi /3,2\pi /3)`$ and $`k=(4\pi /3,4\pi /3)`$, so again the FS becomes disconnected, this time for fillings close to $`n=1`$. In the multiband version of FLEX, the quantities $`G`$, $`\chi `$, $`\mathrm{\Sigma }`$, and $`\varphi `$ have $`2\times 2`$ matrix forms, e.g., $`G_{\alpha \beta }(𝐤,i\omega _n)`$, where $`\alpha ,\beta `$ denote A or B sites. The band representation of the Green’s function and the order parameters is obtained by using the relation between the annihilation operators of upper $`(u)`$ and lower $`(l)`$ band electrons ($`c_𝐤^u`$, $`c_𝐤^l`$) and those of A and B site electrons ($`c_𝐤^\mathrm{A}`$, $`c_𝐤^\mathrm{B}`$). As for $`\chi `$, we diagonalize the $`2\times 2`$ matrix $`\chi _{\alpha \beta }`$ to obtain $`\chi _\pm =(\chi _{\mathrm{AA}}+\chi _{\mathrm{BB}})/2\pm \sqrt{[(\chi _{\mathrm{AA}}\chi _{\mathrm{BB}})/2]^2+\|\chi _{\mathrm{AB}}\|^2}`$. In Fig.5, we plot $`\|G^l(𝐤,i\pi k_BT)\|^2`$ (a) and $`\chi _+(𝐤,0)`$ (b) for $`n=0.95`$, $`U=8`$ and $`T=0.01`$. Since $`\chi _{\mathrm{AB}}(\mathrm{𝟎},0)`$ is found to be negative, the peak around $`𝐤=\mathrm{𝟎}`$ in $`\chi _+(𝐤,0)`$ is an indication of antiferromagnetic fluctuations, as expected for a nearly half-filled bipartite lattice system. Note that $`\chi _+(𝐤,0)`$ has a broad structure compared to the case for the triangular lattice (Fig.3(b)). If we turn to $`\lambda _s`$ and $`\lambda _t`$ as functions of $`T`$ in Fig.6(a), $`\lambda _t`$ is again large, but this time $`\lambda _s`$ is in fact larger. We can trace this to the broad spin structure, for which spin fluctuations with relatively large momentum can be exchanged to mediate singlet pairing at wave vectors away from of the nodes. Nevertheless, we can still observe that $`\lambda _t`$ is enhanced above $`\lambda _s/3`$ (recall that $`\|V_t^{(2)}\|\|V_s^{(2)}\|/3`$), which should be due to the fact that the nodes in $`\varphi _s^l`$ intersect the FS, while those in $`\varphi _t^l`$ do not as seen by comparing Figs.5(a) and (c)/(d). We have found that $`\|\varphi _{\mathrm{AB}}\|>(<)\|\varphi _{\mathrm{AA}}\|`$ for singlet (triplet) pairing, meaning that singlet (triplet) pairing mainly takes place on different (same) sublattices. This is in fact consistent with the antiferromagnetic alignment of the spins. Then, we can intuitively expect that triplet can dominate over singlet if we introduce a level offset, $`\mathrm{\Delta }_{\mathrm{AB}}`$, between A and B sites. In Fig.6(b), we plot $`\lambda _s`$ and $`\lambda _t`$ for the same parameter values as in Fig.6(a), except for a finite $`\mathrm{\Delta }_{\mathrm{AB}}=5`$. Triplet pairing indeed dominates over singlet pairing at low temperatures, and here again a possible finite $`T_c`$ for triplet superconductivity is suggested. The intuitive picture can be paraphrased in the momentum space that the peak structure of $`\chi `$ (not shown) becomes sharper when $`\mathrm{\Delta }_{\mathrm{AB}}0`$. Finally, let us make some remarks concerning possible relevance to real superconductors. As for UPt<sub>3</sub>, if we examine the FS calculated by FLAPW method, we notice that there are two disconnected pockets (band 37 in ref.), which, in our view, is favorable for triplet pairing. It would be an interesting future problem to investigate in detail the applicability of the present mechanism. Disconnected FS can arise similarly in graphite intercalation compounds (GIC), except that the FS is cylindrical (quasi 2D). This is because graphite is a system where honeycomb sheets of carbon atoms are stacked. Although spin fluctuations in GIC may not be strong enough to induce superconductivity purely electronically, the disconnectivity of the FS itself should be favorable for triplet pairing, so a cooperation between certain phonon modes and (weak) spin fluctuations might lead to triplet superconductivity (even in the absence of $`\mathrm{\Delta }_{\mathrm{AB}}`$ considered above). Namely, if attractive intra-FS pair scatterings mediated by phonons are present, antiferromagnetic spin fluctuations as considered here would work constructively with phonons to enhance intra-FS pair scatterings for triplet pairing, while the converse is true for singlet pairing. Experimentally, although triplet pairing has not been claimed to our knowledge, a large value of $`H_{c2}`$ (extrapolated to $`T=0`$) observed in C<sub>4</sub>KTl<sub>1.5</sub> is in fact reminiscent of a large H<sub>c2</sub> in (TMTSF)<sub>2</sub>X. As for (TMTSF)<sub>2</sub>X and Sr<sub>2</sub>RuO<sub>4</sub>, Kohmoto and Sato have recently proposed that disconnectivity (quasi-one dimensionality) of the FS, along with the presence of spin fluctuations originating from the nesting of the FS, plays an essential role in stabilizing phonon-mediated triplet p-wave pairing. Our study is related to this proposal in that disconnectivity is important, but in these systems, as seen in Kohmoto and Sato’s argument, the dominant spin fluctuations have wave-vectors that bridge the two pieces of the FS, so that they mediate inter-FS pair scatterings rather than intra-FS ones. Thus, our purely spin-fluctuation-mediated pairing mechanism does not directly apply to these materials, although we do believe that the disconnectivity of the FS may be playing a certain role in realizing triplet superconductivity. We wish to thank Prof. H. Aoki for illuminating discussions. We also thank Prof. M. Kohmoto and Dr. M. Sato for fruitful discussions and sending us ref. prior to publication. K.K. thanks Dr. H. Kontani for discussions on multiband FLEX. Numerical calculations were performed at the Computer Center, University of Tokyo, and at the Supercomputer Center, ISSP, University of Tokyo. K.K. acknowledges support by the Grant-in-Aid for Scientific Research from the Ministry of Education of Japan, while R.A. acknowledges support by the JSPS Research Fellowships for Young Scientists.
warning/0004/hep-th0004148.html	ar5iv	text	# 1 Introduction ## 1 Introduction The AdS/CFT correspondence provides a powerful tool for studying $`𝒩=4`$ supersymmetric Yang–Mills theory in four dimensions at large $`N`$ and large ’t Hooft coupling. In particular, there exist precise prescriptions to calculate correlation functions, spectra of gauge invariant operators, Wilson loops and $`c`$-functions in supergravity. These data can be compared with field theory or provide non-trivial predictions for strongly coupled field theories. A natural question is whether this correspondence can be extended to theories with spontaneously or manifestly broken superconformal symmetry. Such theories arise either by giving vacuum expectation values to scalar fields , - or by deformations of the conformal theory with relevant operators -. These issues can all be treated efficiently in the context of five-dimensional gauged supergravity and the resulting backgrounds are kink-type solutions with four-dimensional Poincaré invariance which approach AdS asymptotically. A related question concerns the uplifting of these solutions to solutions of type-IIB supergravity/string theory, which can be quite involved, as we will see in this work. It is also of interest for the program of consistent truncations which, as yet, has not been completed in the case of $`𝒩=8`$ gauged supergravity in five dimensions. In this note we will study a supergravity dual of a particular deformation of the $`𝒩=4`$ SYM theory by a mass term that preserves $`𝒩=2`$ supersymmetry . This is simply $`𝒩=2`$ supersymmetric gauge theory with gauge group $`SU(N)`$ coupled to a massive hypermultiplet in the adjoint representation of the gauge group. One can think of this model also as $`𝒩=1`$ supersymmetric QCD with three chiral multiplets $`\mathrm{\Phi }_{i=1,2,3}`$ in the adjoint, out of which, one is massless and the other two have equal masses. Other choices of the massterms are of course possible and lead to models with $`𝒩=1`$ supersymmetry. Such models have been studied previously in the context of the AdS/CFT correspondence . The outline of this paper is as follows: in section 2 we present the background of gauged supergravity that is dual to $`𝒩=4`$ SYM deformed by a mass term for one hypermultiplet. We find a family of solutions that is parametrized by one real constant $`c`$, whose value determines the physics and study fluctuation in this background. For $`c0`$ it turns out that we describe the strong coupling regime of part of the Coulomb branch of an $`𝒩=2`$ theory, whereas flows with $`c>0`$ are unphysical. In section 3, we compute the uplifted metric in ten dimensions. With the help of this type-IIB background we calculate expectation values of Wilson loops that correspond to the potential of an external heavy quark-antiquark pair. Finally, in section 4 we present some concluding remarks and comment on the singularities that are typical for backgrounds of non-conformal theories. Note added In the final stages of our work, the paper appeared that has considerable overlap with ours. In addition to the uplifted ten-dimensional metric, that we also computed independently, these authors presented the axion/dilaton and the complex two-form in ten dimensions. A numerical investigation has also appeared before in . ## 2 A dual of $`𝒩=2`$ supersymmetric gauge theories Our starting point is the action of five-dimensional gauged supergravity $$S=d^5x\left\{\frac{1}{4}\frac{1}{2}\underset{I}{}_\mu \alpha _I^\mu \alpha _IP(\alpha )\right\},$$ (1) where we have chosen canonical kinetic terms for the scalars. This is possible only in certain sub-sectors of the full $`E_{6(6)}/USp(8)`$ coset space sigma-model parameterizing the 42 scalars of the theory. We have also set to zero all other tensor fields. For the applications we have in mind we need the scalars that correspond to dimension 2 and 3 operators, which are in the $`\mathrm{𝟐𝟎}`$ and $`\mathrm{𝟏𝟎}`$ representation of $`SO(6)SU(4)`$, respectively. The massless 5-dim dilaton and axion will not play a rôle and are constant. They correspond to exactly marginal deformations in the gauge theory. Furthermore, we are interested in solutions of (1) that preserve part of the supersymmetries. For such supersymmetric flows it is known that the scalar potential $`P`$ can be written in terms of a superpotential, which we will denote as $`W(\alpha )`$: $$P(\alpha )=\frac{1}{8}\underset{I}{}\left(\frac{W}{\alpha _I}\right)^2\frac{1}{3}W^2.$$ (2) Furthermore, we demand that the solutions of (1) preserve four-dimensional Poincaré invariance along the brane directions, and that the metric asymptotes the $`AdS_5`$ space-time near the boundary which we take to be at $`r\mathrm{}`$. For the metric we make the ansatz $$ds_5^2=e^{2A(r)}dx_{\|\|}^2+dr^2,$$ (3) or, if we need the metric in its conformally flat version then $$ds_5^2=e^{2A(z)}\left(dx_{\|\|}^2+dz^2\right),$$ (4) where the relation between the different coordinate choices is $`dr=e^Adz`$. The boundary at $`r=\mathrm{}`$ corresponds to $`z=0`$. On the supergravity side two scalars, denoted by $`\alpha _2`$ and $`\alpha _3`$, are involved and belong to the $`\mathrm{𝟐𝟎}`$ and to the $`\mathrm{𝟏𝟎}+\overline{\mathrm{𝟏𝟎}}`$ representation of $`SO(6)`$, respectively. They are dual to operator-bilinears in scalars and fermions as: $`\alpha _2:𝒪_2=\mathrm{Tr}(\overline{Z}_1Z_1+\overline{Z}_2Z_22\overline{Z}_3Z_3),`$ $`\alpha _3:𝒪_3=\mathrm{Tr}(\lambda _1\lambda _1+\lambda _2\lambda _2+\mathrm{}+h.c.),`$ (5) where the $`Z_{i=1,2,3}`$ denote the complex scalar components of the chiral superfields $`\mathrm{\Phi }_i`$, the $`\lambda _{i=1,2,3}`$ are the corresponding fermionic components, and the $`\mathrm{}`$ denote scalar trilinear terms. The non-vanishing scalar $`\alpha _2`$ reduces the gauge group to $`SU(2)\times SU(2)\times U(1)SU(4)`$ under which the scalars in the $`\mathrm{𝟔}`$ and the fermions in the $`\mathrm{𝟒}`$ decompose as: $`\mathrm{𝟔}`$ $``$ $`(1,1)_2+(1,1)_2+(2,2)_0,`$ $`\mathrm{𝟒}`$ $``$ $`(2,1)_1+(1,2)_1.`$ (6) We see that the $`U(1)`$ is to be identified with the $`U(1)_R`$ symmetry of the field theory, since the two scalars in the vector multiplet correspond to the two $`SU(2)`$ singlets in the $`\mathrm{𝟔}`$ which have charge $`\pm 2`$ under the $`U(1)`$ as in field theory, whereas the scalars of the hypermultiplet are singlets. The scalar $`\alpha _3`$ breaks one of the $`SU(2)`$’s to $`U(1)`$, and the unbroken one can be identified with the $`SU(2)_R`$ symmetry of the gauge theory under which the scalars in the vector multiplet are singlets and the scalars in the hypermultiplet are doublets. Furthermore, the two fermions in the hypermultiplet are singlets under $`SU(2)_R`$ and have $`U(1)_R`$ charge $`1`$, the two fermions in the vector multiplet are $`SU(2)_R`$ doublets and their $`U(1)_R`$ charge is $`+1`$. So in our decomposition the second $`SU(2)`$ is broken to $`U(1)`$. The potential in the scalar field space is obtained by a truncation of the four scalar potential that has been computed in and reads $`P`$ $`=`$ $`{\displaystyle \frac{1}{16}}e^{\frac{4}{\sqrt{6}}\alpha _2}\left(48e^{\sqrt{6}\alpha _2}\mathrm{cosh}\sqrt{2}\alpha _3+e^{2\sqrt{6}\alpha _2}\mathrm{sinh}^2\sqrt{2}\alpha _3\right),`$ $`W`$ $`=`$ $`e^{2\alpha _2/\sqrt{6}}{\displaystyle \frac{1}{2}}e^{4\alpha _2/\sqrt{6}}\mathrm{cosh}\sqrt{2}\alpha _3,`$ (7) where $`W`$ is the corresponding superpotential (cf. (2)). For a supersymmetric flow the equations of motion of (1) reduce to a set of first order equations: $`\dot{\alpha }_2`$ $`=`$ $`{\displaystyle \frac{1}{2}}_{\alpha _2}W={\displaystyle \frac{1}{\sqrt{6}}}e^{2\alpha _2/\sqrt{6}}{\displaystyle \frac{1}{\sqrt{6}}}e^{4a_2/\sqrt{6}}\mathrm{cosh}\sqrt{2}\alpha _3,`$ $`\dot{\alpha }_3`$ $`=`$ $`{\displaystyle \frac{1}{2}}_{\alpha _3}W={\displaystyle \frac{1}{2\sqrt{2}}}e^{4\alpha _2/\sqrt{6}}\mathrm{sinh}\sqrt{2}\alpha _3,`$ (8) $`\dot{A}`$ $`=`$ $`{\displaystyle \frac{1}{3}}W,`$ where the dot denotes the derivative with respect to the variable $`r`$ introduced in (3). As a consistency check, we note that we may set the scalar $`\alpha _3`$ to zero, as it is obvious from (2). Then, the remaining scalar describes the part of the Coulomb branch of the $`𝒩=4`$ SYM theory corresponding to the background of D3-branes distributed on a disc and $`SO(4)\times U(1)`$ symmetry.<sup>1</sup><sup>1</sup>1There should be a generalization of our discussion so far that includes a third scalar that further breaks $`SU(2)\times U(1)\times U(1)`$ to $`SU(2)\times U(1)`$. When the scalar in the $`\mathrm{𝟏𝟎}`$ is turned off this should be the part of the Coulomb branch of the $`𝒩=4`$ SYM theory corresponding to the background of D3-branes distributed on a ellipsoid, which has an $`SO(4)`$ symmetry. This solution was constructed in . The systems of equations in (2) form a rather complicated set of coupled first order differential equations, which, at first sight, do not have a solution that can be written down in a closed form. This is true for the coordinate $`r`$, but if we take the scalar $`\alpha _3`$ as a new radial coordinate<sup>2</sup><sup>2</sup>2From (12) and the remark after (4) it is clear that $`\alpha _3`$ is a monotonously decreasing function of $`r`$. a solution can be given in terms of elementary functions. Note that the equations in (2) have a $`𝐙_2`$ symmetry which transforms $`\alpha _3\alpha _3`$, so that we may restrict our analysis to $`\alpha _30`$ without loss of generality. For the scalar field $`\alpha _3`$ we find $$e^{\sqrt{6}\alpha _2}=\mathrm{sinh}^2\sqrt{2}\alpha _3\left[c+\mathrm{ln}\mathrm{tanh}\left(\frac{\alpha _3}{\sqrt{2}}\right)\right]+\mathrm{cosh}\sqrt{2}\alpha _3,$$ (9) where $`c`$ is an integration constant. As we will see the physical picture depends crucially on whether $`c`$ is positive, negative or zero. For the 5-dim metric we find $$ds_5^2=\frac{1}{\mathrm{sinh}^2\sqrt{2}\alpha _3}\left(2^{2/3}e^{4\alpha _2/\sqrt{6}}dx_{\|\|}^2+8e^{8\alpha _2/\sqrt{6}}d\alpha _3^2\right).$$ (10) The transition to the conformally flat metric (4) with conformal factor $$e^{2A(z)}=\frac{2^{2/3}e^{4\alpha _2/\sqrt{6}}}{\mathrm{sinh}^2\sqrt{2}\alpha _3},$$ (11) is given by the relation of the differentials $$dz=2^{7/6}e^{\sqrt{6}\alpha _2}d\alpha _3.$$ (12) We note that, for generic values of the constant $`c`$, the explicit dependence of $`\alpha _3`$ on $`z`$ cannot be obtained from (12), since the corresponding integral cannot be explicitly evaluated. Let us consider the behaviour near the boundary where the background becomes $`AdS_5`$. The boundary is approached as $`\alpha _30`$ which in turn, using (9), forces $`\alpha _20`$ as well. In this limit we find $`z2^{7/6}\alpha _3`$. Furthermore, we want to check that the scalars behave as expected from the AdS/CFT correspondence. The scalar $`\alpha _3`$ is dual to the fermion bilinear operator $`𝒪_3`$. Our field theory is obtained as a mass deformation of the $`𝒩=4`$ theory and therefore the solution near the horizon should behave as a pure mass term, i.e. $`\alpha _3z^1`$ and there should be no $`z^3`$ term corresponding to a gluino condensate. On the other hand, the behaviour of the scalar $`\alpha _2`$ is a beautiful example of the double rôle that the scalars play in the AdS/CFT correspondence — they can appear as deformations of the conformal theory, or they parameterize states in the theory. Here we expect both since we have to give mass to the scalar components of the chiral multiplets $`\mathrm{\Phi }_{1,2}`$ as dictated by supersymmetry, and the gauge theory has a Coulomb branch parametrized by the vev of $`Z_3`$, i.e. $`\alpha _2`$ should have contributions of the type $`z^2`$ and $`z^2\mathrm{log}z`$. From our explicit solution we immediately find that $`\alpha _3`$ $``$ $`z`$ $`\alpha _2`$ $``$ $`{\displaystyle \frac{1\mathrm{ln}2+2c}{\sqrt{6}}}z^2+{\displaystyle \frac{2}{\sqrt{6}}}z^2\mathrm{ln}z.`$ (13) From these expressions we see that the correct terms appear, as expected, and that there is a one parameter family of solutions labelled by $`c`$.<sup>3</sup><sup>3</sup>3As was noted in , there is a critical line in the space of solutions corresponding to $`c=0`$, in our notation. In the notation of , this gives $`b_{\mathrm{cr}}=(1\mathrm{ln}2)/\sqrt{6}0.125`$, thus confirming the critical value found numerically in . In this paper we will also be interested in the spectrum corresponding to the massless scalar equation in a background metric of the type (4) and with a plane wave ansatz along the flat coordinates denoted by $`x_{\|\|}`$. Within the AdS/CFT correspondence the solutions and eigenvalues of the massless scalar equation have been associated, on the gauge theory side, with the spectrum of the operator $`\mathrm{Tr}F^2`$, whereas those of the graviton fluctuations polarized in the directions parallel to the brane, with the energy momentum tensor $`T_{\mu \nu }`$ . Though a priori not to be expected, these two spectra and the corresponding eigenfunctions coincide (for a more recent discussion see ). It is well known that the entire analysis can be cast into an equivalent Schrödinger problem with potential $$V=\frac{9}{4}A_{}^{}{}_{}{}^{2}+\frac{3}{2}A^{\prime \prime }$$ (14) and eigenvalue equal to the mass squared ($`M^2`$). Since, as already noted, for generic values of the constant $`c`$ we cannot find the explicit dependence of $`\alpha _3`$ on $`z`$, it is not possible to explicitly evaluate the potential using (11). However, one easily can show that it is a monotonously decreasing function of $`z`$. As in the potential (14) has the same form as the potentials appearing in supersymmetric quantum mechanics and, therefore, the spectrum is bounded from below. Its supersymmetric partner is given by an equation similar to (14), but with a relative minus sign between the two terms. ### 2.1 $`c<0`$ In the following we will study the case $`c<0`$ in more detail. In particular, we will determine the spectrum of fluctuations of a minimal scalar in this background. First, note that $`e^{\sqrt{6}\alpha _2}`$ is a monotonically decreasing function of $`\alpha _3`$ and, as it turns out, it has only a finite range (see fig. 1) $$0\alpha _3\alpha _3^{\mathrm{max}}.$$ (15) In certain limits $`\alpha _3^{\mathrm{max}}`$ can be found analytically: $$\alpha _3^{\mathrm{max}}=\{\begin{array}{cc}\frac{1}{3\sqrt{2}}\mathrm{ln}\left(\frac{3c}{16}\right),\hfill & c0^{},\hfill \\ \frac{1}{\sqrt{2c}},\hfill & c\mathrm{}.\hfill \end{array}$$ (16) Using (9) and (12) we find that $$e^{\sqrt{6}\alpha _2}\frac{2\sqrt{2}}{\mathrm{sinh}\sqrt{2}\alpha _3^{\mathrm{max}}}\left(\alpha _3^{\mathrm{max}}\alpha _3\right),\mathrm{as}\alpha _3\alpha _3^{\mathrm{max}}$$ (17) and that $$z2^{1/3}\mathrm{sinh}\sqrt{2}\alpha _3^{\mathrm{max}}\mathrm{ln}\left(\alpha _3^{\mathrm{max}}\alpha _3\right),\mathrm{as}\alpha _3\alpha _3^{\mathrm{max}}.$$ (18) Therefore $`z`$ takes values in the whole semi-infinite real line, i.e. $`0z<\mathrm{}`$. From (11), (17) and (18) we find that $$A\frac{2^{1/3}z}{3\mathrm{sinh}\sqrt{2}\alpha _3^{\mathrm{max}}},\mathrm{as}z\mathrm{}.$$ (19) From (14) one concludes that the mass spectrum is continuous with a mass gap. The value of the latter is found by evaluating the potential for $`z\mathrm{}`$. The result is $$M_{\mathrm{gap}}^2=\frac{1}{2^{4/3}\mathrm{sinh}^2\sqrt{2}\alpha _3^{\mathrm{max}}},$$ (20) and in the two limits corresponding to (16) $$M_{\mathrm{gap}}^2=\{\begin{array}{cc}\frac{1}{4}\left(3c\right)^{2/3},\hfill & c0^{},\hfill \\ 2^{4/3}(c),\hfill & c\mathrm{}.\hfill \end{array}$$ (21) Hence, the mass gap vanishes as $`c`$ tends to zero and grows large as $`c`$ approaches large negative values. It is also worth noting that certain solutions corresponding to points on the Coulomb branch of $`𝒩=4`$ SYM, which were obtained previously in the literature, can be obtained in special limits of our solution. The limit $`c\mathrm{}`$ corresponds to a distribution of D3-branes on a disc of radius 2, in our units. In order to see that let us change variables as $$\alpha _3=\frac{1}{\sqrt{2c}}\mathrm{tanh}\left(z/2\right),$$ (22) and then send $`c\mathrm{}`$. As we can see from (16), this effectively shrinks the range of $`\alpha _3`$ to zero. After some elementary algebra we find $`ds_5^2`$ $`=`$ $`{\displaystyle \frac{\mathrm{cosh}^{2/3}(z/2)}{\mathrm{sinh}^2(z/2)}}\left(dz^2+2^{2/3}(c)dx_{\|\|}^2\right),`$ $`e^{\sqrt{6}\alpha _2}`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{cosh}^2(z/2)}}.`$ (23) This model has been studied on its own and also exhibits a mass gap. In fact in order to compute it using our formulae we first re-scale the energies by the factor $`2^{2/3}(c)`$ as it is evident from the expression for the metric in (23). Then, the second line in (21) gives $`M_{\mathrm{gap}}^2=1/4`$ which is in agreement with the result for the mass gap found in , when it is expressed in our units. We finally note that, from a physical point of view a continuous spectrum with a mass gap should be associated with a complete screening of charges in a quark-antiquark pair at a finite separation distance (inversely proportional to $`M_{\mathrm{gap}}`$). This will be further discussed in subsection 3.1. ### 2.2 $`c>0`$ Contrary to the previous case, $`e^{\sqrt{6}\alpha _2}`$ is now not a monotonic function (see fig. 1). For increasing $`\alpha _3`$ it first decreases, takes a minimum at $`\alpha _3=\alpha _3^{\mathrm{min}}`$ with $$e^{\sqrt{6}\alpha _2}\|_{\mathrm{min}}=\frac{1}{\mathrm{cosh}\alpha _3^{\mathrm{min}}},$$ (24) and then increases monotonically. The coordinate $`z`$ on the other hand can easily be seen to have finite range: $$0zz_0=2^{7/6}_0^{\mathrm{}}e^{\sqrt{6}\alpha _2}𝑑\alpha _3.$$ (25) In certain limits one can work out the value for $`\alpha _3^{\mathrm{min}}`$: $$\alpha _3^{\mathrm{min}}\{\begin{array}{cc}\frac{1}{6\sqrt{2}}\mathrm{ln}\frac{8}{3c},\hfill & c0^+,\hfill \\ \sqrt{2}e^{c1},\hfill & c+\mathrm{}.\hfill \end{array}$$ (26) We also obtain from (9) and (12) that $$e^{\sqrt{6}\alpha _2}\frac{c}{4}e^{2\sqrt{2}\alpha _3},\mathrm{as}\alpha _3\mathrm{}$$ (27) and that $$e^{2\sqrt{2}\alpha _3}2^{5/3}c(z_0z),\mathrm{as}zz_0^{}.$$ (28) Then from (11), (27) and (28) the warp factor of the metric becomes $$e^A2^{7/18}c^{1/6}(z_0z)^{1/6},\mathrm{as}zz_0^{},$$ (29) which corresponds to a naked singularity. It turns out to be a bad singularity, but we will return to this point later and discuss it in more detail. The spectrum of scalar fluctuations is now discrete since the range of $`z`$ is finite. Knowing the warp factor we can easily determine the Schrödinger potential in the limit $`zz_0^{}`$. The result is: $$V\frac{3}{16}\frac{1}{(zz_0)^2},\mathrm{as}zz_0^{}.$$ (30) A WKB calculation yields for the spectrum $$M_n^2=\frac{\pi ^2}{z_0^2}m(m+1)+𝒪(m^0),n=1,2,\mathrm{}$$ (31) For $`c\mathrm{}`$ the geometry reduces to that of a distribution of D3-branes smeared on a three-sphere of radius 2 in our units. In order to illustrate that, similarly to before, we first change variables as $$\alpha _3=\frac{1}{\sqrt{2c}}\mathrm{tan}\left(z/2\right),$$ (32) and then let $`c\mathrm{}`$ with the result $`ds_5^2`$ $`=`$ $`{\displaystyle \frac{\mathrm{cos}^{2/3}(z/2)}{\mathrm{sin}^2(z/2)}}\left(dz^2+2^{2/3}cdx_{\|\|}^2\right),`$ $`e^{\sqrt{6}\alpha _2}`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{cos}^2(z/2)}}.`$ (33) In the limit $`c\mathrm{}`$ we find using (25) that the maximum value of $`z`$ behaves as $`z_0\pi /(2^{1/3}c^{1/2})`$. Then, after rescaling of the energies by the factor $`2^{2/3}c`$, the WKB mass spectrum is just $`M_{\mathrm{gap}}^2=m(m+1)`$ and is in fact the exact result (as noted in ), in the units that we are using. ### 2.3 $`c=0`$ In this case the ranges of $`z`$ and $`\alpha _3`$ are infinite and $`e^{\sqrt{6}\alpha _2}`$ is a monotonously decreasing function (see fig. 1). In particular, from (9) and (12) we find that $$e^{\sqrt{6}\alpha _2}\frac{4}{3}e^{\sqrt{2}\alpha _3},\mathrm{as}\alpha _3\mathrm{}$$ (34) and that $$e^{\sqrt{2}\alpha _3}\frac{2^{4/3}}{3}z,\mathrm{as}z\mathrm{}.$$ (35) Hence, $$e^Az^{4/3},\mathrm{as}z\mathrm{}$$ (36) and the Schrödinger potential becomes $`V6/z^2`$ for $`z\mathrm{}`$. Therefore the spectrum is continuous and has no gap. Accordingly, we expect that the screening of charges in a quark-antiquark pair will be perfect only at an infinite separation. This will be confirmed in subsection 3.1. ## 3 The lift to ten dimensions Recently, it has been shown how general solutions of 5-dim gauged supergravity can be uplifted to type-IIB supergravity. This is part of the program to prove that a consistent truncation exists. This has not been shown for all fields, but for the metric the full non-linear KK ansatz has been conjectured and it has passed several non-trivial tests. The inverse of the deformed metric on $`S^5`$ is given in term of the $`27`$-bein $`𝒱`$ which has global $`E_{6(6)}`$ and local $`USp(8)`$ indices. In the $`SL(6,𝐑)\times SL(2,𝐑)`$ basis the vielbein is decomposed as $`𝒱(𝒱^{IJab},𝒱_{I\alpha }^{}{}_{}{}^{ab})`$ in terms of which the inverse metric becomes $$\widehat{g}^{mn}=\mathrm{\Delta }^{2/3}g^{mn}=2K_{IJ}^mK_{KL}^n\stackrel{~}{𝒱}_{IJab}\stackrel{~}{𝒱}_{KLcd}\mathrm{\Omega }^{ac}\mathrm{\Omega }^{bd},$$ (37) with $`\stackrel{~}{𝒱}`$ being the inverse vielbein and $`\mathrm{\Delta }^2=\mathrm{det}(g_{mn})/\mathrm{det}(g_{mn}^{(0)})`$, where $`(g_{mn}^{(0)})`$ is the undeformed five-sphere metric. Consequently, the ten-dimensional metric takes the form of a warped product space $$ds_{10}^2=\mathrm{\Delta }^{2/3}ds_5^2+g_{mn}dy^mdy^n=\mathrm{\Delta }^{2/3}\left(ds_5^2+\widehat{g}_{mn}dy^mdy^n\right).$$ (38) Then using the parameterization presented in we compute the 5-dim metric to be $`d\widehat{s}^2`$ $`=`$ $`{\displaystyle \frac{e^{\frac{2}{\sqrt{6}}\alpha _2}}{\mathrm{cosh}\sqrt{2}\alpha _3}}d\theta ^2+{\displaystyle \frac{e^{\frac{8}{\sqrt{6}}\alpha _2}\mathrm{sin}^2\theta }{\mathrm{\Delta }_1}}d\varphi _1^2`$ (39) $`+e^{\frac{2}{\sqrt{6}}\alpha _2}\mathrm{cos}^2\theta \left({\displaystyle \frac{1}{\mathrm{\Delta }_1\mathrm{cosh}\sqrt{2}\alpha _3}}\sigma _3^2+{\displaystyle \frac{1}{\mathrm{\Delta }_2}}(\sigma _1^2+\sigma _2^2)\right),`$ where $`\mathrm{\Delta }_1`$ $`=`$ $`e^{\sqrt{6}\alpha _2}\mathrm{cos}^2\theta \mathrm{cosh}\sqrt{2}\alpha _3+\mathrm{sin}^2\theta ,`$ $`\mathrm{\Delta }_2`$ $`=`$ $`e^{\sqrt{6}\alpha _2}\mathrm{cos}^2\theta +\mathrm{sin}^2\theta \mathrm{cosh}\sqrt{2}\alpha _3.`$ (40) The Maurer–Cartan forms $`\sigma _i`$, $`i=1,2,3`$ for $`SU(2)`$ are defined to obey $`d\sigma _i=ϵ_{ijk}\sigma _j\sigma _k`$. A convenient parameterization in terms of three Euler angles $`\varphi _2`$, $`\varphi _3`$ and $`\psi `$ is given by $`\sigma _1`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(\mathrm{cos}\varphi _2d\psi +\mathrm{sin}\varphi _2\mathrm{sin}\psi d\varphi _3\right),`$ $`\sigma _2`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(\mathrm{sin}\varphi _2d\psi +\mathrm{cos}\varphi _2\mathrm{sin}\psi d\varphi _3\right),`$ (41) $`\sigma _3`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(d\varphi _2+\mathrm{cos}\psi d\varphi _3\right).`$ Finally, the warp factor in (38) is easily computed to be $$\mathrm{\Delta }^{2/3}=e^{\frac{2}{\sqrt{6}}\alpha _2}\left(\mathrm{\Delta }_1\mathrm{\Delta }_2\mathrm{cosh}\sqrt{2}\alpha _3\right)^{1/4}.$$ (42) The metric (39) is manifestly $`SU(2)`$ invariant being written in terms of the corresponding Maurer–Cartan forms. In addition, there is an $`U(1)\times U(1)`$ invariance corresponding to the isometry with respect to the commuting Killing vectors $`/\varphi _1`$ and $`/\varphi _2`$. ### 3.1 Wilson loops and screening We now calculate string probes in the ten-dimensional background (39) that correspond to Wilson loops in field theory associated with the potential between an external quark-antiquark pair . The relevant probe action is that of a fundamental string which involves only the string frame metric and the NS $`B`$-field. As in a similar computation in , involving a rotating D3-brane, we consider trajectories with constant values for the angles $`\theta `$ and $`\varphi _1`$, as well as for the Euler angles entering into the definition of the Maurer–Cartan forms. For such constant values the pull-back of the $`B`$-field, which was recently calculated , vanishes, so that the action consists only of the Nambu–Goto term. Since there is an explicit dependence of the metric on $`\theta `$, consistency requires that the variation of this action with respect to $`\theta `$ is zero, in order for the equations of motion to be obeyed. It is easy to see that this procedure allows $`\theta =0`$ or $`\theta =\pi /2`$. The action can be written as $$S_{NG}=\frac{T}{2\pi \alpha ^{}}𝑑x\sqrt{g(\alpha _3)\left(\frac{d\alpha _3}{dx}\right)^2+f(\alpha _3)},$$ (43) with $`T`$ coming from the trivial integration in the Euclidean-time direction and $`x`$ is one of the spatial directions along the brane. The functions $`g(\alpha _3,\theta )=g_{\tau \tau }g_{\alpha _3\alpha _3}`$ , $`f(\alpha _3,\theta )=g_{\tau \tau }g_{xx}`$, where eventually only the values $`\theta =0`$ and $`\theta =\pi /2`$ are allowed. Following standard procedures one can find integral expressions for the separation of the two sources and the potential as function of the maximal distance of the string from the boundary. For small separations we find, as expected, a Coulomb potential $`V_{q\overline{q}}1/L`$. However, for larger values of $`L`$ we have to distinguish between the cases of negative and zero $`c`$. Let us begin with the $`c<0`$ case: as we increase $`L`$ the energy increases until it vanishes at a certain maximal length. Beyond that there doesn’t exist a geodesic connecting the two sources, the configuration will be that of two disconnected straight strings and the potential is zero. We can interpret this as complete screening with the maximum length given by $$L_{\mathrm{max}}=\{\begin{array}{cc}\frac{\pi }{M_{\mathrm{gap}}},\hfill & \theta =0,\hfill \\ \frac{\pi }{2M_{\mathrm{gap}}},\hfill & \theta =\frac{\pi }{2},\hfill \end{array}$$ (44) where the mass gap $`M_{\mathrm{gap}}`$ is given by (20). This behaviour of Wilson loops is precisely the one found in the context of the Coulomb branch of theories with sixteen supercharges for D3-branes distributed uniformly over a disc. For $`c=0`$ the behaviour is quite different. There doesn’t exist a finite screening length and the potential vanishes only at an infinite separation. For a short separation the potential exhibits the usual Coulombic behaviour, but we can also work out the behaviour for a very large separation. The result depends on the direction on the sphere: for $`\theta =0`$ we find again $`V_{q\overline{q}}1/L`$ which is reminiscent of a conformal theory, although the metric is not that for the $`AdS_5`$ space-time. For $`\theta =\pi /2`$ the potential is screened, but it still exhibits a power-law behaviour as $`V_{q\overline{q}}1/L^2`$. This is consistent with the fact that there is a ring-type singularity for the metric at $`\theta =\pi /2`$ as for the enhancon in and for a uniform distribution of D3-branes in a ring in . For the latter case the Wilson loop was analyzed in and also exhibits a $`1/L^2`$ behaviour for large separations. Note also that, for the qualitative features that we have presented for the Wilson loops the value of the 10-dim dilaton computed in plays no rôle except for the case with $`c=0`$ and $`\theta =\pi /2`$. Then the inclusion of the dilaton factor is indeed necessary in order to obtain the $`1/L^2`$ behaviour that we have mentioned. ## 4 Concluding remarks The 5-dim backgrounds that we described have naked singularities in the interior (IR in field theory). Hence, the question arises whether the physics or more specifically string theory is singular or not in such backgrounds. Well known examples of singular geometries that are resolved by string theory are orbifolds, orientifolds and conifolds. More recently singular geometries appeared in the supergravity duals of non-conformal theories. Examples are dilatonic branes, backgrounds dual to theories with sixteen supercharges on the Coulomb branch, for which the singularity is the source of a distribution of branes, and duals of theories with eight supercharges where the singularity is removed by a mechanism explained in . We think that our geometries are similar in nature to the examples of the Coulomb branch with sixteen supercharges. Furthermore, the distinction between good and bad singularities should be consistent with the AdS/CFT correspondence. This means that only a well-defined deformation and vacuum-state in field theory is dual to a geometry with a good singularity. We have encountered this issue in our paper. For $`c>0`$ we seemingly try to give a vev to a massive scalar field in the field theory. However, since we have turned on a non-zero source (mass term) for these scalars, its vev has to be fixed uniquely. Therefore, the case $`c>0`$ should correspond to a bad singularity, and $`c0`$, on the other hand, should be an acceptable one. Unfortunately, we do not know enough about string theory in such backgrounds to answer the question whether the singularity is physical or not. Instead, we can ask a somewhat simpler question: is the propagation of a quantum test particle well defined in the presence of the singularity? This criterion is identical to finding a unique self-adjoint extension of the wave-operator. This can be quickly answered by looking at the two solutions of the wave-operator locally near the singularity. The relevant part of the wave equation is $$\frac{d}{dr}\left(\sqrt{g}g^{rr}\frac{d\psi }{dr}\right)=0$$ (45) and the norm is $`\frac{q^2}{2}\sqrt{g}g^{tt}\psi ^{}\psi +\frac{1}{2}\sqrt{g}g^{ij}D_i\psi ^{}D_j\psi `$. This is the Sobolev norm and it is bounded from above by a constant times the energy of the fluctuation . This is a physical sensible norm because it guarantees that the backreaction of the fluctuation is small. For $`c0`$ only one solution is normalizable and the non-normalizable is discarded. Therefore, there exists a unique self-adjoint extension and the singularity is wave-regular , in accord with our expectations. This is in accordance with the fact that, in these cases the singularity is null (as it occurs at $`z=\mathrm{}`$) and the evolution of initial date is the ordinary Cauchy which is unique. However, for $`c>0`$, the metric near $`z=z_0`$ is approximated by $$ds^2=\stackrel{~}{x}^{1/3}(dx_{\|\|}^2+d\stackrel{~}{x}^2),$$ (46) where $`\stackrel{~}{x}=z_0z0^+`$. Now both solutions $`\psi 1,\stackrel{~}{x}^{1/2}`$ are normalizable and the singularity is wave-singular. This is consistent with field theory expectations and a different criterion presented in . Actually, there is a subtlety in the choice of the norm of the wavefunction. As refs. we use the Sobolev norm which contains derivatives of the wavefunction, whereas in a norm without derivatives is used. This distinction is important for example in the case of the negative mass Schwarzschild black hole: using the Sobolev norm it is wave-regular, but it is singular in the convention of . We have also checked the criterion using the Sobolev norm in all backgrounds corresponding to continuous D3-, M2- and M5-brane distributions describing states in Coulomb branch of the corresponding supersymmetric field theories. Referring to the conventions in , the naked curvature singularities occuring close to the brane distributions are physical (unphysical) for $`n2`$ ($`n=1`$). This is in accordance with field theory expectations since in the unphysical cases the brane distributions contain a component corresponding to negative tension branes . The criterion using the Sobolev norm is in accord with the criterion of for all cases we are aware of, but it seems to us that it is more natural to impose. There seems to be a consistent picture emerging using the simple criterion of wave-regularity to distinguish acceptable singularities from bad ones. Evidently a better understanding of string theory in such backgrounds is desireable, and in hindsight it is interesting to note that this is one of the few examples in the AdS/CFT correspondence where the field theory side can teach us something about the associated string theory. ## Acknowledgements A.B. and K.S. would like to thank the Erwin Schrödinger Institute in Vienna and the CIT-USC center for Theoretical Physics at USC in Los Angeles, respectively, for hospitality and generous financial support during the final stages of this research, as well as for the opportunity to present these results in seminars prior to submitting the paper at hep-th. Also A.B. thanks S.J. Rey for extensive discussions.
warning/0004/cond-mat0004226.html	ar5iv	text	# Untitled Document Renormalization Group, hidden symmetries and approximate Ward identities in the XYZ model, I. G. Benfatto, V. Mastropietro Supported by MURST, Italy, and EC HCM contract number CHRX-CT94-0460. e-mail: benfatto@mat.uniroma2.it, mastropi@mat.uniroma2.it. Dipartimento di Matematica, Università di Roma “Tor Vergata” Via della Ricerca Scientifica, I-00133, Roma Abstract. Using renormalization group methods, we study the Heisenberg-Ising $`XYZ`$ chain in an external magnetic field directed as the $`z`$ axis, in the case of small coupling $`J_3`$ in the $`z`$ direction. We study the asymptotic behaviour of the spin space-time correlation function in the direction of the magnetic field and the singularities of its Fourier transform. The work is organized in two parts. In the present paper an expansion for the ground state energy and the effective potential is derived, which is convergent if the running coupling constants are small enough. In the subsequent paper, by using hidden symmetries of the model, we show that this condition is indeed verified, if $`J_3`$ is small enough, and we derive an expansion for the spin correlation function. We also prove, by means of an approximate Ward identity, that a critical index, related with the asymptotic behaviour of the correlation function, is exactly vanishing. 1. Introduction 1.1 If $`(S_x^1,S_x^2,S_x^3)=\frac{1}{2}(\sigma _x^1,\sigma _x^2,\sigma _x^3)`$, for $`i=1,2,\mathrm{},L`$, $`\sigma _i^\alpha `$, $`\alpha =1,2,3`$, being the Pauli matrices, the Hamiltonian of the Heisenberg-Ising $`XYZ`$ chain is given by $$H=\underset{x=1}{\overset{L1}{}}[J_1S_x^1S_{x+1}^1+J_2S_x^2S_{x+1}^2+J_3S_x^3S_{x+1}^3+hS_x^3]hS_L^3+U_L^1,$$ $`(1.1)`$ where the last term, to be fixed later, depends on the boundary conditions. The space-time spin correlation function at temperature $`\beta ^1`$ is given by $$\mathrm{\Omega }_{L,\beta }^\alpha (𝐱)=<S_𝐱^\alpha S_\mathrm{𝟎}^\alpha >_{L,\beta }<S_𝐱^\alpha >_{L,\beta }<S_\mathrm{𝟎}^\alpha >_{L,\beta },$$ $`(1.2)`$ where $`𝐱=(x,x_0)`$, $`S_𝐱^\alpha =e^{Hx_0}S_x^\alpha e^{Hx_0}`$ and $`<.>_{L,\beta }=Tr[e^{\beta H}.]/Tr[e^{\beta H}]`$ denotes the expectation in the grand canonical ensemble. We shall use also the notation $`\mathrm{\Omega }^\alpha (𝐱)lim_{L,\beta \mathrm{}}\mathrm{\Omega }_{L,\beta }^\alpha (𝐱)`$. The Hamiltonian (1.1) can be written \[LSM\] as a fermionic interacting spinless Hamiltonian. In fact, it is easy to check that the operators $$a_x^\pm \left[\underset{y=1}{\overset{x1}{}}(\sigma _y^3)\right]\sigma _x^\pm $$ $`(1.3)`$ are a set of anticommuting operators and that, if $`\sigma _x^\pm =(\sigma _x^1\pm i\sigma _x^2)/2`$, we can write $$\sigma _x^{}=e^{i\pi _{y=1}^{x1}a_y^+a_y^{}}a_x^{},\sigma _x^+=a_x^+e^{i\pi _{y=1}^{x1}a_y^+a_y^{}},\sigma _x^3=2a_x^+a_x^{}1.$$ $`(1.4)`$ Hence, if we fix the units so that $`J_1+J_2=2`$ and we introduce the anisotropy $`u=(J_1J_2)/(J_1+J_2)`$, we get $$\begin{array}{cc}\hfill H& =\underset{x=1}{\overset{L1}{}}\{\frac{1}{2}[a_x^+a_{x+1}^{}+a_{x+1}^+a_x^{}]\frac{u}{2}[a_x^+a_{x+1}^++a_{x+1}^{}a_x^{}]\hfill \\ & J_3(a_x^+a_x^{}\frac{1}{2})(a_{x+1}^+a_{x+1}^{}\frac{1}{2})\}h_{x=1}^L(a_x^+a_x^{}\frac{1}{2})+U_L^2,\hfill \end{array}$$ $`(1.5)`$ where $`U_L^2`$ is the boundary term in the new variables. We choose it so that the fermionic Hamiltonian (1.5) coincides with the Hamiltonian of a fermion system on the lattice with periodic boundary conditions, that is we put $`U_L^2`$ equal to the term in the first sum in the r.h.s. of (1.5) with $`x=L`$ and $`a_{L+1}^\pm =a_1^\pm `$ (in \[LMS\] this choice for the $`XY`$ chain is called “c-cyclic”). It is easy to see that this choice corresponds to fix the boundary conditions for the spin variables so that $$U_L^1=\frac{1}{2}[\sigma _L^+e^{i\pi 𝒩}\sigma _1^{}+\sigma _L^{}e^{i\pi 𝒩}\sigma _1^+]\frac{u}{2}[\sigma _L^+e^{i\pi 𝒩}\sigma _1^++\sigma _L^{}e^{i\pi 𝒩}\sigma _1^{}]\frac{J_3}{4}\sigma _L^3\sigma _1^3,$$ $`(1.6)`$ where $`𝒩=_{x=1}^La_x^+a_x`$. Strictly speaking, with this choice $`U_L^1`$ does not look really like a boundary term, because $`𝒩`$ depends on all the spins of the chain. However $`[(1)^𝒩,H]=0`$; hence the Hilbert space splits up in two subspaces on which $`(1)^𝒩`$ is equal to $`1`$ or to $`1`$ and on each of these subspaces $`U_L^1`$ really depends only on the boundary spins. One expects that, in the $`L\mathrm{}`$ limit, the correlation functions are independent on the boundary term, but we shall not face here this problem. 1.2 The Heisenberg $`XYZ`$ chain has been the subject of a very active research over many years with a variety of methods. A first class of results is based on the exact solutions. If one of the three parameters is vanishing (e.g. $`J_3=0`$), the model is called $`XY`$ chain. Its solution is based on the fact that the hamiltonian, in the fermionic form (1.5), is quadratic in the fermionic fields, so that it can be diagonalized (see \[LSM\], \[LSM1\]) by a Bogoliubov transformation. If $`u=0`$, we get the free Fermi gas with Fermi momentum $`p_F=\mathrm{arccos}(h)`$; if $`\|u\|>0`$, it turns out that the energy spectrum has a gap at $`p_F`$. The equal time correlation functions $`\mathrm{\Omega }^\alpha (x,0)`$ were explicitly calculated in \[Mc\] (even at finite $`L`$ and $`\beta `$), in the case $`h=0`$, that is $`p_F=\pi /2`$. Note that, while $`\mathrm{\Omega }^3(𝐱)`$ coincides with the correlation function of the density in the fermionic representation of the model, $`\mathrm{\Omega }^1(𝐱)`$ and $`\mathrm{\Omega }^2(𝐱)`$ are given by quite complicated expressions. It turns out, for example, that, if $`\|u\|<1`$, $`\mathrm{\Omega }^3(x,0)`$ is of the following form: $$\mathrm{\Omega }^3(x,0)=\frac{\alpha ^{\|x\|}}{\pi ^2x^2}\mathrm{sin}^2\left(\frac{\pi x}{2}\right)F(\|x\|\mathrm{log}\alpha ,\|x\|),\alpha =(1\|u\|)/(1+\|u\|),$$ $`(1.7)`$ where $`F(\gamma ,n)`$ is a bounded function, such that, if $`\gamma 1`$, $`F(\gamma ,n)=1+O(\gamma \mathrm{log}\gamma )+O(1/n)`$, while, if $`\gamma 1`$ and $`n2\gamma `$, $`F(\gamma ,n)=\pi /2+O(1/\gamma )`$. For $`\|h\|>0`$, it is not possible to get a so explicit expression for $`\mathrm{\Omega }^3(x,0)`$. However, it is not difficult to prove that, if $`\|u\|<\mathrm{sin}p_F`$, $`\|\mathrm{\Omega }^3(x,0)\|\alpha ^{\|x\|}`$ and, if $`x0`$ and $`\|ux\|1`$ $$\mathrm{\Omega }^3(x,0)=\frac{1}{\pi ^2x^2}\mathrm{sin}^2(p_Fx)[1+O(\|ux\|\mathrm{log}\|ux\|)+O(1/\|x\|)].$$ $`(1.8)`$ Note that, if $`u=0`$, a very easy calculation shows that $`\mathrm{\Omega }^3(x,0)=(\pi ^2x^2)^2\mathrm{sin}^2(p_Fx)`$. We want to stress that the only case in which the correlation functions and their asymptotic behaviour can be computed explicitly in a rigorous way is just the $`J_3=0`$ case. If two parameters are equal (e.g. $`J_1=J_2`$), but $`J_30`$, the model is called $`XXZ`$ model. In the case $`h=0`$, it was solved in \[YY\] via the Bethe-ansatz, in the sense that the Hamiltonian was diagonalized. However, it was not possible till now to obtain the correlation functions from the exact solution. Such solution is a particular case of the general solution of the XYZ model by Baxter \[B\], but again only in the case of zero magnetic field. The ground state energy has been computed and it has been proved that there is a gap in the spectrum, which, if $`J_1J_2`$ and $`J_3`$ are not too large, is given approximately by (see \[LP\]) $$\mathrm{\Delta }=8\pi \frac{\mathrm{sin}\mu }{\mu }\|J_1\|\left(\frac{\|J_1^2J_2^2\|}{16(J_1^2J_3^2)}\right)^{\frac{\pi }{2\mu }}$$ $`(1.9)`$ with $`\mathrm{cos}\mu =J_3/J_1`$. The solution is based on the fact that the $`XYZ`$ chain with periodic boundary conditions is equivalent to the eight vertex model, in the sense that $`H`$ is proportional to the logarithmic derivative with respect to a parameter of the eight vertex transfer matrix, if a suitable identification of the parameters is done, see \[S\], \[B\]. The eight vertex model is obtained by putting arrows in a suitable way on a two-dimensional lattice with $`M`$ rows, $`L`$ columns and periodic boundary conditions. There are eight allowed vertices, and with each of them an energy is associated in a suitable way (there are four different values of the energy). With the above choice of the parameters and $`TT_c<0`$ and small, $`u=O(\|TT_c\|)`$, so that the critical temperature of the eight vertex model corresponds to no anisotropy in the $`XYZ`$ chain. Moreover, see \[JKM\], the correlation function $`C_x`$ between two vertical arrows in a row, separated by $`x`$ vertices, is given, in the limit $`M\mathrm{}`$, by $`C_x=<S_0^2S_x^2>`$. However, an explicit expression for the correlation functions cannot be derived for the $`XYZ`$ or the eight vertex model. In \[JKM\] the correlation length of $`C_x`$ was computed heuristically under some physical assumptions (an exact computation is difficult because it does not depend only on the largest and the next to the largest eigenvalues). The result is $`\xi ^1=(TT_c)^{\frac{\pi }{2\mu }}`$, if $`\xi `$ is the correlation length. One sees that the critical index of the correlation length is non universal. Another interesting observation is that the $`XYZ`$ model is equivalent to two interpenetrating two-dimensional Ising lattices with nearest-neighbor coupling, interacting via a four spins coupling (which is proportional to $`J_3`$). The four spin correlation function is identical to $`C_x`$. In the decoupling limit $`J_3=0`$ the two Ising lattices are independent and one can see that the Ising model solution can be reduced to the diagonalization, via a Bogoliubov transformation, of a quadratic Fermi Hamiltonian, see \[LSM1\]. Recent new results using the properties of the transfer matrix can be found in \[EFIK\], in which an integro-difference equation for the correlation function of the XXZ chain is obtained. It is however not clear how to deduce the physical properties of the correlation function from this equation. 1.3 Since it is very difficult to extract detailed information on the behaviour of the correlation functions from the above exact solutions, the XYZ model has been studied by quantum field theory methods, see \[LP\]. The idea is to approximate the fermionic hamiltonian (1.5) by the hamiltonian of the massive Thirring model, describing a massive relativistic spinning particle on the continuum $`d=1`$ space interacting with a local current-current potential (for a heuristic justification of this approximation, see \[A\]). As a relativistic field theory, the massive Thirring model is plagued by ultraviolet divergences, which were absent in the original model, defined on a lattice; one can heuristically remove this problem by introducing ”by hand” an ultraviolet cut-off. A way to introduce it could be to consider a short-ranged instead of a local potential; if $`J_1=J_2`$, this means that we have approximated the $`XXZ`$-chain with the Luttinger model, whose correlation functions can be explicitly computed, see \[ML\], \[BGM\]. The Luttinger model is defined in terms of two fields $`\psi _{𝐱,\omega }`$, $`\omega =\pm 1`$, and one expects that, if $`\|h\|<1`$ and $`J_3`$ is small enough, the large distance asymptotic behaviour of $`\mathrm{\Omega }^3(𝐱)`$ is qualitatively similar to that of the truncated correlation of the operator $`\rho _𝐱=\psi _𝐱^+\psi _𝐱^{}`$, where $`\psi _𝐱^\sigma =_\omega \mathrm{exp}(i\sigma \omega p_Fx)\psi _{𝐱,\omega }`$, if some “reasonable” relationship between the parameters of the two models is assumed. One can make for instance the substitutions $`\lambda J_3`$ and $`p_0^1a=1`$, if $`\lambda `$ is the coupling in the Luttinger model, $`a`$ is the chain step and $`p_0^1`$ is the potential range. Moreover, one expects that it is possible to choose a constant $`\nu `$ of order $`J_3`$, so that $`h=h_0+\nu `$ and $`p_F=\mathrm{arccos}(J_3h_0)`$, see §$`\mathrm{}`$1.4 below. Of course such identification is completely arbitrary, but one can hope that for large distances the function $`\mathrm{\Omega }^3(𝐱)`$ has something to do with the truncated correlation of $`\rho _𝐱`$, which can be obtained by the general formula (2.5) of \[BGM\], based on the exact solution of \[ML\]. There is apparently a problem, since the expectation of $`\rho _𝐱`$ is infinite; however, it is possible to see that there exists the limit, as $`\epsilon _1,\epsilon _20^+`$, of $`[<\rho _{𝐱,\epsilon _1}\rho _{𝐲,\epsilon _2}><\rho _{𝐱,\epsilon _1}><\rho _{𝐲,\epsilon _2}>]`$, where $`\rho _{𝐱,\epsilon }=\psi _{(x,x_0+\epsilon )}^+\psi _{(x,x_0)}^{}`$, and it is natural to take this quantity, let us call it $`G(𝐱𝐲)`$, as the truncated correlation of $`\rho _𝐱`$. Let us define $`v_0=\mathrm{sin}p_F`$; from (2.5) of \[BGM\] (by inserting a missing $`(\epsilon _i\epsilon _j)`$ in the last sum), it follows that, for $`\|𝐱\|\mathrm{}`$ $$G(𝐱)[1+\lambda a_1(\lambda )]\frac{\mathrm{cos}(2p_Fx)}{2\pi ^2[(v_0^{}x_0)^2+x^2]^{1+\lambda a_3(\lambda )}}+\frac{(v_0x_0)^2x^2}{2\pi ^2[(v_0x_0)^2+x^2]^2},$$ $`(1.10)`$ where $`v_0^{}=v_0[1+\lambda a_2(\lambda )]`$ and $`a_i(\lambda )`$, $`i=1,2,3`$, are bounded functions. Note that, in the second term in the r.h.s. of (1.10), the bare Fermi velocity $`v_0`$ appears, instead of the renormalized one, $`v_0^{}`$, as one could maybe expect. In the physical literature, it is more usual the introduction of other ultraviolet cutoffs, such that the resulting model is not exactly soluble, even if $`J_1=J_2`$; however, it can be studied heuristically, see \[LP\], and the resulting density-density correlation function is more or less of the form (1.10). If $`J_1J_2`$, there is no soluble model suitable for a similar analysis of the large distance behaviour of $`\mathrm{\Omega }^3(𝐱)`$. However, one can guess that the asymptotic behaviour is still of the form (1.10), if $`1<<\|𝐱\|<<1/\|u\|^\alpha `$, for some $`\alpha `$. We shall prove that this is indeed true, with $`\alpha =1+O(J_3)`$. 1.4 In this paper we develop a rigorous renormalization group analysis for the $`XYZ`$ Hamiltonian in its fermionic form (some “not optimal” bounds for the correlation function $`\mathrm{\Omega }^3(𝐱)`$ were already found in \[M2\]). As we said before, $`\mathrm{\Omega }^3(𝐱)`$ can be obtained from the exact solution only in the case $`J_3=0`$, when the fermionic theory is a non interacting one. In particular, if $`𝐱=(x,0)`$ and $`\|ux\|<<1`$, (1.8) and a more detailed analysis of the “small” terms in the r.h.s. (in order to prove that their derivatives of order $`n`$ decay as $`\|x\|^n`$), show that $`\mathrm{\Omega }^3(x,0)`$ is a sum of “oscillating” functions with frequencies $`(np_F)/\pi \mathrm{mod}\mathrm{\hspace{0.17em}1}`$, $`n=0,\pm 1`$, where $`p_F=\mathrm{arccos}(h)`$; this means that its Fourier transform has to be a smooth function, even for $`u=0`$, in the neighborhood of any momentum $`k0,\pm 2p_F`$. These frequencies are proportional to $`p_F`$, so they depend only on the external magnetic field $`h`$. If $`J_30`$, a similar property is satisfied for the leading terms in the asymptotic behaviour, as we shall prove, but the value of $`p_F`$ depends in general also on $`u`$ and $`J_3`$. For example, if $`u=0`$, the Hamiltonian (1.5) is equal, up to a constant, to the Hamiltonian of a free fermion gas with Fermi momentum $`p_F=\mathrm{arccos}(J_3h)`$ plus an interaction term proportional to $`J_3`$. As it is well known, the interaction modifies the Fermi momentum of the system by terms of order $`J_3`$ and it is convenient (see \[BG\], for example), in order to study the interacting model, to fix the Fermi momentum to an interaction independent value, by adding a counterterm to the hamiltonian. We proceed here in a similar way, that is we fix $`p_F`$ and $`h_0`$ so that $$h=h_0\nu ,\mathrm{cos}p_F=J_3h_0,$$ $`(1.11)`$ and we look for a value of $`\nu `$, depending on $`u,J_3,h_0`$, such that, as in the $`J_3=0`$ case, the leading terms in the asymptotic behaviour of $`\mathrm{\Omega }_{L,\beta }^3(𝐱)`$ can be represented as a sum of oscillating functions with frequencies $`(np_F)/\pi \mathrm{mod}\mathrm{\hspace{0.17em}1}`$, $`n=0,\pm 1`$. As we shall see, we can realize this program only if $`J_3`$ is small enough and it turns out that $`\nu `$ is of order $`J_3`$. It follows that we can only consider magnetic fields such that $`\|h\|<1`$. Moreover, it is clear that the equation $`h=h_0\nu (u,J_3,h_0)`$ can be inverted, once the function $`\nu (u,J_3,h_0)`$ has been determined, so that $`p_F`$ is indeed a function of the parameters appearing in the original model. If $`J_1=J_2`$, it is conjectured, on the base of heuristic calculations, that to fix $`p_F`$ is equivalent to the impose the condition that, in the limit $`L,\beta \mathrm{}`$, the density is fixed (“Luttinger Theorem”) to the free model value $`\rho =p_F/\pi `$. Remembering that $`\rho \frac{1}{2}`$ is the magnetization in the $`3`$-direction for the original spin variables, this would mean that to fix $`p_F`$ is equivalent to fix the magnetization in the $`3`$ direction, by suitably choosing the magnetic field. If $`J_1J_2`$, there is in any case no simple relation between $`p_F`$ and the mean magnetization, as one can see directly in the case $`J_3=0`$, where one can do explicit calculations. The only exception is the case $`p_F=\pi /2`$, where one can see that, in the limit $`L\mathrm{}`$, $`\nu =J_3`$ (so that $`h=0`$ by (1.11)) and that $`<S_x^3>=0`$. This last property easily follows from the observation that, if one choose $`h=0`$ in the original Hamiltonian (1.1), then the expectation of $`S_x^3`$ has to be equal to zero, by symmetry reasons, up to terms which go to $`0`$ for $`L\mathrm{}`$. Our main achievement is an expansion of $`\mathrm{\Omega }_{L,\beta }^3(𝐱)`$, to be derived in paper II, which provides a very detailed and explicit description of it. We state in the following theorem some of its properties, but we stress that many other interesting properties of $`\mathrm{\Omega }_{L,\beta }^3(𝐱)`$ can be extracted from the expansion. 1.5 Theorem. Suppose that the equations (1.11) are satisfied and that $`v_0=\mathrm{sin}p_F\overline{v}_0>0`$, for some value of $`\overline{v}_0`$ fixed once for all, and let us define $`a_0=\mathrm{min}\{p_F/2,(\pi p_F)/2\}`$; then the following is true. a) There exists a constant $`\epsilon `$, such that, if $`(u,J_3)𝒜`$, with $$𝒜=\{(u,J_3):\|u\|\frac{a_0}{8(1+\sqrt{2})},\|J_3\|\epsilon \},$$ $`(1.12)`$ it is possible to choose $`\nu `$, so that $`\|\nu \|c\|J_3\|`$, for some constant $`c`$ independent of $`L`$, $`\beta `$, $`u`$, $`J_3`$, $`p_F`$, and the spin correlation function $`\mathrm{\Omega }_{L,\beta }^3(𝐱)`$ is a bounded (uniformly in $`L`$, $`\beta `$, $`p_F`$ and $`(u,J_3)𝒜`$) function of $`𝐱=(x,x_0)`$, $`x=1,\mathrm{},L`$, $`x_0[0,\beta ]`$, periodic in $`x`$ and $`x_0`$ of period $`L`$ and $`\beta `$ respectively, continuous as a function of $`x_0`$. b) We can write $$\mathrm{\Omega }_{L,\beta }^3(𝐱)=\mathrm{cos}(2p_Fx)\mathrm{\Omega }_{L,\beta }^{3,a}(𝐱)+\mathrm{\Omega }_{L,\beta }^{3,b}(𝐱)+\mathrm{\Omega }_{L,\beta }^{3,c}(𝐱),$$ $`(1.13)`$ with $`\mathrm{\Omega }_{L,\beta }^{3,i}(𝐱)`$, $`i=a,b,c`$, continuous bounded functions, which are infinitely times differentiable as functions of $`x_0`$, if $`i=a,b`$. Moreover, there exist two constants $`\eta _1`$ and $`\eta _2`$ of the form $$\eta _1=a_1J_3+O(J_3^2),\eta _2=a_2J_3+O(J_3^2),$$ $`(1.14)`$ $`a_1`$ and $`a_2`$ being positive constants, uniformly bounded in $`L`$, $`\beta `$, $`p_F`$ and $`(u,J_3)𝒜`$, such that the following is true. Let us define $$𝐝(𝐱)=(\frac{L}{\pi }\mathrm{sin}(\frac{\pi x}{L}),\frac{\beta }{\pi }\mathrm{sin}(\frac{\pi x_0}{\beta }))$$ $`(1.15)`$ and suppose that $`\|𝐝(𝐱)\|1`$. Then, given any positive integers $`n`$ and $`N`$, there exist positive constants $`\vartheta <1`$ and $`C_{n,N}`$, independent of $`L`$, $`\beta `$, $`p_F`$ and $`(u,J_3)𝒜`$, so that, for any integers $`n_0,n_10`$ and putting $`n=n_0+n_1`$, $$\|_{x_0}^{n_0}\overline{}_x^{n_1}\mathrm{\Omega }_{L,\beta }^{3,a}(𝐱)\|\frac{1}{\|𝐝(𝐱)\|^{2+2\eta _1+n}}\frac{C_{n,N}}{1+[\mathrm{\Delta }\|𝐝(𝐱)\|]^N},$$ $`(1.16)`$ $$\|_{x_0}^{n_0}\overline{}_x^{n_1}\mathrm{\Omega }_{L,\beta }^{3,b}(𝐱)\|\frac{1}{\|𝐝(𝐱)\|^{2+n}}\frac{C_{n,N}}{1+[\mathrm{\Delta }\|𝐝(𝐱)\|]^N},$$ $`(1.17)`$ $$\|\mathrm{\Omega }_{L,\beta }^{3,c}(𝐱)\|\frac{1}{\|𝐝(𝐱)\|^2}\left[\frac{1}{\|𝐝(𝐱)\|^\vartheta }+\frac{(\mathrm{\Delta }\|𝐝(𝐱)\|)^\vartheta }{\|𝐝(𝐱)\|^{\mathrm{min}\{0,2\eta _1\}}}\right]\frac{C_{0,N}}{1+[\mathrm{\Delta }\|𝐝(𝐱)\|]^N},$$ $`(1.18)`$ where $`\overline{}_x`$ denotes the discrete derivative and $$\mathrm{\Delta }=\mathrm{max}\{\|u\|^{1+\eta _2},\sqrt{(v_0\beta )^2+L^2}\}.$$ $`(1.19)`$ c) There exist the limits $`\mathrm{\Omega }^{3,i}(𝐱)=lim_{L,\beta \mathrm{}}\mathrm{\Omega }_{L,\beta }^{3,i}(𝐱)`$, $`𝐱\text{}\times \text{}`$; they satisfy the bounds (1.16), with $`\|𝐱\|`$ in place of $`\|𝐝(𝐱)\|`$. Moreover, $`\mathrm{\Omega }^{3,a}(𝐱)`$ and $`\mathrm{\Omega }^{3,b}(𝐱)`$ are even functions of $`𝐱`$ and there exists a constant $`\delta ^{}`$, of order $`J_3`$, such that, if $`1\|𝐱\|\mathrm{\Delta }^1`$ and $`v_0^{}=v_0(1+\delta ^{})`$, given any $`N>0`$ $$\begin{array}{cc}\hfill \mathrm{\Omega }^{3,a}(𝐱)& =\frac{1+A_1(𝐱)}{2\pi ^2[x^2+(v_0^{}x_0)^2]^{1+\eta _1}},\hfill \\ \hfill \mathrm{\Omega }^{3,b}(𝐱)& =\frac{1}{2\pi ^2[x^2+(v_0^{}x_0)^2]}\left\{\frac{x_0^2(x/v_0^{})^2}{x^2+(v_0^{}x_0)^2}+A_2(𝐱)\right\},\hfill \end{array}$$ $`(1.20)`$ $$\|A_i(𝐱)\|C_N\left\{\frac{1}{1+\|𝐱\|^N}+\|J_3\|+(\mathrm{\Delta }\|𝐱\|)^{1/2}\right\},$$ $`(1.21)`$ for some constant $`C_N`$. The function $`\mathrm{\Omega }^{3,a}(𝐱)`$ is the restriction to $`\text{}\times \text{}`$ of a function on $`\text{}^2`$, satisfying the symmetry relation $$\mathrm{\Omega }^{3,a}(x,x_0)=\mathrm{\Omega }^{3,a}(x_0v_0^{},\frac{x}{v_0^{}}).$$ $`(1.22)`$ d) Let $`\widehat{\mathrm{\Omega }}^3(𝐤)`$, $`𝐤=(k,k_0)[\pi ,\pi ]\times \text{}^1`$, the Fourier transform of $`\mathrm{\Omega }^3(𝐱)`$. For any fixed $`𝐤`$ with $`𝐤(0,0),(\pm 2p_F,0)`$, $`\widehat{\mathrm{\Omega }}^3(𝐤)`$ is uniformly bounded as $`u0`$; moreover, for some constant $`c_2`$, $$\begin{array}{cc}\hfill \|\widehat{\mathrm{\Omega }}^3(0,0)\|& c_2\left[1+\|J_3\|\mathrm{log}\frac{1}{\mathrm{\Delta }}\right],\hfill \\ \hfill \|\widehat{\mathrm{\Omega }}^3(\pm 2p_F,0)\|& c_2\frac{1\mathrm{\Delta }^{2\eta _1}}{2\eta _1}.\hfill \end{array}$$ $`(1.23)`$ Finally, if $`u=0`$, $`\|\widehat{\mathrm{\Omega }}^3(𝐤)\|c_2[1+\|J_3\|\mathrm{log}\|𝐤\|^1]`$ near $`𝐤=(0,0)`$, and, at $`𝐤=(\pm 2p_F,0)`$, it is singular only if $`J_3<0`$; in this case it diverges as $`\|𝐤(\pm 2p_F,0)\|^{2\eta _1}/\|\eta _1\|`$. e) Let $`G(x)=\mathrm{\Omega }^3(x,0)`$ and $`\widehat{G}(k)`$ its Fourier transform. For any fixed $`k0,\pm 2p_F`$, $`\widehat{G}(k)`$ is uniformly bounded as $`u0`$, together with its first derivative; moreover $$\begin{array}{cc}\hfill \|_k\widehat{G}(0)\|& c_2,\hfill \\ \hfill \|_k\widehat{G}(\pm 2p_F)\|& c_2(1+\mathrm{\Delta }^{2\eta _1}).\hfill \end{array}$$ $`(1.24)`$ Finally, if $`u=0`$, $`_k\widehat{G}(k)`$ has a first order discontinuity at $`k=0`$, with a jump equal to $`1+O(J_3)`$, and, at $`k=\pm 2p_F`$, it is singular only if $`J_3<0`$; in this case it diverges as $`\|k(\pm 2p_F)\|^{2\eta _1}`$. 1.6 Remarks. a)The above theorem holds for any magnetic field $`h`$ such that $`\mathrm{sin}p_F>0`$; remember that the exact solution given in \[B\] is valid only for $`h=0`$. Moreover $`u`$ has not to be very small, but we only need a bound of order $`1`$ on its value, see (1.12); the only perturbative parameter is $`J_3`$. However the interesting (and more difficult) case is when also $`u`$ is small. b)A naive estimate of $`\epsilon `$ is $`\epsilon =c(\mathrm{sin}p_F)^\alpha `$, with $`c,\alpha `$ positive numbers; in other words we must take smaller and smaller $`J_3`$ for $`p_F`$ closer and closer to $`0`$ or $`\pi `$, i.e. for magnetic fields of size close to $`1`$. It is unclear at the moment if this is only a technical problem or a property of the model. c)If $`J_1J_2`$ and $`J_30`$, one can distinguish, like in the $`J_3=0`$ case (1.7), two different regimes in the asymptotic behaviour of the correlation function $`\mathrm{\Omega }^3(𝐱)`$, discriminated by an intrinsic length $`\xi `$, which is approximately given by the inverse of spectral gap, whose size, is of order $`\|u\|^{1+\eta _2}`$, see (1.19), in agreement with (1.9), found by the exact solution. If $`1<<\|𝐱\|<<\xi `$, the bounds for the correlation function are the same as in the gapless $`J_1=J_2`$ case; if $`\xi <<\|𝐱\|`$, there is a faster than any power decay with rate of order $`\xi ^1`$. In the first region we can obtain the exact large distance asymptotic behaviour of $`\mathrm{\Omega }^3(𝐱)`$, see (1.20),(1.21); in the second region only an upper bound is obtained. Note that, even in the $`J_3=0`$ case, it is not so easy to obtain a more precise result, if $`h0`$, see §1.2. The spin interaction in the $`z`$ direction has the effect that the gap becomes anomalous, in the sense that it acquires a critical index $`\eta _2`$; the ratio between the “renormalized” and the “bare” gap is very small or very large, if $`u`$ is small, depending on the sign of $`J_3`$. d)It is useful to compare the expression for the large distance behaviour of $`\mathrm{\Omega }^3(𝐱)`$ in the case $`u=0`$ with its analogous for the Luttinger model, see §1.3. A first difference is that, while in the Luttinger model the Fermi momentum is independent of the interaction, in the $`XYZ`$ model in general it is changed non trivially by the interaction, unless the magnetic external field is zero, i.e. $`p_F=\frac{\pi }{2}`$. The reason is that the Luttinger model has special parity properties which are not satisfied by the $`XYZ`$ chain (except if the magnetic field is vanishing). e)Another peculiar property of the Luttinger model correlation function is that it depends on $`p_F`$ only through the factor $`\mathrm{cos}(2p_Fx)`$; this is true not only for the asymptotic behaviour (1.10), but also for the complete expression given in \[BGM\], and is due to a special symmetry of the Luttinger model (the Fermi momentum disappears from the Hamiltonian if a suitable redefinition of the fermionic fields is done, see \[BGM\]). This property is of course not true in the $`XYZ`$ model and in fact the dependence on $`p_F`$ of $`\mathrm{\Omega }^3(𝐱)`$ is very complicated. However we prove that $`\mathrm{\Omega }^3(𝐱)`$ can be written as sum of three terms, see (1.13), and the first two terms are very similar to the two terms in the r.h.s. of (1.10). In particular, the functions $`\mathrm{\Omega }^{3,a}(𝐱)`$ and $`\mathrm{\Omega }^{3,b}(𝐱)`$ have the same power decay as the analogous functions in the Luttinger model and are “free of oscillations”, in the sense that each derivative increases the decay power of one unit, see (1.16),(1.17). This is not true for the third term $`\mathrm{\Omega }^{3,c}(𝐱)`$, which does not satisfy a similar bound, because of the presence of oscillating contributions. However we can prove that such term, if $`u=0`$, is negligible for large distances, see (1.18) (note that $`\vartheta `$ is $`J_3`$ and $`u`$ independent, unlike $`\eta _1`$). Of course this is true only for small $`J_3`$ and it could be that $`\mathrm{\Omega }^{3,c}(𝐱)`$ plays an important role for larger $`J_3`$. If we compare, in the case $`u=0`$, the functions $`\mathrm{\Omega }^{3,a}(𝐱)`$ and $`\mathrm{\Omega }^{3,b}(𝐱)`$, see (1.20), with the corresponding ones in the Luttinger model, see (1.10), we see that they differ essentially for the non oscillating functions $`A_i(𝐱)`$, containing terms of higher order in our expansion. However, this difference is not important in the case of $`\mathrm{\Omega }^{3,a}(𝐱)`$, which also satisfies the same symmetry property (1.22) as the analogue in the Luttinger model, of course with different values of $`v_0^{}`$; note that the validity of (1.22) allows to interpret $`v_0^{}`$ as the renormalized Fermi velocity. Guided by the analogy with the Luttinger model, one would like to prove a similar property for $`\mathrm{\Omega }^{3a}(𝐱)`$ with $`v_0`$ replacing $`v_0^{}`$; such property holds in fact for the Luttinger model, see (1.10). However we were not able to prove a similar properties for $`A_2(𝐱)`$, and this has some influence on our results, see below. f)Another important property of the Luttinger model correlation function is the fact that the “not oscillating term”, that is the term corresponding to $`\mathrm{\Omega }^{3,b}(𝐱)`$, does not acquire a critical index, contrary to what happens for the term corresponding to $`\mathrm{cos}(2p_Fx)\mathrm{\Omega }^{3,b}(𝐱)`$. Hence one is naturally led to the conjecture that the critical index of $`\mathrm{\Omega }_{L,\beta }^{3,b}(𝐱)`$ is still vanishing, see for instance \[Sp\]. In our expansion, the critical index of $`\mathrm{\Omega }^{3,b}(𝐱)`$ is represented as a convergent series, but, even if an explicit computation of the first order term gives a vanishing result, it is not obvious that this is true at any order. However, due to some hidden symmetries of the model (i.e. symmetries approximately enjoyed by the relevant part of the effective interaction), we can prove a suitable approximate Ward identity, implying that all the coefficients of the series are indeed vanishing. g) The above properties can be used to study the Fourier transform $`\widehat{G}(k)`$ of the equal time correlation function $`G(x)=\mathrm{\Omega }^3(x,0)`$. If $`J_3=0`$, $`\widehat{G}(k)`$ is bounded together with its first order derivative up to $`u=0`$; in fact, the possible logarithmic divergence at $`k=\pm 2p_F`$ and $`k=0`$ (if $`u=0`$) of $`\widehat{G}(k)`$ is changed by the parity properties of $`G(x)`$ in a first order discontinuity. If $`J_30`$, $`\widehat{G}(k)`$ behaves near $`k=\pm 2p_F`$ in a completely different way. In fact it is bounded and continuous if $`J_3>0`$, while it has a power like singularity, if $`u=0`$ and $`J_3<0`$, see item e) of Theorem (1.5). This is due to the fact that the critical index $`\eta _1`$, characterizing the asymptotic behaviour of $`\mathrm{\Omega }^{3,a}(𝐱)`$, has the same sign of $`J_3`$ (note that $`\eta _1`$ has nothing to do with the critical index $`\eta `$ related with the two point fermionic Schwinger function, which is $`O((J_3)^2)`$). On the other hand, the behaviour of $`\widehat{G}(k)`$ near $`k=0`$ is the same for the Luttinger model, the $`XYZ`$ model and the free fermionic gas ($`J_1=J_2`$, $`J_3=0`$) (see also \[Sp\] for a heuristic explanation). This is due to the vanishing of the critical index related with $`\mathrm{\Omega }^{3,b}(𝐱)`$ and to the parity properties of the leading terms, which change, as in the $`J_3=0`$ case, the apparent dimensional logarithmic divergence in a first order discontinuity, h) If $`u=0`$, the (two dimensional) Fourier transform can be singular only at $`𝐤=(0,0)`$ and $`𝐤=(\pm 2p_F,0)`$. If $`J_3=0`$, the singularity is logarithmic at $`𝐤=(\pm 2p_F,0)`$; if $`J_30`$, the singularity is removed if $`J_3>0`$, while it is enhanced to a power like singularity if $`J_3<0`$, see item d) in the Theorem (1.5). Hence, the singularity at $`𝐤=(\pm 2p_F,0)`$ is of the same type as in the Luttinger model, see (1.10). However, we can not conclude that the same is true for the Fourier transform at $`𝐤=0`$, which is bounded in the Luttinger model, while we can not exclude a logarithmic divergence. In order to get such a stronger result, it would be sufficient to prove that the function $`\mathrm{\Omega }^{3,b}(𝐱)`$ is odd in the exchange of $`(x,x_0)`$ with $`(x_0v,x/v)`$, for some $`v`$; this property is true for the leading term corresponding to $`\mathrm{\Omega }^{3,b}(𝐱)`$ in (1.10), with $`v=v_0`$, but seems impossible to prove on the base of our expansion. We can only see this symmetry for the leading term, with $`v=v_0^{}`$ (or any other value $`v`$ differing for terms of order $`J_3`$, since the substitution of $`v_0^{}`$ with $`v`$ would not affect the bound (1.21)), but this is only sufficient to prove that the singularity has to be of order $`J_3`$, at least. i) Our theorem cannot be proved by building a multiscale renormalized expansion, neither by taking as the “free model” the $`XY`$ one and $`J_3`$ as the perturbative parameter, nor by taking as the free model the $`XXY`$ one and $`u`$ as the perturbative parameter. In fact, in order to solve the model, one cannot perform a single Bogoliubov transformation as in the $`J_3=0`$ case; the gap has a non trivial flow and one has to perform a different Bogoliubov transformation for each renormalization group integration. This can be seen clearly in $`\mathrm{}`$2.66, which is the fermionic integration of a fermionic theory with gap $`\sigma _h`$ and wave function renormalization $`Z_h`$. If $`J_3=0`$, then $`\sigma _h=u`$ and $`Z_h=1`$, but, if $`J_30`$, they are rapidly varying functions of $`h`$. l) If $`u=0`$, the critical indices and $`\nu `$ can be computed with any prefixed precision; we write explicitly in the theorem only the first order for simplicity. However, if $`u0`$, they are not fixed uniquely; for what concerns $`\nu `$, this means that, in the gapped case, the system is insensitive to variations of the magnetic field much smaller than the gap size. m) There is no reason to restrict the analysis to a nearest-neighbor Hamiltonian like (1.1); it will be clear in the following that our results still holds for non nearest-neighbor spin hamiltonians, as such hamiltonians differ from (1.1) for irrelevant (in the RG sense) terms; see also \[Spe\], where the case $`J_3=0`$ is studied. n) The same techniques could perhaps be used to study $`\mathrm{\Omega }_{L,\beta }^1(𝐱)`$ and $`\mathrm{\Omega }_{L,\beta }^2(𝐱)`$, however this problem is more difficult, as one has to study the average of the exponential of the sum of fermionic density operators, see(1.4). In the $`J_3=0`$ case the evaluation of $`\mathrm{\Omega }_{L,\beta }^1(𝐱)`$ and $`\mathrm{\Omega }_{L,\beta }^2(𝐱)`$ was done in \[Mc\]. 1.7 The proof of the theorem is organized into two parts. In the present paper we define a Renormalization Group expansion for the effective potential and the ground state energy of the $`XYZ`$ model, see §(2). One has to perform a multiscale analysis with a different Bogoliubov transformation for each renormalization group integration. A definition of localization operator is introduced, which is different with respect to the one suggested by a simple power counting argument. In §(3) we prove that such expansion is convergent if the running coupling constants are small enough. Despite we are interested in $`\mathrm{\Omega }_𝐱^3`$, we study in detail the convergence of the effective potential and the ground state energy for pedagogical reasons as the expansion for $`\mathrm{\Omega }_𝐱^3`$ is clearer once the expansion for the effective potential is understood. The proof of the convergence requires some care as the power counting has to be improved. Moreover we pay attention to perform all the estimates taking finite $`L,\beta `$; this requires some care, as the preceding analysis of similar problems were not so careful about this point. While in this paper we deal essentially with convergence problems of the renormalized expansions, in the subsequent one we have to analyze carefully the expansions in order to exploit the cancellations, based on symmetry properties, which allow to complete the proof of the theorem; the convergence of the expansion for $`\mathrm{\Omega }_{L,\beta }^3(𝐱)`$ is a corollary of the analogous proof given in this paper for the effective potential or the ground state energy. 2. Multiscale decomposition and anomalous integration 2.1 The Hamiltonian (1.5) can be written, if $`U_L^2`$ is chosen as explained in §1.1 and the definitions (1.11) are used, in the following way (by neglecting a constant term): $$\begin{array}{cc}\hfill H& =\underset{x\mathrm{\Lambda }}{}\{(\mathrm{cos}p_F+\nu )a_x^+a_x^{}\frac{1}{2}[a_x^+a_{x+1}^{}+a_{x+1}^+a_x^{}]\hfill \\ & \frac{u}{2}[a_x^+a_{x+1}^++a_{x+1}^{}a_x^{}]+\lambda (a_x^+a_x^{})(a_{x+1}^+a_{x+1}^{})\},\hfill \end{array}$$ $`(2.1)`$ where $`\mathrm{\Lambda }`$ is an interval of $`L`$ points on the one-dimensional lattice of step one, which will chosen equal to $`([L/2],[(L1)/2])`$, the fermionic field $`a_x^\pm `$ satisfies periodic boundary conditions and $$\lambda =J_3.$$ $`(2.2)`$ The Hamiltonian (2.1) will be considered as a perturbation of the Hamiltonian $`H_0`$ of a system of free fermions in $`\mathrm{\Lambda }`$ with unit mass and chemical potential $`\mu =1\mathrm{cos}p_F`$ ($`u=J_3=\nu =0`$); $`p_F`$ is the Fermi momentum. This system will have, at zero temperature, density $`\rho =p_F/\pi `$, corresponding to magnetization $`\rho 1/2`$ in the $`3`$-direction for the original spin system. Since $`p_F`$ is not uniquely defined at finite volume, we choose it so that $$p_F=\frac{2\pi }{L}(n_F+\frac{1}{2}),n_F\text{},\underset{L\mathrm{}}{lim}p_F=\pi \rho $$ $`(2.3)`$ This means, in particular, that $`p_F`$ is not an allowed momentum of the fermions. We consider also the operators $`a_𝐱^\pm =e^{x_0H}a_x^\pm e^{Hx_0}`$, with $$𝐱=(x,x_0),\beta /2x_0\beta /2,$$ $`(2.4)`$ for some $`\beta >0`$; on $`x_0`$, which we shall call the time variable, antiperiodic boundary conditions are imposed. Many interesting physical properties of the fermionic system at inverse temperature $`\beta `$ can be expressed in terms of the Schwinger functions, that is the truncated expectations in the Grand Canonical Ensemble of the time order product of the field $`a_𝐱^\pm `$ at different space-time points. There is of course a relation between these functions and the expectations of some suitable observables in the spin system. However, by looking at (1.4), one sees that this relation is simple enough only in the case of the truncated expectations of the time order product of the fermionic density operator $`\rho _𝐱=a_𝐱^+a_𝐱^{}`$ at different space-time points, which we shall call the density Schwinger functions; they coincide with the truncated expectations of the time order product of the operator $`S_𝐱^3=e^{x_0H}S_x^3e^{Hx_0}`$ at different space-time points. As it is well known, the Schwinger functions can be written as power series in $`\lambda `$ and $`u`$, convergent for $`\|\lambda \|,\|u\|\epsilon _\beta `$, for some constant $`\epsilon _\beta `$ (the only trivial bound of $`\epsilon _\beta `$ goes to zero, as $`\beta \mathrm{}`$). This power expansion is constructed in the usual way in terms of Feynman graphs, by using as free propagator the function $$\begin{array}{cc}\hfill g^{L,\beta }(𝐱𝐲)& =\frac{\mathrm{Tr}\left[e^{\beta H_0}𝐓(a_𝐱^{}a_𝐲^+)\right]}{\mathrm{Tr}[e^{\beta H_0}]}=\hfill \\ & =\frac{1}{L}\underset{k𝒟_L}{}e^{ik(xy)}\left\{\frac{e^{\tau e(k)}}{1+e^{\beta e(k)}}\text{1}(\tau >0)\frac{e^{(\beta +\tau )e(k)}}{1+e^{\beta e(k)}}\text{1}(\tau 0)\right\},\hfill \end{array}$$ $`(2.5)`$ where $`𝐓`$ is the time order product, $`N=_{x\mathrm{\Lambda }}a_x^+a_x^+`$, $`\tau =x_0y_0`$, $`\text{1}(E)`$ denotes the indicator function ($`\text{1}(E)=1`$, if $`E`$ is true, $`\text{1}(E)=0`$ otherwise), $$e(k)=\mathrm{cos}p_F\mathrm{cos}k,$$ $`(2.6)`$ and $`𝒟_L\{k=2\pi n/L,n\text{},[L/2]n[(L1)/2]\}`$. It is also well known that, if $`x_0y_0`$, $`g^{L,\beta }(𝐱𝐲)=lim_M\mathrm{}g^{L,\beta ,M}(𝐱𝐲)`$, where $$g^{L,\beta ,M}(𝐱𝐲)=\frac{1}{L\beta }\underset{𝐤𝒟_{L,\beta }}{}\frac{e^{i𝐤(𝐱𝐲)}}{ik_0+\mathrm{cos}p_F\mathrm{cos}k},$$ $`(2.7)`$ $`𝐤=(k,k_0)`$, $`𝐤𝐱=k_0x_0+kx`$, $`𝒟_{L,\beta }𝒟_L\times 𝒟_\beta `$, $`𝒟_\beta \{k_0=2(n+1/2)\pi /\beta ,nZ,MnM1\}`$. Note that $`g^{L,\beta ,M}(𝐱𝐲)`$ is real, $`M`$. Hence, if we introduce a finite set of Grassmanian variables $`\{\widehat{a}_𝐤^\pm \}`$, one for each $`𝐤𝒟_{L,\beta }`$, and a linear functional $`P(da)`$ on the generated Grassmanian algebra, such that $$P(da)\widehat{a}_{𝐤_1}^{}\widehat{a}_{𝐤_2}^+=L\beta \delta _{𝐤_1,𝐤_2}\widehat{g}(𝐤_1),\widehat{g}(𝐤)=\frac{1}{ik_0+\mathrm{cos}p_F\mathrm{cos}k},$$ $`(2.8)`$ we have $$\underset{M\mathrm{}}{lim}\frac{1}{L\beta }\underset{𝐤𝒟_{L,\beta }}{}e^{i𝐤(𝐱𝐲)}\widehat{g}(𝐤)=\underset{M\mathrm{}}{lim}P(da)a_𝐱^{}a_𝐲^+g^{L,\beta }(𝐱;𝐲),$$ $`(2.9)`$ where the Grassmanian field $`a_𝐱`$ is defined by $$a_𝐱^\pm =\frac{1}{L\beta }\underset{𝐤𝒟_{L,\beta }}{}\widehat{a}_𝐤^\pm e^{\pm i𝐤𝐱}.$$ $`(2.10)`$ The “Gaussian measure” $`P(da)`$ has a simple representation in terms of the “Lebesgue Grassmanian measure” $`_{𝐤𝒟_{L,\beta }}da_𝐤^+da_𝐤^{}`$, defined as the linear functional on the Grassmanian algebra, such that, given a monomial $`Q(a^{},a^+)`$ in the variables $`a_𝐤^{},a_𝐤^+`$, $`𝐤𝒟_{L,\beta }`$, its value is $`0`$, except in the case $`Q(a^{},a^+)=_𝐤\widehat{a}_𝐤^{}\widehat{a}_𝐤^+`$, up to a permutation of the variables. In this case the value of the functional is determined, by using the anticommuting properties of the variables, by the condition $$\left[\underset{𝐤𝒟_{L,\beta }}{}da_𝐤^+da_𝐤^{}\right]\underset{𝐤𝒟_{L,\beta }}{}\widehat{a}_𝐤^{}\widehat{a}_𝐤^+=1$$ $`(2.11)`$ We have $$P(da)=\left\{\underset{𝐤}{}(L\beta \widehat{g}_𝐤)\widehat{a}_𝐤^+\widehat{a}_𝐤^{}\right\}\mathrm{exp}\left\{\underset{𝐤}{}(L\beta \widehat{g}_𝐤)^1\widehat{a}_𝐤^+\widehat{a}_𝐤^{}\right\}.$$ $`(2.12)`$ Note that, since $`(\widehat{a}_𝐤^{})^2=(\widehat{a}_𝐤^+)^2=0`$, $`e^{z\widehat{a}_𝐤^+\widehat{a}_𝐤}=1z\widehat{a}_𝐤^+\widehat{a}_𝐤`$, for any complex $`z`$. Remark. The ultraviolet cutoff $`M`$ on the $`k_0`$ variable was introduced so that the Grassman algebra is finite; this implies that the Grassmanian integration is indeed a simple algebraic operation and all quantities that appear in the calculations are finite sums. However, $`M`$ does not play any essential role in this paper, since all bounds will be uniform with respect to $`M`$ and they easily imply the existence of the limit. Hence, we shall not stress the dependence on $`M`$ of the various quantities we shall study. By using standard arguments (see, for example, \[NO\], where a different regularization of the propagator is used), one can show that the partition function and the Schwinger functions can be calculated as expectations of suitable functions of the Grassmanian field with respect to the “Gaussian measure” $`P(da)`$. In particular the partition function $`\mathrm{Tr}[e^{\beta H}]`$ is equal to $`𝒵_{L,\beta }\mathrm{Tr}[e^{\beta H_0}]`$, with $$𝒵_{L,\beta }=P(da)e^{𝒱(a)},$$ $`(2.13)`$ where $$\begin{array}{ccc}& 𝒱(a)=uV_u(a)+\lambda V_\lambda (a)+\nu N(a),\hfill & \\ \hfill V_\lambda (a)& =\underset{x,y\mathrm{\Lambda }}{}_{\beta /2}^{\beta /2}𝑑x_0_{\beta /2}^{\beta /2}𝑑y_0v_\lambda (𝐱𝐲)a_𝐱^+a_𝐲^+a_𝐲^{}a_𝐱^{},N(a)=\underset{x\mathrm{\Lambda }}{}_{\beta /2}^{\beta /2}𝑑x_0a_𝐱^+a_𝐱^{},\hfill & \\ \hfill V_u(a)& =\underset{x,y\mathrm{\Lambda }}{}_{\beta /2}^{\beta /2}𝑑x_0_{\beta /2}^{\beta /2}𝑑y_0v_u(𝐱𝐲)\left[a_𝐱^+a_𝐲^+a_𝐱^{}a_𝐲^{}\right]\hfill & (2.14)\hfill \end{array}$$ where $$v_\lambda (𝐱𝐲)=\frac{1}{2}\delta _{1,\|xy\|}\delta (x_0y_0),v_u(𝐱𝐲)=\frac{1}{2}\delta _{x,y+1}\delta (x_0y_0).$$ $`(2.15)`$ Note that the parameter $`\nu `$ has been introduced in order to fix the singularities of the interacting propagator to the values of the free model, that is $`𝐤=(0,\pm p_F)`$. Hence $`\nu `$ is a function of $`\lambda ,u,p_F`$, which has to be fixed so that the perturbation expansion is convergent (uniformly in $`L,\beta `$). This choice of $`\nu `$ has also the effect of fixing the singularities of the spin correlation function Fourier transform, as we explained in the introduction, see §1.4. Note that, if $`p_F=\pi /2`$, one can prove that $`\nu =\lambda `$, by using simple symmetry properties of our expansion; this implies, by using (1.11), that $`h=0`$. If $`u=0`$, it is conjectured, on the base of heuristic calculations, that this condition is equivalent to the condition that, in the limit $`L,\beta \mathrm{}`$, the density is fixed (“Luttinger Theorem”) to the free model value $`\rho =p_F/\pi `$. If $`u0`$, there is no simple relation between the value of $`p_F`$ and the density, as one can see directly in the case $`\lambda =0`$, where one can do explicit calculations. 2.2 We shall begin our analysis by rewriting the potential $`𝒱(a)`$ as $$𝒱(a)=𝒱^{(1)}(a)+uV_u(a)+\delta ^{}V_\delta (a),$$ $`(2.16)`$ where $$𝒱^{(1)}(a)=\lambda V_\lambda (a)+\nu N(a)\delta ^{}V_\delta (a),$$ $`(2.17)`$ and $$V_\delta (a)=\frac{1}{L\beta }\underset{𝐤}{}e(k)\widehat{a}_𝐤^+\widehat{a}_𝐤^{}.$$ $`(2.18)`$ $`\delta ^{}`$ is an arbitrary parameter, to be fixed later, of modulus smaller than $`1/2`$; its introduction is not really necessary, but allows to simplify the discussion of the spin correlation function asymptotic behaviour. In terms of the Fermionic system, it will describe the modification of the Fermi velocity due to the interaction. Afterwards we “move” the terms $`uV_u(a)`$ and $`\delta ^{}V_\delta (a)`$ from the interaction to the Gaussian measure. In order to describe the properties of the new Gaussian measure, it is convenient to introduce a new set of Grassmanian variables $`\widehat{b}_{𝐤,\omega }^\sigma `$, $`\omega =\pm 1`$, $`𝐤𝒟_{L,\beta }^+`$, by defining $$𝒟_{L,\beta }^\omega =\{𝐤𝒟_{L,\beta }:\omega k>0\}\{𝐤𝒟_{L,\beta }:k=0,\omega k_0>0\},$$ $`(2.19)`$ $$\widehat{b}_{𝐤,\omega }^\sigma =\widehat{a}_{\omega 𝐤}^{\sigma \omega },$$ $`(2.20)`$ so that, by using (2.10) $$a_𝐱^\sigma =\frac{1}{L\beta }\underset{𝐤𝒟_{L,\beta }^+,\omega =\pm 1}{}\widehat{b}_{𝐤,\omega }^{\sigma \omega }e^{i\sigma \omega 𝐤𝐱}.$$ $`(2.21)`$ It is easy to see that $$𝒵_{L,\beta }=e^{L\beta t_1}P(db)e^{\stackrel{~}{𝒱}^{(1)}(b)},$$ $`(2.22)`$ with $`\stackrel{~}{𝒱}^{(1)}(b)=𝒱^{(1)}(a)`$, where $`a`$ has to be interpreted as the r.h.s. of (2.21), $$\begin{array}{cc}\hfill P(db)& =\left\{\underset{𝐤𝒟_{L,\beta }^+}{}\frac{(L\beta )^2}{k_0^2(1+\delta ^{})^2e(k)^2u^2\mathrm{sin}^2k}\underset{\omega =\pm 1}{}\widehat{b}_{𝐤,\omega }^+\widehat{b}_{𝐤,\omega }^{}\right\}\hfill \\ & \mathrm{exp}\left\{\frac{1}{L\beta }\underset{𝐤𝒟_{L,\beta }^+}{}\underset{\omega ,\omega ^{}}{}\widehat{b}_{𝐤,\omega }^+T_{\omega ,\omega ^{}}(𝐤)\widehat{b}_{𝐤,\omega }^{}\right\},\hfill \end{array}$$ $`(2.23)`$ $$T(𝐤)=\left(\begin{array}{cc}ik_0+(1+\delta ^{})e(k)& iu\mathrm{sin}k\\ iu\mathrm{sin}k& ik_0(1+\delta ^{})e(k)\end{array}\right),$$ $`(2.24)`$ $$t_1=\frac{1}{L\beta }\underset{𝐤𝒟_{L,\beta }^+}{}\mathrm{log}\frac{k_0^2+(1+\delta ^{})e(k)^2+u^2\mathrm{sin}^2k}{k_0^2+e(k)^2}.$$ $`(2.25)`$ Note that $`t_1`$ is uniformly bounded as $`L,\beta \mathrm{}`$, if $`\|\delta ^{}\|1/2`$, as we are supposing. For $`\lambda =\nu =\delta ^{}=0`$, it represents the free energy for lattice site of $`HH_0`$. 2.3 For $`\lambda =\nu =0`$, all the properties of the model can be analyzed in terms of the Grassmanian measure (2.23). In particular, we have $$P(db)a_𝐱^{\sigma _1}a_𝐲^{\sigma _2}=\frac{1}{L\beta }\underset{𝐤𝒟_{L,\beta }^+}{}\left[e^{i𝐤(𝐱𝐲)}T^1(𝐤)_{\sigma _1,\sigma _2}e^{i𝐤(𝐱𝐲)}T^1(𝐤)_{\sigma _2,\sigma _1}\right],$$ $`(2.26)`$ where $`T^1(𝐤)`$ denotes the inverse of the matrix $`T(𝐤)`$. This matrix is defined for any $`𝐤𝒟_{L,\beta }`$ and satisfies the symmetry relation $$T^1(𝐤)_{\sigma _2,\sigma _1}=T^1(𝐤)_{\sigma _1,\sigma _2},$$ $`(2.27)`$ so that we can write (2.26) also in the form $$P(db)a_𝐱^{\sigma _1}a_𝐲^{\sigma _2}=\frac{1}{L\beta }\underset{𝐤𝒟_{L,\beta }}{}e^{i𝐤(𝐱𝐲)}T^1(𝐤)_{\sigma _1,\sigma _2}.$$ $`(2.28)`$ If $`\lambda 0`$, we shall study the model, for $`\lambda `$ small, in terms of a perturbative expansion, based on a multiscale decomposition of the measure (2.23), by using the methods introduced in \[BG\] and extended in various other papers (\[BGPS\], \[BM1\], \[M1\]). In order to discuss the structure of the expansion, it is convenient to explain first how it works in the case of the free energy for site of $`HH_0`$ $$E_{L,\beta }=\frac{1}{L\beta }\mathrm{log}𝒵_{L,\beta }.$$ $`(2.29)`$ Let $`T^1`$ be the one dimensional torus, $`kk^{}_{T^1}`$ the usual distance between $`k`$ and $`k^{}`$ in $`T^1`$ and $`k=k0`$. We introduce a scaling parameter $`\gamma >1`$ and a positive function $`\chi (𝐤^{})C^{\mathrm{}}(T^1\times R)`$, $`𝐤^{}=(k^{},k_0)`$, such that $$\chi (𝐤^{})=\chi (𝐤^{})=\{\begin{array}{cc}1\hfill & \text{if }\|𝐤^{}\|<t_0a_0v_0^{}/\gamma ,\hfill \\ 0\hfill & \text{if }\|𝐤^{}\|>a_0v_0^{},\hfill \end{array}$$ $`(2.30)`$ where $$\|𝐤^{}\|=\sqrt{k_0^2+(v_0^{}k^{}_{T^1})^2},$$ $`(2.31)`$ $$a_0=\mathrm{min}\{p_F/2,(\pi p_F)/2\},$$ $`(2.32)`$ $$v_0^{}=v_0(1+\delta ^{}),v_0=\mathrm{sin}p_F.$$ $`(2.33)`$ In order to give a well defined meaning to the definition (2.30), $`v_0^{}>0`$ has to be positive. Hence we shall suppose that $$v_0\overline{v}_0>0,\|\delta ^{}\|\frac{1}{2},$$ $`(2.34)`$ where $`\overline{v}_0`$ is fixed once for all. All our results will be uniform in $`v_0`$, under the conditions (2.34), but we shall not stress this fact anymore in the following. The definition (2.30) is such that the supports of $`\chi (kp_F,k_0)`$ and $`\chi (k+p_F,k_0)`$ are disjoint and the $`C^{\mathrm{}}`$ function on $`T^1\times R`$ $$\widehat{f}_1(𝐤)1\chi (kp_F,k_0)\chi (k+p_F,k_0)$$ $`(2.35)`$ is equal to $`0`$, if $`[v_0^{}(\|k\|p_F)_{T^1}]^2+k_0^2<t_0^2`$. We define also, for any integer $`h0`$, $$f_h(𝐤^{})=\chi (\gamma ^h𝐤^{})\chi (\gamma ^{h+1}𝐤^{});$$ $`(2.36)`$ we have, for any $`\overline{h}<0`$, $$\chi (𝐤^{})=\underset{h=\overline{h}+1}{\overset{0}{}}f_h(𝐤^{})+\chi (\gamma ^{\overline{h}}𝐤^{}).$$ $`(2.37)`$ Note that, if $`h0`$, $`f_h(𝐤^{})=0`$ for $`\|𝐤^{}\|<t_0\gamma ^{h1}`$ or $`\|𝐤^{}\|>t_0\gamma ^{h+1}`$, and $`f_h(𝐤^{})=1`$, if $`\|𝐤^{}\|=t_0\gamma ^h`$, so that $$f_{h_1}(𝐤^{})f_{h_2}(𝐤^{})=0,\text{if }\|h_1h_2\|>1.$$ $`(2.38)`$ We finally define, for any $`h0`$: $$\widehat{f}_h(𝐤)=f_h(kp_F,k_0)+f_h(k+p_F,k_0);$$ $`(2.39)`$ This definition implies that, if $`h0`$, the support of $`\widehat{f}_h(𝐤)`$ is the union of two disjoint sets, $`A_h^+`$ and $`A_h^{}`$. In $`A_h^+`$, $`k`$ is strictly positive and $`kp_F_{T^1}a_0\gamma ^ha_0`$, while, in $`A_h^{}`$, $`k`$ is strictly negative and $`k+p_F_{T^1}a_0\gamma ^h`$. The label $`h`$ is called the scale or frequency label. Note that, if $`𝐤𝒟_{L,\beta }`$, then $`\|𝐤\pm (p_F,0)\|\sqrt{(\pi \beta ^1)^2+(v_0^{}\pi L^1)^2}`$, by (2.3) and the definition of $`𝒟_{L,\beta }`$. Therefore $$\widehat{f}_h(𝐤)=0h<h_{L,\beta }=\mathrm{min}\{h:t_0\gamma ^{h+1}>\sqrt{(\pi \beta ^1)^2+(v_0^{}\pi L^1)^2}\},$$ $`(2.40)`$ and, if $`𝐤𝒟_{L,\beta }`$, the definitions (2.35) and (2.39), together with the identity (2.37), imply that $$1=\underset{h=h_{L,\beta }}{\overset{1}{}}\widehat{f}_h(𝐤).$$ $`(2.41)`$ We now introduce, for each scale label $`h`$, such that $`h_{L,\beta }h1`$, a set of Grassmanian variables $`b_{𝐤,\omega }^{(h)\sigma }`$ and a corresponding Gaussian measure $`P(db^{(h)})`$, such that, if $`h=1`$, then $`𝐤𝒟_{L,\beta }`$ and $$P(db^{(1)})b_{𝐤_1,\omega _1}^{(1)\sigma _1}b_{𝐤_2,\omega _2}^{(1)\sigma _2}=L\beta \sigma _1\delta _{\sigma _1,\sigma _2}\delta _{𝐤_1,𝐤_2}\frac{1}{2}T^1(𝐤_1)_{\omega _1,\omega _2}\widehat{f}_1(𝐤_1),$$ $`(2.42)`$ while, if $`h0`$, then $`𝐤𝒟_{L,\beta }^+`$ and $$P(db^{(h)})b_{𝐤_1,\omega _1}^{(h)\sigma _1}b_{𝐤_2,\omega _2}^{(h)\sigma _2}=L\beta \sigma _1\delta _{\sigma _1,\sigma _2}\delta _{𝐤_1,𝐤_2}T^1(𝐤_1)_{\omega _1,\omega _2}f_h(k_1p_F,k_0).$$ $`(2.43)`$ The support properties of the r.h.s. of (2.42) and (2.43) allow to impose the condition $$b_{𝐤,\omega }^{(h)\sigma }=0,\text{if}\widehat{f}_h(𝐤)=0.$$ $`(2.44)`$ By using (2.26) and (2.27), it is easy to see that $$P(db)a_𝐱^{\sigma _1}a_𝐲^{\sigma _2}=\underset{h=h_{L,\beta }}{\overset{1}{}}\underset{\omega _1,\omega _2}{}P(db^{(h)})b_{𝐱,\omega _1}^{(h)\sigma _1\omega _1}b_{𝐲,\omega _2}^{(h)\sigma _2\omega _2},$$ $`(2.45)`$ where, if $`h0`$, $$b_{𝐱,\omega }^{(h)\sigma }=\frac{1}{L\beta }\underset{𝐤𝒟_{L,\beta }^+}{}\widehat{b}_{𝐤,\omega }^{(h)\sigma }e^{i\sigma 𝐤𝐱},$$ $`(2.46)`$ while, if $`h=1`$, a similar definition is used, with $`𝒟_{L,\beta }`$ in place of $`𝒟_{L,\beta }^+`$. Note that this different definition, which is at the origin of the factor $`1/2`$ in the r.h.s. of (2.42), is not really necessary, but implies that $`P(db^{(1)})b_{𝐱,\omega _1}^{(1)}b_{𝐲,\omega _2}^{(h)+}`$ is bounded for $`M\mathrm{}`$, a property which should otherwise be true only for $`_{\omega _1,\omega _2}P(db^{(1)})b_{𝐱,\omega _1}^{(1)\sigma _1\omega _1}b_{𝐲,\omega _2}^{(h)\sigma _2\omega _2}`$. In the following, we shall use this property in order to simplify the discussion in some minor points. The identity (2.45), as it is well known, implies that, if $`F(a)`$ is any function of the variables $`a_𝐱^\sigma `$, then $$P(da)F(a)=\underset{h=h_{L,\beta }}{\overset{1}{}}P(db^{(h)})F\left(\underset{h=h_{L,\beta }}{\overset{1}{}}a^{(h)}\right),$$ $`(2.47)`$ where $$a_𝐱^{(h)\sigma }=\underset{\omega =\pm 1}{}b_{𝐱,\omega }^{(h)\sigma \omega }.$$ $`(2.48)`$ It is now convenient to introduce a variable which measures the distance of the momentum from the Fermi surface, by putting $`k=k^{}+p_F`$, with $`k^{}𝒟_L^{}=\{k^{}=2(n+1/2)\pi /L,n\text{},[L/2]n[(L1)/2]\}`$. Moreover, we rename the Grassmanian variables, by defining $$\widehat{\psi }_{𝐤^{},\omega }^{(h)\sigma }=\widehat{b}_{𝐤^{}+𝐩_F,\omega }^{(h)\sigma },\psi _{𝐱,\omega }^{(h)\sigma }=\frac{1}{L\beta }\underset{𝐤^{}𝒟_{L,\beta }^{}}{}e^{i\sigma 𝐤^{}𝐱}\widehat{\psi }_{𝐤^{},\omega }^{(h)\sigma },$$ $`(2.49)`$ where $`𝒟_{L,\beta }^{}=𝒟_L^{}\times 𝒟_\beta `$, $`𝐤^{}=(k^{},k_0)`$ and $`𝐩_F=(p_F,0)`$. Note that, by (2.44), $$\widehat{\psi }_{𝐤^{},\omega }^{(h)\sigma }=0\text{ if }\widehat{f}_h(𝐤^{}+𝐩_F)=0.$$ $`(2.50)`$ The definition (2.49) allows to write (2.48) in the form $$a_𝐱^{(h)\sigma }=\underset{\omega }{}e^{i\sigma \omega 𝐩_F𝐱}\psi _{𝐱,\omega }^{(h)\sigma \omega }.$$ $`(2.51)`$ The measure $`P(db^{(h)})`$ can be thought in a natural way as a measure on the variables $`\psi _{𝐱,\omega }^{(h)\sigma }`$, that we shall denote $`P(d\psi ^{(h)})`$. Then, (2.43) and (2.49) imply that, if $`h1`$, $$P(d\psi ^{(h)})\widehat{\psi }_{𝐤_1^{},\omega _1}^{(h)\sigma _1}\widehat{\psi }_{𝐤_2^{},\omega _2}^{(h)\sigma _2}=(1\frac{1}{2}\delta _{h,1})L\beta \sigma _1\delta _{\sigma _1,\sigma _2}\delta _{𝐤_1^{},𝐤_2^{}}\stackrel{~}{g}_{\omega _1,\omega _2}^{(h)}(𝐤_1^{}),$$ $`(2.52)`$ where, if $`f_1(𝐤^{})\widehat{f}_1(𝐤^{}+𝐩_F)`$, $$\stackrel{~}{g}^{(h)}(𝐤^{})=\frac{f_h(𝐤^{})}{k_0^2E(k^{})^2u^2\mathrm{sin}^2(k^{}+p_F)}\left(\begin{array}{cc}ik_0E(k^{})& iu\mathrm{sin}(k^{}+p_F)\\ iu\mathrm{sin}(k^{}+p_F)& ik_0+E(k^{})\end{array}\right),$$ $`(2.53)`$ $$E(k^{})=v_0^{}\mathrm{sin}k^{}+(1+\delta ^{})(1\mathrm{cos}k^{})\mathrm{cos}p_F.$$ $`(2.54)`$ In the following we shall use also the notation $$\psi _{𝐱,\omega }^{(h)\sigma }=\underset{h^{}=h_{L,\beta }}{\overset{h}{}}\psi _{𝐱,\omega }^{(h^{})\sigma },P(d\psi ^{(h)})=\underset{h^{}=h_{L,\beta }}{\overset{h}{}}P(d\psi ^{(h^{})}),$$ $`(2.55)`$ which allows to write the identity (2.47) as $$P(da)F(a)=P(d\psi ^{(1)})\stackrel{~}{F}(\psi ^{(1)}),$$ $`(2.56)`$ where $`\stackrel{~}{F}(\psi ^{(1)})`$ is obtained from $`F(_ha^{(h)})`$, by using (2.51). Remark. Note that the sum over $`k_0`$ in (2.49) can be thought as a finite sum for any $`M`$, if $`h0`$, because of the support properties of $`\widehat{\psi }_{𝐤^{},\omega }^{(h)\sigma }`$. Hence, all quantities that we shall calculate will depend on $`M`$ only trough the propagator $`\stackrel{~}{g}^{(1)}(𝐤^{})`$, if $`M`$ is large enough. 2.4 If we apply (2.56) to $`𝒵_{L,\beta }`$ and we use (2.29) and (2.22), we get $$e^{L\beta E_{L,\beta }}=e^{L\beta t_1}P(d\psi ^{(1)})e^{𝒱^{(1)}(\psi ^{(1)})},$$ $`(2.57)`$ where $$𝒱^{(1)}(\psi ^{(1)})=\lambda V_\lambda (\underset{h=h_{L,\beta }}{\overset{1}{}}a^{(h)})+\nu N(\underset{h=h_{L,\beta }}{\overset{1}{}}a^{(h)})\delta ^{}V_\delta (\underset{h=h_{L,\beta }}{\overset{1}{}}a^{(h)}).$$ $`(2.58)`$ Let us now perform the integration over $`\psi ^{(1)}`$; we get $$e^{L\beta E_{L,\beta }}=e^{L\beta (\stackrel{~}{E}_1+t_1)}P(d\psi ^{(0)})e^{\overline{𝒱}^{(0)}(\psi ^{(0)})},\overline{𝒱}^{(0)}(0)=0,$$ $`(2.59)`$ $$e^{\overline{𝒱}^{(0)}(\psi ^{(0)})L\beta \stackrel{~}{E}_1}=P(d\psi ^{(+1)})e^{𝒱^{(1)}(\psi ^{(0)}+\psi ^{(+1)}}).$$ $`(2.60)`$ This step is essentially trivial. In fact, it is easy to see that $`\stackrel{~}{g}_{\omega ,\omega ^{}}^{(1)}(𝐤^{})`$ is bounded, for $`M\mathrm{}`$, uniformly in $`L,\beta `$, and that its Fourier transform $`\stackrel{~}{g}_{\omega ,\omega ^{}}^{(1)}(𝐱)`$ is a bounded function with fast decaying properties (uniformly in $`L,\beta `$). Hence, by using standard perturbation theory, it is easy to see that $`\overline{𝒱}^{(0)}(\psi ^{(0)})`$ can be written in the form $$\begin{array}{cc}& \overline{𝒱}^{(0)}(\psi ^{(0)})=\underset{n=1}{\overset{\mathrm{}}{}}\frac{1}{(L\beta )^{2n}}\underset{\underset{¯}{\sigma },\underset{¯}{\omega }}{}\underset{𝐤_1^{},\mathrm{},𝐤_{2n}^{}}{}\underset{i=1}{\overset{2n}{}}\widehat{\psi }_{𝐤_i^{},\omega _i}^{(0)\sigma _i}\hfill \\ & \widehat{W}_{2n,\underset{¯}{\sigma },\underset{¯}{\omega }}^{(0)}(𝐤_1^{},\mathrm{},𝐤_{2n1}^{})\delta (\underset{i=1}{\overset{2n}{}}\sigma _i(𝐤_i^{}+𝐩_F)),\hfill \end{array}$$ $`(2.61)`$ where $`\underset{¯}{\sigma }=(\sigma _1,\mathrm{},\sigma _{2n})`$, $`\underset{¯}{\omega }=(\omega _1,\mathrm{},\omega _{2n})`$ and we used the notation $$\delta (𝐤)=\delta (k)\delta (k_0),\delta (k)=L\underset{n\text{}}{}\delta _{k,2\pi n},\delta (k_0)=\beta \delta _{k_0,0}.$$ $`(2.62)`$ As we shall prove in §$`\mathrm{}`$3, the kernels $`\widehat{W}_{2n,\underset{¯}{\sigma },\underset{¯}{\omega }}^{(0)}(𝐤_1^{},\mathrm{},𝐤_{2n1}^{})`$, as well as $`\stackrel{~}{E}_1`$, are expressed as power series of $`\lambda ,\nu `$, convergent for $`\epsilon Max(\|\lambda \|,\|\nu \|)\epsilon _0`$, for $`\epsilon _0`$ small enough. Moreover there exists a constant $`C`$, such that, uniformly in $`L,\beta `$, $`\|\stackrel{~}{E}_1\|C\epsilon `$ and $`\|\widehat{W}_{2n,\underset{¯}{\sigma },\underset{¯}{\omega }}^{(0)}\|C^n\epsilon ^{\mathrm{max}(1,n1)}`$. Remark - The conservation of momentum and the support property (2.50) of $`\widehat{\psi }_{𝐤^{},\omega }^{(0)\sigma }`$ imply that, if $`n=1`$, only the terms with $`\sigma _1+\sigma _2=0`$ contribute to the sum in (2.61). Let us now define $`𝐤^{}=(k,k_0)`$. It is possible to show, by using the symmetries of the interaction and of the covariance $`\stackrel{~}{g}^{(1)}(𝐤^{})`$, that $$\begin{array}{cc}\hfill \widehat{W}_{n,\underset{¯}{\sigma },\underset{¯}{\omega }}^{(0)}(𝐤_1^{},\mathrm{},𝐤_{n1}^{})& =(1)^{\frac{1}{2}_{i=1}^n\sigma _i\omega _1}[\widehat{W}_{n,\underset{¯}{\sigma },\underset{¯}{\omega }}^{(0)}(𝐤_1,\mathrm{},𝐤_{n1})]^{}=\hfill \\ & =(1)^{\frac{1}{2}_{i=1}^n\sigma _i\omega _i}\widehat{W}_{n,\underset{¯}{\sigma },\underset{¯}{\omega }}^{(0)}(𝐤_1,\mathrm{},𝐤_{n1}).\hfill \end{array}$$ $`(2.63)`$ 2.5 The integration of the fields of scale $`h0`$ is performed iteratively.We define a sequence of positive constants $`Z_h`$, $`h=h_{L,\beta },\mathrm{},0`$, a sequence of effective potentials $`𝒱^{(h)}(\psi )`$, a sequence of constants $`E_h`$ and a sequence of functions $`\sigma _h(𝐤^{})`$, such that $$Z_0=1,E_0=\stackrel{~}{E}_1+t_1,\sigma _0(𝐤^{})=u\mathrm{sin}(k^{}+p_F),$$ $`(2.64)`$ and $$e^{L\beta E_{L,\beta }}=P_{Z_h,\sigma _h,C_h}(d\psi ^{(h)})e^{𝒱^{(h)}(\sqrt{Z_h}\psi ^{(h)})L\beta E_h},𝒱^{(h)}(0)=0,$$ $`(2.65)`$ where $$\begin{array}{cc}& P_{Z_h,\sigma _h,C_h}(d\psi ^{(h)})=\underset{𝐤^{}:C_h^1(𝐤^{})>0}{}\underset{\omega =\pm 1}{}\frac{d\widehat{\psi }_{𝐤^{},\omega }^{(h)+}d\widehat{\psi }_{𝐤^{},\omega }^{(h)}}{𝒩(𝐤^{})}\hfill \\ & \mathrm{exp}\left\{\frac{1}{L\beta }\underset{𝐤^{}:C_h^1(𝐤^{})>0}{}C_h(𝐤^{})Z_h\underset{\omega ,\omega ^{}=\pm 1}{}\widehat{\psi }_{𝐤^{},\omega }^{(h)+}T_{\omega ,\omega ^{}}^{(h+1)}\widehat{\psi }_{𝐤^{},\omega ^{}}^{(h)}\right\},\hfill \end{array}$$ $`(2.66)`$ $$𝒩(𝐤^{})=\frac{C_h(𝐤^{})Z_h}{L\beta }[k_0^2+E(k^{})^2+\sigma _h(𝐤^{})^2]^{1/2},$$ $`(2.67)`$ $$C_h(𝐤^{})^1=\underset{j=h_{L,\beta }}{\overset{h}{}}f_j(𝐤^{}),$$ $`(2.68)`$ and the $`2\times 2`$ matrix $`T_h(𝐤^{})`$ is given by $$T_h(𝐤^{})=\left(\begin{array}{cc}ik_0+E(k^{})& i\sigma _{h1}(𝐤^{})\\ i\sigma _{h1}(𝐤^{})& ik_0E(k^{})\end{array}\right).$$ $`(2.69)`$ We shall also prove that the $`𝒱^{(h)}`$ can be represented as $$\begin{array}{cc}& 𝒱^{(h)}(\psi ^{(h)})=\underset{n=1}{\overset{\mathrm{}}{}}\frac{1}{(L\beta )^{2n}}\underset{\genfrac{}{}{0pt}{}{𝐤_1^{},\mathrm{},𝐤_{2n}^{},}{\underset{¯}{\sigma },\underset{¯}{\omega }}}{}\underset{i=1}{\overset{2n}{}}\widehat{\psi }_{𝐤_i^{},\omega _i}^{(h)\sigma _i}\hfill \\ & \widehat{W}_{2n,\underset{¯}{\sigma },\underset{¯}{\omega }}^{(h)}(𝐤_1^{},\mathrm{},𝐤_{2n1}^{})\delta (\underset{i=1}{\overset{2n}{}}\sigma _i(𝐤_i^{}+𝐩_F)),\hfill \end{array}$$ $`(2.70)`$ with the kernels $`\widehat{W}_{2n,\underset{¯}{\sigma },\underset{¯}{\omega }}^{(h)}`$ verifying the symmetry relations $$\begin{array}{cc}\hfill \widehat{W}_{n,\underset{¯}{\sigma },\underset{¯}{\omega }}^{(h)}(𝐤_1^{},\mathrm{},𝐤_{n1}^{})& =(1)^{\frac{1}{2}_{i=1}^n\sigma _i\omega _i}[\widehat{W}_{n,\underset{¯}{\sigma },\underset{¯}{\omega }}^{(h)}(𝐤_1,\mathrm{},𝐤_{n1})]^{}=\hfill \\ & =(1)^{\frac{1}{2}_{i=1}^n\sigma _i\omega _i}\widehat{W}_{n,\underset{¯}{\sigma },\underset{¯}{\omega }}^{(h)}(𝐤_1,\mathrm{},𝐤_{n1}).\hfill \end{array}$$ $`(2.71)`$ The previous claims are true for $`h=0`$, by (2.59), (2.61), (2.64) and (2.53). In order to prove them for any $`hh_{L,\beta }`$, we must explain how $`𝒱^{(h1)}(\psi )`$ is calculated, given $`𝒱^{(h)}(\psi )`$. It is convenient, for reasons which will be clear below, to split $`𝒱^{(h)}`$ as $`𝒱^{(h)}+𝒱^{(h)}`$, where $`=1`$ and $``$, the localization operator, is a linear operator on functions of the form (2.70), defined in the following way by its action on the kernels $`\widehat{W}_{2n,\underset{¯}{\sigma },\underset{¯}{\omega }}^{(h)}`$. 1) If $`2n=4`$, then $$\widehat{W}_{4,\underset{¯}{\sigma },\underset{¯}{\omega }}^{(h)}(𝐤_1^{},𝐤_2^{},𝐤_3^{})=\widehat{W}_{4,\underset{¯}{\sigma },\underset{¯}{\omega }}^{(h)}(\overline{𝐤}_{++},\overline{𝐤}_{++},\overline{𝐤}_{++}),$$ $`(2.72)`$ where $$\overline{𝐤}_{\eta \eta ^{}}=(\eta \frac{\pi }{L},\eta ^{}\frac{\pi }{\beta }).$$ $`(2.73)`$ Note that this definition depends on the the field variables order in the r.h.s. of (2.70), if $`_{i=1}^4\sigma _i0`$. In fact, since $`\sigma _4𝐤_4^{}=_{i=1}^3\sigma _i𝐤_i^{}𝐩_F_{i=1}^4\sigma _i`$ (modulo $`(2\pi ,0)`$), if $`𝐤_i^{}=\overline{𝐤}_{++}`$ for $`i=1,2,3`$, $`𝐤_4^{}=\overline{𝐤}_{++}`$ only if $`_{i=1}^4\sigma _i=0`$. This is apparently a problem, because the representation (2.70) is not uniquely defined (the terms which differ by a common permutation of the $`\underset{¯}{\sigma }`$ and $`\underset{¯}{\omega }`$ indices are equivalent). However, it is easy to see, by using the anticommuting property of the field variables, that the contribution to $`𝒱^{(h)}`$ of the terms with $`2n=4`$ is equal to $`0`$, unless, after a suitable permutation of the fields, $`\underset{¯}{\sigma }=(+,,+,)`$, $`\underset{¯}{\omega }=(+1,1,1,+1)`$. The previous discussion implies that we are free to change the order of the field variables as we like, before applying the definition (2.72); this freedom will be useful in the construction of the main expansion in §$`\mathrm{}`$3. 2) If $`2n=2`$ and, possibly after a suitable permutation of the fields, $`\underset{¯}{\sigma }=(+,)`$ ($`\sigma _1+\sigma _2=0`$, by the remark following (2.62)), then $$\begin{array}{cc}\hfill \widehat{W}_{2,\underset{¯}{\sigma },\underset{¯}{\omega }}^{(h)}(𝐤^{})& =\frac{1}{4}\underset{\eta ,\eta ^{}=\pm 1}{}\widehat{W}_{2,\underset{¯}{\sigma },\underset{¯}{\omega }}^{(h)}(\overline{𝐤}_{\eta \eta ^{}})\hfill \\ & \left\{1+\delta _{\omega _1,\omega _2}\left[\eta \frac{L}{\pi }\left(b_L+a_L\frac{E(k^{})}{v_0^{}}\right)+\eta ^{}\frac{\beta }{\pi }k_0\right]\right\},\hfill \end{array}$$ $`(2.74)`$ where $$a_L\frac{L}{\pi }\mathrm{sin}\frac{\pi }{L}=1,\frac{\mathrm{cos}p_F}{v_0}(1\mathrm{cos}\frac{\pi }{L})+b_L\frac{L}{\pi }\mathrm{sin}\frac{\pi }{L}=0.$$ $`(2.75)`$ In order to better understand this definition, note that, if $`L=\beta =\mathrm{}`$, $$\widehat{W}_{2,\underset{¯}{\sigma },\underset{¯}{\omega }}^{(h)}(𝐤^{})=\widehat{W}_{2,\underset{¯}{\sigma },\underset{¯}{\omega }}^{(h)}(0)+\delta _{\omega _1,\omega _2}\left[\frac{E(k^{})}{v_0^{}}\frac{\widehat{W}_{2,\underset{¯}{\sigma },\underset{¯}{\omega }}^{(h)}}{k^{}}(0)+k_0\frac{\widehat{W}_{2,\underset{¯}{\sigma },\underset{¯}{\omega }}^{(h)}}{k_0}(0)\right].$$ $`(2.76)`$ Hence, $`\widehat{W}_{2,\underset{¯}{\sigma },\underset{¯}{\omega }}^{(h)}(𝐤^{})`$ has to be understood as a discrete version of the Taylor expansion up to order $`1`$. Since $`a_L=1+O(L^2)`$ and $`b_L=O(L^2)`$, this property would be true also if $`a_L=1`$ and $`b_L=0`$; however the choice (2.75) has the advantage to share with (2.76) another important property, that is $`^2\widehat{W}_{2,\underset{¯}{\sigma },\underset{¯}{\omega }}^{(h)}(𝐤^{})=\widehat{W}_{2,\underset{¯}{\sigma },\underset{¯}{\omega }}^{(h)}(𝐤^{})`$. 3) In all the other cases $$\widehat{W}_{2n,\underset{¯}{\sigma },\underset{¯}{\omega }}^h(𝐤_1^{},\mathrm{},𝐤_{2n1}^{})=0.$$ $`(2.77)`$ By (2.72) and the remark following (2.76), the operator $``$ satisfies the relation $$=0.$$ $`(2.78)`$ By using the anticommuting properties of the Grassmanian variables (see discussion in item 1) above) and the symmetry relations (2.71), we can write $`𝒱^{(h)}`$ in the following way: $$𝒱^{(h)}(\psi ^{(h)})=\gamma ^hn_hF_\nu ^{(h)}+s_hF_\sigma ^{(h)}+z_hF_\zeta ^{(h)}+a_hF_\alpha ^{(h)}+l_hF_\lambda ^{(h)},$$ $`(2.79)`$ where $`n_h`$, $`s_h`$, $`z_h`$, $`a_h`$ and $`l_h`$ are real numbers and $$\begin{array}{ccc}\hfill F_\nu ^{(h)}& =\underset{\omega =\pm 1}{}\frac{\omega }{L\beta }\underset{𝐤^{}𝒟_{L,\beta }^{}}{}\widehat{\psi }_{𝐤^{},\omega }^{(h)+}\widehat{\psi }_{𝐤^{},\omega }^{(h)},\hfill & \\ \hfill F_\sigma ^{(h)}& =\underset{\omega =\pm 1}{}\frac{i\omega }{(L\beta )}\underset{𝐤^{}𝒟_{L,\beta }^{}}{}\widehat{\psi }_{𝐤^{},\omega }^{(h)+}\widehat{\psi }_{𝐤^{},\omega }^{(h)},\hfill & \\ \hfill F_\alpha ^{(h)}& =\underset{\omega =\pm 1}{}\frac{\omega }{(L\beta )}\underset{𝐤^{}𝒟_{L,\beta }^{}}{}\frac{E(k^{})}{v_0^{}}\widehat{\psi }_{𝐤^{},\omega }^{(h)+}\widehat{\psi }_{𝐤^{},\omega }^{(h)},\hfill & (2.80)\hfill \\ \hfill F_\zeta ^{(h)}& =\underset{\omega =\pm 1}{}\frac{1}{(L\beta )}\underset{𝐤^{}𝒟_{L,\beta }^{}}{}(ik_0)\widehat{\psi }_{𝐤^{},\omega }^{(h)+}\widehat{\psi }_{𝐤^{},\omega }^{(h)},\hfill & \\ \hfill F_\lambda ^{(h)}& =\frac{1}{(L\beta )^4}\underset{𝐤_1^{},\mathrm{},𝐤_4^{}𝒟_{L,\beta }^{}}{}\widehat{\psi }_{𝐤_1^{},+1}^{(h)+}\widehat{\psi }_{𝐤_2^{},1}^{(h)}\widehat{\psi }_{𝐤_3^{},1}^{(h)+}\widehat{\psi }_{𝐤_4^{},+1}^{(h)}\delta (𝐤_1^{}𝐤_2^{}+𝐤_3^{}𝐤_4^{}).\hfill & \end{array}$$ By using (2.72) and (2.74), it is easy to see that, if $`\epsilon \mathrm{max}\{\|\lambda \|,\|\nu \|\}`$, $$\begin{array}{cc}\hfill l_0=4\lambda \mathrm{sin}^2(p_F+\pi /L)+O(\epsilon ^2),& a_0=\delta ^{}v_0+c_0^\delta \lambda _1+O(\epsilon ^2),\hfill \\ \hfill s_0=O(u\epsilon ),z_0=O(\epsilon ^2),& n_0=\nu +O(\epsilon ),\hfill \end{array}$$ $`(2.81)`$ where $`c_0^\delta `$ is a constant, bounded uniformly in $`L,\beta `$. We now renormalize the free measure $`P_{Z_h,\sigma _h,C_h}(d\psi ^{(h)})`$, by adding to it part of the r.h.s. of (2.79). We get $$\begin{array}{cc}\hfill P_{Z_h,\sigma _h,C_h}(d\psi ^{(h)})& e^{𝒱^{(h)}(\sqrt{Z_h}\psi ^{(h)})}=\hfill \\ & =e^{L\beta t_h}P_{\stackrel{~}{Z}_{h1},\sigma _{h1},C_h}(d\psi ^{(h)})e^{\stackrel{~}{𝒱}^{(h)}(\sqrt{Z_h}\psi ^{(h)})},\hfill \end{array}$$ $`(2.82)`$ where $`P_{\stackrel{~}{Z}_{h1},\sigma _{h1},C_h}(d\psi ^{(h)})`$ is obtained from $`P_{Z_h,\sigma _h,C_h}(d\psi ^{(h)})`$ by substituting $`Z_h`$ with $$\stackrel{~}{Z}_{h1}(𝐤^{})=Z_h[1+C_h^1(𝐤^{})z_h]$$ $`(2.83)`$ and $`\sigma _h(𝐤^{})`$ with $$\sigma _{h1}(𝐤^{})=\frac{Z_h}{\stackrel{~}{Z}_{h1}(𝐤^{})}[\sigma _h(𝐤^{})+C_h^1(𝐤^{})s_h];$$ $`(2.84)`$ moreover $$\stackrel{~}{𝒱}^{(h)}(\sqrt{Z_h}\psi ^{(h)})=𝒱^{(h)}(\sqrt{Z_h}\psi ^{(h)})s_hZ_hF_\sigma ^{(h)}z_hZ_h[F_\zeta ^{(h)}+v_0^{}F_\alpha ^{(h)}]$$ $`(2.85)`$ and the factor $`\mathrm{exp}(L\beta t_h)`$ in (2.82) takes into account the different normalization of the two measures, so that $$t_h=\frac{1}{L\beta }\underset{𝐤^{}:C_h^1(𝐤^{})>0}{}\mathrm{log}\left\{[1+z_hC_h^1(𝐤^{})]^2\frac{k_0^2+E(k^{})^2+\sigma _{h1}(𝐤^{})^2}{k_0^2+E(k^{})^2+\sigma _h(𝐤^{})^2}\right\}.$$ $`(2.86)`$ Note that $$\stackrel{~}{𝒱}^{(h)}(\psi )=\gamma ^hn_hF_\nu ^{(h)}+(a_hz_hv_0^{})F_\alpha ^{(h)}+l_hF_\lambda ^{(h)}.$$ $`(2.87)`$ The r.h.s of (2.82) can be written as $$e^{L\beta t_h}P_{Z_{h1},\sigma _{h1},C_{h1}}(d\psi ^{(h1)})P_{Z_{h1},\sigma _{h1},\stackrel{~}{f}_h^1}(d\psi ^{(h)})e^{\stackrel{~}{𝒱}^{(h)}(\sqrt{Z_h}\psi ^{(h)})},$$ $`(2.88)`$ where $$Z_{h1}=Z_h(1+z_h),\stackrel{~}{f}_h(𝐤^{})=Z_{h1}[\frac{C_h^1(𝐤^{})}{\stackrel{~}{Z}_{h1}(𝐤^{})}\frac{C_{h1}^1(𝐤^{})}{Z_{h1}}].$$ $`(2.89)`$ Note that $`\stackrel{~}{f}_h(𝐤^{})`$ has the same support of $`f_h(𝐤^{})`$; in fact, by using (2.38), it is easy to see that $$\stackrel{~}{f}_h(𝐤^{})=f_h(𝐤^{})\left[1+\frac{z_hf_{h+1}(𝐤^{})}{1+z_hf_h(𝐤^{})}\right].$$ $`(2.90)`$ Moreover, by (2.49), $$P_{Z_{h1},\sigma _{h1},\stackrel{~}{f}_h^1}(d\psi ^{(h)})\psi _{𝐱,\omega }^{(h)}\psi _{𝐲,\omega ^{}}^{(h)+}=\frac{g_{\omega ,\omega ^{}}^{(h)}(𝐱𝐲)}{Z_{h1}},$$ $`(2.91)`$ where $$g_{\omega ,\omega ^{}}^{(h)}(𝐱𝐲)=\frac{1}{L\beta }\underset{𝐤^{}}{}e^{i𝐤^{}(𝐱𝐲)}\stackrel{~}{f}_h(𝐤^{})[T_h^1(𝐤^{})]_{\omega ,\omega ^{}},$$ $`(2.92)`$ and $`T_h^1(𝐤^{})`$ is the inverse of the $`T_h(𝐤^{})`$ defined in (2.69). $`T_h^1(𝐤^{})`$ is well defined on the support of $`\stackrel{~}{f}_h(𝐤^{})`$ and, if we set $$A_h(𝐤^{})=detT_h(𝐤^{})=k_0^2E(k^{})^2[\sigma _{h1}(𝐤^{})]^2,$$ $`(2.93)`$ then $$T_h^1(𝐤^{})=\frac{1}{A_h(𝐤^{})}\left(\begin{array}{cc}ik_0E(k^{})& i\sigma _{h1}(𝐤^{})\\ i\sigma _{h1}(𝐤^{})& ik_0+E(k^{})\end{array}\right).$$ $`(2.94)`$ The propagator $`g_{\omega ,\omega ^{}}^{(h)}(𝐱)`$ is an antiperiodic function of $`x`$ and $`x_0`$, of period $`L`$ and $`\beta `$, respectively. Its large distance behaviour is given by the following lemma (see also \[BM2\]), where we use the definitions $$\sigma _h\sigma _h(0),$$ $`(2.95)`$ $$d_L(x)=\frac{L}{\pi }\mathrm{sin}(\frac{\pi x}{L}),d_\beta (x_0)=\frac{\beta }{\pi }\mathrm{sin}(\frac{\pi x_0}{\beta }),$$ $`(2.96)`$ $$𝐝(𝐱𝐲)=(d_L(xy),d_\beta (x_0y_0)).$$ $`(2.97)`$ 2.6 Lemma. Let us suppose that $`h_{L,\beta }h0`$ and $$\|z_h\|\frac{1}{2},\|s_h\|\frac{1}{2}\|\sigma _h\|,\|\delta ^{}\|\frac{1}{2}.$$ $`(2.98)`$ We can write $$g_{\omega ,\omega }^{(h)}(𝐱𝐲)=g_{L,\omega }^{(h)}(𝐱𝐲)+r_1^{(h)}(𝐱𝐲)+r_2^{(h)}(𝐱𝐲),$$ $`(2.99)`$ where $$g_{L,\omega }^{(h)}(𝐱𝐲)=\frac{1}{L\beta }\underset{𝐤^{}}{}\frac{e^{i𝐤^{}(𝐱𝐲)}}{ik_0+\omega v_0^{}k^{}}\stackrel{~}{f}_h(𝐤^{}).$$ $`(2.100)`$ Moreover, given the positive integers $`N,n_0,n_1`$ and putting $`n=n_0+n_1`$, there exist a constant $`C_{N,n}`$ such that $$\begin{array}{cc}\hfill \|_{x_0}^{n_0}\overline{}_x^{n_1}r_1^{(h)}(𝐱𝐲)\|& C_{N,n}\frac{\gamma ^{2h+n}}{1+(\gamma ^h\|𝐝(𝐱𝐲))\|^N},\hfill \\ \hfill \|_{x_0}^{n_0}\overline{}_x^{n_1}r_2^{(h)}(𝐱𝐲)\|& C_{N,n}\|\frac{\sigma ^h}{\gamma ^h}\|^2\frac{\gamma ^{h+n}}{1+(\gamma ^h\|𝐝(𝐱𝐲)\|)^N},\hfill \end{array}$$ $`(2.101)`$ $$\|_{x_0}^{n_0}\overline{}_x^{n_1}g_{\omega ,\omega }^{(h)}(𝐱𝐲)\|C_{N,n}\|\frac{\sigma ^h}{\gamma ^h}\|\frac{\gamma ^{h+n}}{1+(\gamma ^h\|𝐝(𝐱𝐲)\|)^N}.$$ $`(2.102)`$ where $`\overline{}_x`$ denotes the discrete derivative. Note that $`g_{L,\omega }^{(h)}(𝐱𝐲)`$ coincides, in the limit $`\beta \mathrm{}`$, with the propagator “at scale $`\gamma ^h`$” of the Luttinger model, see \[BGM\], with $`\stackrel{~}{f}_h`$ in place of $`f_h`$. This remark will be crucial for studying the renormalization group flow in \[BeM\]. 2.7 Proof of Lemma 2.6. By using (2.38), it is easy to see that $`\sigma _h(𝐤^{})=\sigma _h(0)`$ on the support of $`f_h(𝐤^{})`$; hence, by (2.83) and (2.84), we have $$\sigma _{h1}(𝐤^{})=\frac{\sigma _h+C_h(𝐤^{})^1s_h}{1+z_hC_h(𝐤^{})^1},$$ $`(2.103)`$ implying, together with (2.98), that there exist two constants $`c_1,c_2`$ such that: $$c_1\|\sigma _h\|\|\sigma _{h1}(𝐤^{})\|c_2\|\sigma _h\|.$$ $`(2.104)`$ Let us now consider two integers $`N_0,N_10`$, such that $`N=N_0+N_1`$, and note that $$\begin{array}{cc}& d_L(xy)^{N_1}d_\beta (x_0y_0)^{N_0}g_{\omega ,\omega ^{}}^{(h)}(𝐱𝐲)=\hfill \\ & e^{i\pi (xL^1N_1+x_0\beta ^1N_0)}(i)^{N_0+N_1}\frac{1}{L\beta }\underset{𝐤^{}}{}e^{i𝐤^{}(𝐱𝐲)}_k^{}^{N_1}_{k_0}^{N_0}\left[\stackrel{~}{f}_h(𝐤^{})[T_h^1(𝐤^{})]_{\omega ,\omega ^{}}\right],\hfill \end{array}$$ $`(2.105)`$ where $`_k^{}`$ and $`_{k_0}`$ denote the discrete derivatives. If $`\omega =\omega ^{}`$, the decomposition (2.99) is related to the following identity: $$\begin{array}{cc}\hfill [T_h^1(𝐤^{})]_{\omega ,\omega }& =\frac{1}{ik_0+\omega v_0^{}k^{}}+\left[\frac{1}{ik_0+\omega E(k^{})}\frac{1}{ik_0+\omega v_0^{}k^{}}\right]+\hfill \\ & +\left[\frac{ik_0+\omega E(k^{})}{k_0^2+E(k^{})^2+[\sigma _{h1}(𝐤^{})]^2}\frac{1}{ik_0+\omega E(k^{})}\right].\hfill \end{array}$$ $`(2.106)`$ The bounds (2.101) and (2.102) easily follow from (2.98), (2.104), the support properties of $`f_h(𝐤^{})`$ and the observation that $`\stackrel{~}{f}_h(𝐤^{})`$ and $`\sigma _h(𝐤^{})`$ are smooth functions of $`𝐤^{}`$ in $`R^2`$, in the support of $`f_h(𝐤^{})`$, so that the discrete derivatives can be bounded as the continuous derivatives. The main point is of course the fact that, on the support of $`f_h(𝐤^{})`$, $`\|ik_0+\omega E(k^{})\|`$, $`\|ik_0+\omega v_0^{}k^{}\|`$ and $`\sqrt{k_0^2+E(k^{})^2+[\sigma _{h1}(𝐤^{})]^2}`$ are of order $`\gamma ^h`$. 2.8 We now rescale the field so that $$\stackrel{~}{𝒱}^{(h)}(\sqrt{Z_h}\psi ^{(h)})=\widehat{𝒱}^{(h)}(\sqrt{Z_{h1}}\psi ^{(h)});$$ $`(2.107)`$ it follows that $$\widehat{𝒱}^{(h)}(\psi )=\gamma ^h\nu _hF_\nu ^{(h)}+\delta _hF_\alpha ^{(h)}+\lambda _hF_\lambda ^{(h)},$$ $`(2.108)`$ where $$\nu _h=\frac{Z_h}{Z_{h1}}n_h,\delta _h=\frac{Z_h}{Z_{h1}}(a_hv_0^{}z_h),\lambda _h=(\frac{Z_h}{Z_{h1}})^2l_h.$$ $`(2.109)`$ We call the set $`\stackrel{}{v}_h=(\nu _h,\delta _h,\lambda _h)`$ the running coupling constants. If now define $$e^{𝒱^{(h1)}(\sqrt{Z_{h1}}\psi ^{(h1)})L\beta \stackrel{~}{E}_h}=P_{Z_{h1},\sigma _{h1},\stackrel{~}{f}_h^1}(d\psi ^{(h)})e^{\widehat{𝒱}^{(h)}(\sqrt{Z_{h1}}\psi ^{(h)})},$$ $`(2.110)`$ it is easy to see that $`𝒱^{(h1)}(\sqrt{Z_{h1}}\psi ^{(h1)})`$ is of the form (2.70) and that $$E_{h1}=E_h+t_h+\stackrel{~}{E}_h.$$ $`(2.111)`$ It is sufficient to use the well known identity $$L\beta \stackrel{~}{E}_h+𝒱^{(h1)}(\sqrt{Z_{h1}}\psi ^{(h1)})=\underset{n=1}{\overset{\mathrm{}}{}}\frac{1}{n!}(1)^{n+1}_h^{T,n}(\widehat{𝒱}^{(h)}(\sqrt{Z_{h1}}\psi ^{(h)})),$$ $`(2.112)`$ where $`_h^{T,n}`$ denotes the truncated expectation of order $`n`$ with propagator $`Z_{h1}^1g_{\omega ,\omega ^{}}^{(h)}`$, see (2.91), and observe that $`\psi ^{(h)}=\psi ^{(h1)}+\psi ^{(h)}`$. Moreover, the symmetry relations (2.71) are still satisfied, because the symmetry properties of the free measure are not modified by the renormalization procedure, so that the effective potential on scale $`h`$ has the same symmetries as the effective potential on scale $`0`$. Let us now define $`\stackrel{~}{E}_{h_{L,\beta }}`$, so that $$e^{L\beta \stackrel{~}{E}_{h_{L,\beta }}}=P_{Z_{h_{L,\beta }1},\sigma _{h_{L,\beta }1},\stackrel{~}{f}_{h_{L,\beta }}^1}(d\psi ^{(h_{L,\beta })})e^{\widehat{𝒱}^{(h_{L,\beta })}(\sqrt{Z_{h_{L,\beta }1}}\psi ^{(h_{L,\beta })})}.$$ $`(2.113)`$ We have $$E_{L,\beta }=\underset{h=h_{L,\beta }}{\overset{1}{}}[\stackrel{~}{E}_h+t_h].$$ $`(2.114)`$ Note that the above procedure allows us to write the running coupling constants $`\stackrel{}{v}_h`$, $`h0`$, in terms of $`\stackrel{}{v}_h^{}`$, $`0h^{}h+1`$, and $`\lambda ,\nu ,u`$: $$\stackrel{}{v}_h=\stackrel{}{\beta }(\stackrel{}{v}_{h+1},\mathrm{},\stackrel{}{v}_0,\lambda ,\nu ,u,\delta ^{}).$$ $`(2.115)`$ The function $`\stackrel{}{\beta }(\stackrel{}{v}_{h+1},\mathrm{},\stackrel{}{v}_0,\lambda ,\nu ,u,\delta ^{})`$ is called the Beta function. 2.9 Let us now explain the main motivations of the integration procedure discussed above. In a renormalization group approach one has to identify the relevant, marginal and irrelevant interactions. By a power counting argument one sees that the terms bilinear in the fields are relevant, hence one should extract from them the relevant and marginal local contributions by a Taylor expansion of the kernel up to order $`1`$ in the external momenta. Since $`\sigma _1+\sigma _2=0`$ by the remark following (2.62), we have to consider only two kinds of bilinear terms: those with $`\omega _1=\omega _2`$ and those with $`\omega _1=\omega _2`$. It turns out that, for the bilinear terms with $`\omega _1=\omega _2`$, a Taylor expansion up to order $`0`$ is sufficient; the reason is that the Feynman graphs contributing to such terms contain at least one non diagonal propagator and, by lemma 2.6, such propagators are smaller than the diagonal ones by a factor $`\sigma _h\gamma ^h`$; as we shall see, this is sufficient to improve the power counting by $`1`$. The previous discussion implies that the regularization of the bilinear terms produces four local terms. One of them, that proportional to $`F_\nu `$, is relevant; it reflects the renormalization of the Fermi momentum and is faced in a standard way \[BG\], by fixing properly the counterterm $`\nu `$ in the Hamiltonian, i.e. by fixing properly the chemical potential, so that the corresponding running coupling $`\nu _h`$ goes to $`0`$ for $`h\mathrm{}`$. The term proportional to $`F_\alpha `$ is marginal, but, as we shall see, stays bounded and of order $`\lambda `$ as $`h\mathrm{}`$, if $`\delta ^{}`$ is of order $`\lambda `$; hence the convergence of the flow is not related to the exact value of $`\delta ^{}`$. However, in order to get a detailed description of the spin correlation function asymptotic behaviour, it is convenient to choose $`\delta ^{}`$ so that $`\delta _h0`$ as $`h\mathrm{}`$. This choice implies that $`v_0^{}=v_0(1+\delta ^{})`$ is the “effective” Fermi velocity of the fermion system. The other two terms are marginal, but have to be treated in different ways. The term proportional to $`F_\zeta `$ is absorbed in the free measure and produces a field renormalization, as in the Luttinger liquid (which is indeed obtained for $`u=0`$). The term proportional to $`F_\sigma `$, related to the presence of a gap in the spectrum, is also absorbed in the free measure, since there is no free parameter in the Hamiltonian to control its flow, as for $`F_\zeta `$. This operation can be seen as the application of a sequence of different Bogoliubov transformations at each integration step, to compare with the single Bogoliubov transformation that it is sufficient to see a gap $`O(u)`$ at the Fermi surface, in the $`XY`$ model ($`\lambda =0`$). It turns out that the gap is deeply renormalized by the interaction, since $`\sigma _h`$ is a sort of “mass terms” with a non trivial renormalization group flow. Let us now consider the quartic terms, which are all marginal. Since there are many of them, depending on the labels $`\omega _i`$ and $`\sigma _i`$ of each field, their renormalization group flow seems difficult to study. However, as we have explained in §2.5, the running couplings corresponding to the quartic terms are all exactly equal to $`0`$ for trivial reasons, unless, after a suitable permutation of the fields, $`\underset{¯}{\sigma }=(+,,+,)`$, $`\underset{¯}{\omega }=(+1,1,1,+1)`$. Hence, by a Taylor expansion of the kernel up to order $`0`$ in the external momenta, all quartic terms can be regularized, by introducing only one running coupling, $`\lambda _h`$. As in the Luttinger liquid \[BGPS, BM1\], the flow of $`\lambda _h`$ and $`\delta _h`$ can be controlled by using some cancellations, due to the fact that the Beta function is “close” (for small $`u`$) to the Luttinger model Beta function. In lemma 2.6 we write the propagator as the Luttinger model propagator plus a remainder, so that the Beta function is equal to the Luttinger model Beta function plus a “remainder”, which is small if $`\sigma _h\gamma ^h`$ is small. Let us define $$h^{}=\mathrm{inf}\{h:0hh_{L,\beta },a_0v_0^{}\gamma ^{\overline{h}1}4\|\sigma _{\overline{h}}\|,\overline{h}:0\overline{h}h\}.$$ $`(2.116)`$ Of course this definition is meaningful only if $`a_0v_0^{}\gamma ^14\|\sigma _0\|=4\|u\|v_0`$ (see (2.64)), that is if $$\|u\|\frac{a_0}{4\gamma }(1+\delta ^{}).$$ $`(2.117)`$ If the condition (2.117) is not satisfied, we shall put $`h^{}=1`$. Lemma 2.6, (2.86) and the definition of $`h^{}`$ easily imply this other Lemma. 2.10 Lemma. If $`h>h^{}0`$ and the conditions (2.98) are satisfied, there is a constant $`C`$ such that $$\|t_h\|C\gamma ^{2h}.$$ $`(2.118)`$ Moreover, given the positive integers $`N,n_0,n_1`$ and putting $`n=n_0+n_1`$, there exist a constant $`C_{N,n}`$ such that $$\|_{x_0}^{n_0}\overline{}_x^{n_1}g_{\omega ,\omega ^{}}^{(h)}(𝐱;𝐲)\|C_{N,n}\frac{\gamma ^{h+n}}{1+(\gamma ^h\|𝐝(𝐱𝐲)\|)^N}.$$ $`(2.119)`$ 2.11 In §$`\mathrm{}`$3 we will see that, using the above lemmas and assuming that the running coupling constants are bounded, the integration of the field $`\psi ^{(h)}`$ in (2.88) is well defined in the limit $`L,\beta \mathrm{}`$, for $`0h>h^{}`$. The integration of the scales from $`h^{}`$ to $`h_{L,\beta }`$ will be performed “in a single step”. This is possible because we shall prove in §$`\mathrm{}`$3 that the integration in the r.h.s. in (2.82) is well defined in the limit $`L,\beta \mathrm{}`$, for $`h=h^{}`$. In order to do that, we shall use the following lemma, whose proof is similar to the proof of lemma 2.6. 2.12 Lemma. Assume that $`h^{}`$ is finite uniformly in $`L,\beta `$, so that $`\|\sigma _{h^{}1}\gamma ^h^{}\|\overline{\kappa }`$, for a suitable constant $`\overline{\kappa }`$ and define $$\frac{\overline{g}_{\omega ,\omega ^{}}^{(h^{})}(𝐱𝐲)}{Z_{h^{}1}}P_{\stackrel{~}{Z}_{h^{}1},\sigma _{h^{}1},C_h^{}}(d\psi ^{(h^{})})\psi _{𝐱,\omega }^{(h^{})}\psi _{𝐲,\omega ^{}}^{(h^{})+}.$$ $`(2.120)`$ Then, given the positive integers $`N,n_0,n_1`$ and putting $`n=n_0+n_1`$, there exist a constant $`C_{N,n}`$ such that $$\|_{x_0}^{n_0}\overline{}_x^{n_1}g_{\omega ,\omega ^{}}^{(h^{})}(𝐱;𝐲)\|C_{N,n}\frac{\gamma ^{h^{}+n}}{1+(\gamma ^h^{}\|𝐝(𝐱𝐲)\|)^N}.$$ $`(2.121)`$ 2.13 Comparing Lemma 2.10 and Lemma 2.12, we see that the propagator of the integration of all the scales between $`h^{}`$ and $`h_{L,\beta }`$ has the same bound as the propagator of the integration of a single scale greater than $`h^{}`$; this property is used to perform the integration of all the scales $`h^{}`$ in a single step. In fact $`\gamma ^h^{}`$ is a momentum scale and, roughly speaking, for momenta bigger than $`\gamma ^h^{}`$ the theory is “essentially” a massless theory (up to $`O(\sigma _h\gamma ^h)`$ terms), while for momenta smaller than $`\gamma ^h^{}`$ it is a “massive” theory with mass $`O(\gamma ^h^{})`$. 3. Analyticity of the effective potential 3.1 We want to study the expansion of the effective potential, which follows from the renormalization procedure discussed in §2. In order to do that, we find it convenient to write $`𝒱^{(h)}`$, $`h1`$, in terms of the variables $`\psi _{𝐱,\omega }^{(h)\sigma }`$. The two contributions to $`𝒱^{(1)}(\psi ^{(1)})`$, see (2.58) and (2.14), become $$\begin{array}{cc}\hfill \lambda V_\lambda (\psi ^1)& =\underset{\underset{¯}{\sigma }}{}d𝐱d𝐲\lambda v_\lambda (𝐱𝐲)e^{i𝐩_F𝐱(\sigma _1+\sigma _4)+i𝐩_F𝐲(\sigma _2+\sigma _3)}\hfill \\ & \psi _{𝐱,\sigma _1}^{(1)\sigma _1}\psi _{𝐲,\sigma _2}^{(1)\sigma _2}\psi _{𝐲,\sigma _3}^{(1)\sigma _3}\psi _{𝐱,\sigma _4}^{(1)\sigma _4},\hfill \\ \hfill \nu N(\psi ^1)& =\underset{\sigma _1,\sigma _2}{}𝑑𝐱e^{i𝐩_F𝐱(\sigma _1+\sigma _2)}\nu \psi _{𝐱,\sigma _1}^{(1)\sigma _1}\psi _{𝐱,\sigma _2}^{(1)\sigma _2},\hfill \end{array}$$ $`(3.1)`$ where $`𝑑𝐱`$ is a shorthand for $`_{x\mathrm{\Lambda }}_{\beta /2}^{\beta /2}𝑑x_0`$. If we define $$\begin{array}{cc}& W_{2n,\underset{¯}{\sigma },\underset{¯}{\omega }}^{(h)}(𝐱_1,\mathrm{},𝐱_{2n})=\hfill \\ & =\frac{1}{(L\beta )^{2n}}\underset{𝐤_1^{},\mathrm{},𝐤_{2n}^{}}{}e^{i_{r=1}^{2n}\sigma _r𝐤_r^{}𝐱_r}\widehat{W}_{2n,\underset{¯}{\sigma },\underset{¯}{\omega }}^{(h)}(𝐤_1^{},\mathrm{},𝐤_{2n1}^{})\delta (\underset{i=1}{\overset{2n}{}}\sigma _i(𝐤_i^{}+𝐩_F)),\hfill \end{array}$$ $`(3.2)`$ we can write (2.70) as $$𝒱^{(h)}(\psi ^{(h)})=\underset{n=1}{\overset{\mathrm{}}{}}\underset{\underset{¯}{\sigma },\underset{¯}{\omega }}{}𝑑𝐱_1\mathrm{}𝑑𝐱_{2n}\left[\underset{i=1}{\overset{2n}{}}\psi _{𝐱_i,\omega _i}^{(h)\sigma _i}\right]W_{2n,\underset{¯}{\sigma },\underset{¯}{\omega }}^{(h)}(𝐱_1,\mathrm{},𝐱_{2n}).$$ $`(3.3)`$ Note that $$W_{2n,\underset{¯}{\sigma },\underset{¯}{\omega }}^{(h)}(𝐱_1+𝐱,\mathrm{},𝐱_{2n}+𝐱)=e^{i𝐩_F𝐱_{r=1}^{2n}\sigma _r}W_{2n,\underset{¯}{\sigma },\underset{¯}{\omega }}^{(h)}(𝐱_1,\mathrm{},𝐱_{2n}),$$ $`(3.4)`$ hence $`W_{2n,\underset{¯}{\sigma },\underset{¯}{\omega }}^{(h)}(𝐱_1,\mathrm{},𝐱_{2n})`$ is translation invariant if and only if $`_{r=1}^{2n}\sigma _r=0`$. The representation of $`𝒱^{(h)}(\psi ^{(h)})`$ in terms of the $`\psi _{𝐱,\omega }^{(h)\sigma }`$ variables is obtained by substituting in the r.h.s. of (2.79) the $`𝐱`$-space representations of the definitions (2.80). We have $$\begin{array}{ccc}\hfill F_\nu ^{(h)}& =\underset{\omega =\pm 1}{}\omega 𝑑𝐱\psi _{𝐱,\omega }^{(h)+}\psi _{𝐱,\omega }^{(h)},\hfill & \\ \hfill F_\sigma ^{(h)}& =\underset{\omega =\pm 1}{}i\omega 𝑑𝐱\psi _{𝐱,\omega }^{(h)+}\psi _{𝐱,\omega }^{(h)},\hfill & \\ \hfill F_\alpha ^{(h)}& =\underset{\omega =\pm 1}{}i\omega 𝑑𝐱\psi _{𝐱,\omega }^{(h)+}[\overline{}_1\psi _{𝐱,\omega }^{(h)}+\frac{i\mathrm{cos}p_F}{2v_0}\overline{}_1^2\psi _{𝐱,\omega }^{(h)}]=\hfill & (3.5)\hfill \\ & =\underset{\omega =\pm 1}{}i\omega 𝑑𝐱[\overline{}_1\psi _{𝐱,\omega }^{(h)+}+\frac{i\mathrm{cos}p_F}{2v_0}\overline{}_1^2\psi _{𝐱,\omega }^{(h)+}]\psi _{𝐱,\omega }^{(h)},\hfill & \\ \hfill F_\zeta ^{(h)}& =\underset{\omega =\pm 1}{}𝑑𝐱\psi _{𝐱,\omega }^{(h)+}_0\psi _{𝐱,\omega }^{(h)}=\underset{\omega =\pm 1}{}𝑑𝐱_0\psi _{𝐱,\omega }^{(h)+}\psi _{𝐱,\omega }^{(h)},\hfill & \\ \hfill F_\lambda ^{(h)}& =𝑑𝐱\psi _{𝐱,+1}^{(h)+}\psi _{𝐱,1}^{(h)}\psi _{𝐱,1}^{(h)+}\psi _{𝐱,+1}^{(h)},\hfill & \end{array}$$ where $`_0`$ is the derivative w.r.t. $`x_0`$, $`\overline{}_1`$ is the symmetric discrete derivative w.r.t. $`x`$, that is, given a function $`f(𝐱)`$, $$\overline{}_1f(𝐱)=[f(x+1,x_0)f(x1,x_0)]/2,$$ $`(3.6)`$ and $`\overline{}_1^2`$ (which is not the square of $`\overline{}_1`$, but has the same properties) is defined by the equation $$\overline{}_1^2f(𝐱)=f(x+1,x_0)+f(x1,x_0)2f(x,x_0).$$ $`(3.7)`$ Let us now discuss the action of the operator $``$ and $`=1`$ on the effective potential in the $`x`$-space representation, by considering the terms for which $`0`$. 1)If $`2n=4`$, by (2.72), $$𝑑\underset{¯}{𝐱}W(\underset{¯}{𝐱})\underset{i=1}{\overset{4}{}}\psi _{𝐱_i,\omega _i}^{(h)\sigma _i}=𝑑\underset{¯}{𝐱}W(\underset{¯}{𝐱})\underset{i=1}{\overset{4}{}}\left[G_{\sigma _i}(𝐱_i𝐱_4)\psi _{𝐱_4,\omega _i}^{(h)\sigma _i}\right],$$ $`(3.8)`$ where $`\underset{¯}{𝐱}=(𝐱_1,\mathrm{},𝐱_4)`$, $`W(\underset{¯}{𝐱})=W_{4,\underset{¯}{\sigma },\underset{¯}{\omega }}^{(h)}(𝐱_1,𝐱_2,𝐱_3,𝐱_4)`$ and $$G_\sigma (𝐱)=e^{i\sigma \overline{𝐤}_{++}𝐱}=e^{i\sigma \pi (\frac{x}{L}+\frac{x_0}{\beta })}.$$ $`(3.9)`$ Note that, as we have discussed in §2.5, the r.h.s. of (3.8) is always equal to $`0`$, unless, after a suitable permutation of the fields, $`\underset{¯}{\sigma }=(+,,+,)`$, $`\underset{¯}{\omega }=(+1,1,1,+1)`$. In this last case the function $`W(\underset{¯}{𝐱})_{i=1}^4G_{\sigma _i}(𝐱_i𝐱_4)=W(\underset{¯}{𝐱})G_+(𝐱_1𝐱_2+𝐱_3𝐱_4)`$ is translation invariant and periodic in the space and time components of all variables $`𝐱_k`$, of period $`L`$ and $`\beta `$, respectively. It follows that the quantities $`G_{\sigma _i}(𝐱_i𝐱_4)\psi _{𝐱_4,\omega _i}^{(h)\sigma _i}`$ in the r.h.s. of (3.8) can be substituted with $`G_{\sigma _i}(𝐱_i𝐱_k)\psi _{𝐱_k,\omega _i}^{(h)\sigma _i}`$, $`k=1,2,3`$. Hence we have four equivalent representations of the localization operation, which differ by the choice of the localization point. The freedom in the choice of the localization point will be useful in the following. If the localization point is chosen as in (3.8), we have $$\begin{array}{cc}& 𝑑\underset{¯}{𝐱}W(\underset{¯}{𝐱})\underset{i=1}{\overset{4}{}}\psi _{𝐱_i,\omega _i}^{(h)\sigma _i}=\hfill \\ & =𝑑\underset{¯}{𝐱}W(\underset{¯}{𝐱})\left[\underset{i=1}{\overset{4}{}}\psi _{𝐱_i,\omega _i}^{(h)\sigma _i}\underset{i=1}{\overset{4}{}}G_{\sigma _i}(𝐱_i𝐱_4)\psi _{𝐱_4,\omega _i}^{(h)\sigma _i}\right].\hfill \end{array}$$ $`(3.10)`$ The term in square brackets in the above equation can be written as $$\begin{array}{cc}& \psi _{𝐱_1,\omega _1}^{(h)\sigma _1}\psi _{𝐱_2,\omega _2}^{(h)\sigma _2}D_{𝐱_3,𝐱_4,\omega _3}^{1,1(h)\sigma _3}\psi _{𝐱_4,\omega _4}^{(h)\sigma _4}+\hfill \\ & +G_{\sigma _3}(𝐱_3𝐱_4)\psi _{𝐱_1,\omega _1}^{(h)\sigma _1}D_{𝐱_2,𝐱_4,\omega _2}^{1,1(h)\sigma _2}\psi _{𝐱_4,\omega _3}^{(h)\sigma _3}\psi _{𝐱_4,\omega _4}^{(h)\sigma _4}+\hfill \\ & +G_{\sigma _3}(𝐱_3𝐱_4)G_{\sigma _2}(𝐱_2𝐱_4)D_{𝐱_1,𝐱_4,\omega _1}^{1,1(h)\sigma _1}\psi _{𝐱_4,\omega _2}^{(h)\sigma _2}\psi _{𝐱_4,\omega _3}^{(h)\sigma _3}\psi _{𝐱_4,\omega _4}^{(h)\sigma _4},\hfill \end{array}$$ $`(3.11)`$ where $$D_{𝐲,𝐱,\omega }^{1,1(h)\sigma }=\psi _{𝐲,\omega }^{(h)\sigma }G_\sigma (𝐲𝐱)\psi _{𝐱,\omega }^{(h)\sigma }.$$ $`(3.12)`$ Similar expressions can be written, if the localization point is chosen in a different way. Note that the decomposition (3.11) corresponds to the following identity: $$\begin{array}{cc}\hfill \widehat{W}_{\tau ,𝐏}^{(h)}(𝐤_1^{},𝐤_2^{},𝐤_3^{})& =\left[\widehat{W}_{\tau ,𝐏}^{(h)}(𝐤_1^{},𝐤_2^{},𝐤_3^{})\widehat{W}_{\tau ,𝐏}^{(h)}(𝐤_1^{},𝐤_2^{},\overline{𝐤}_{++})\right]+\hfill \\ & +\left[\widehat{W}_{\tau ,𝐏}^{(h)}(𝐤_1^{},𝐤_2^{},\overline{𝐤}_{++})\widehat{W}_{\tau ,𝐏}^{(h)}(𝐤_1^{},\overline{𝐤}_{++},\overline{𝐤}_{++})\right]+\hfill \\ & +\left[\widehat{W}_{\tau ,𝐏}^{(h)}(𝐤_1^{},\overline{𝐤}_{++},\overline{𝐤}_{++})\widehat{W}_{\tau ,𝐏}^{(h)}(\overline{𝐤}_{++},\overline{𝐤}_{++},\overline{𝐤}_{++})\right],\hfill \end{array}$$ $`(3.13)`$ and that the $`i`$-th term in the r.h.s. of (3.13) is equal to $`0`$ for $`𝐤_i^{}=\overline{𝐤}_{++}`$. The field $`D_{𝐲,𝐱,\omega }^{1,1(h)\sigma }`$ is antiperiodic in the space and time components of $`𝐱`$ and $`𝐲`$, of period $`L`$ and $`\beta `$, and is equal to $`0`$ if $`𝐱=𝐲`$ modulo $`(L,\beta )`$. This means that it is dimensionally equivalent to the product of $`d(𝐱,𝐲)`$ (see (2.97)) and the derivative of the field, so that the bound of its contraction with another field variable on a scale $`h^{}<h`$ will produce a “gain” $`\gamma ^{(hh^{})}`$ with respect to the contraction of $`\psi _{𝐲,\omega }^{(h)\sigma }`$. If we insert (3.11) in the r.h.s. of (3.10), we can decompose the l.h.s in the sum of three terms, which differ from the term which $``$ acts on mainly because one $`\psi ^{(h)}`$ field is substituted with a $`D^{1,1(h)}`$ field and some of the other $`\psi ^{(h)}`$ fields are “translated” in the localization point. All three terms share the property that the field whose $`𝐱`$ coordinate is equal to the localization point is not affected by the action of $``$. In our approach, the regularization effect of $``$ will be exploited trough the decomposition (3.11). However, for reasons that will become clear in the following, it is convenient to start the analysis by using another representation of the expression resulting from the insertion of (3.11) in (3.10). If $`\psi _{𝐱_i}\psi _{𝐱_i,\omega _i}^{(h)\sigma _i}`$, we can write, if the localization point is $`𝐱_4`$, $$\begin{array}{cc}& 𝑑\underset{¯}{𝐱}\underset{i=1}{\overset{4}{}}\psi _{𝐱_i}W(\underset{¯}{𝐱})=\hfill \\ & =𝑑\underset{¯}{𝐱}\underset{i=1}{\overset{4}{}}\psi _{𝐱_i}\left[W(\underset{¯}{𝐱})\delta (𝐱_3𝐱_4)𝑑𝐲_3W(𝐱_1,𝐱_2,𝐲_3,𝐱_4)G_{\sigma _3}(𝐲_3𝐱_4)\right]+\hfill \\ & +d\underset{¯}{𝐱}\underset{i=1}{\overset{4}{}}\psi _{𝐱_i}\delta (𝐱_3𝐱_4)d𝐲_3[W(𝐱_1,𝐱_2,𝐲_3,𝐱_4)G_{\sigma _3}(𝐲_3𝐱_4)\hfill \\ & \delta (𝐱_2𝐱_4)d𝐲_2W(𝐱_1,𝐲_2,𝐲_3,𝐱_4)G_{\sigma _3}(𝐲_3𝐱_4)G_{\sigma _2}(𝐲_2𝐱_4)]+\hfill \\ & +d\underset{¯}{𝐱}\underset{i=1}{\overset{4}{}}\psi _{𝐱_i}\delta (𝐱_2𝐱_4)\delta (𝐱_3𝐱_4)d𝐲_2d𝐲_3[W(𝐱_1,𝐲_2,𝐲_3,𝐱_4)G_{\sigma _3}(𝐲_3𝐱_4)\hfill \\ & G_{\sigma _2}(𝐲_2𝐱_4)\delta (𝐱_1𝐱_4)d𝐲_1W(𝐲_1,𝐲_2,𝐲_3,𝐱_4)\underset{i=1}{\overset{3}{}}G_{\sigma _i}(𝐲_i𝐱_4)],\hfill \end{array}$$ $`(3.14)`$ where $`\delta (𝐱)`$ is the antiperiodic delta function, that is $$\delta (𝐱)=\frac{1}{L\beta }\underset{𝐤^{}𝒟_{L,\beta }^{}}{}e^{i\sigma 𝐤^{}𝐱}.$$ $`(3.15)`$ Similar expressions are obtained, if the localization point is chosen in a different way. In the new representation, the action of $``$ is seen as the decomposition of the original term in the sum of three terms, which are still of the form (3.3), but with a different kernel, containing suitable delta functions. 2)If $`2n=2`$ and, possibly after a suitable permutation of the fields, $`\underset{¯}{\sigma }=(+,)`$, $`\omega _1=\omega _2=\omega `$, by (2.74), $$\begin{array}{cc}& 𝑑𝐱_1𝑑𝐱_2W_{2,\underset{¯}{\sigma },\underset{¯}{\omega }}^{(h)}(𝐱_1𝐱_2)\psi _{𝐱_1,\omega }^{(h)+}\psi _{𝐱_2,\omega }^{(h)}=\hfill \\ & =𝑑𝐱_1𝑑𝐱_2W_{2,\underset{¯}{\sigma },\underset{¯}{\omega }}^{(h)}(𝐱_1𝐱_2)\psi _{𝐱_1,\omega }^{(h)+}T_{𝐱_2,𝐱_1,\omega }^{1(h)}\hfill \\ & =𝑑𝐱_1𝑑𝐱_2W_{2,\underset{¯}{\sigma },\underset{¯}{\omega }}^{(h)}(𝐱_1𝐱_2)T_{𝐱_1,𝐱_2,\omega }^{1(h)+}\psi _{𝐱_2,\omega }^{(h)},\hfill \end{array}$$ $`(3.16)`$ with $$\begin{array}{cc}& T_{𝐲,𝐱,\omega }^{1(h)\sigma }=\psi _{𝐱,\omega }^{(h)\sigma }c_\beta (y_0x_0)[c_L(yx)+b_Ld_L(yx)]+\hfill \\ & +[\overline{}_1\psi _{𝐱,\omega }^{(h)\sigma }+\frac{i\mathrm{cos}p_F}{2v_0}\overline{}_1^2\psi _{𝐱,\omega }^{(h)\sigma }]c_\beta (y_0x_0)a_Ld_L(yx)+\hfill \\ & +_0\psi _{𝐱,\omega }^{(h)\sigma }d_\beta (y_0x_0)c_L(yx),\hfill \end{array}$$ $`(3.17)`$ where $`d_L(x)`$ and $`d_\beta (x_0)`$ are defined as in (2.96) and $$c_L(x)=\mathrm{cos}(\pi xL^1),c_\beta (x_0)=\mathrm{cos}(\pi x_0\beta ^1).$$ $`(3.18)`$ As in the item 1), we define the localization point as the $`𝐱`$ coordinate of the field which is left unchanged $``$. We are free to choose it equal to $`𝐱_1`$ or $`𝐱_2`$. This freedom affects also the action of $``$, which can be written as $$\begin{array}{cc}& 𝑑𝐱_1𝑑𝐱_2W_{2,\underset{¯}{\sigma },\underset{¯}{\omega }}^{(h)}(𝐱_1𝐱_2)\psi _{𝐱_1,\omega }^{(h)+}\psi _{𝐱_2,\omega }^{(h)}=\hfill \\ & =𝑑𝐱_1𝑑𝐱_2W_{2,\underset{¯}{\sigma },\underset{¯}{\omega }}^{(h)}(𝐱_1𝐱_2)\psi _{𝐱_1,\omega }^{(h)+}D_{𝐱_2,𝐱_1,\omega }^{2(h)}\hfill \\ & =𝑑𝐱_1𝑑𝐱_2W_{2,\underset{¯}{\sigma },\underset{¯}{\omega }}^{(h)}(𝐱_1𝐱_2)D_{𝐱_1,𝐱_2,\omega }^{2(h)+}\psi _{𝐱_2,\omega }^{(h)},\hfill \end{array}$$ $`(3.19)`$ with $$D_{𝐲,𝐱,\omega }^{2(h)\sigma }=\psi _{𝐲,\omega }^{(h)\sigma }T_{𝐲,𝐱,\omega }^{1(h)\sigma }.$$ $`(3.20)`$ Hence the effect of $``$ can be described as the replacement of a $`\psi ^{(h)\sigma }`$ field with a $`D^{2(h)\sigma }`$ field, with a gain in the bounds (see discussion in item 1) above) of a factor $`\gamma ^{2(hh^{})}`$. Also in this case, it is possible to write the regularized term in the form (3.3). We get $$\begin{array}{ccc}& d𝐱d𝐲W_{2,\underset{¯}{\sigma },\underset{¯}{\omega }}^{(h)}(𝐱𝐲)\psi _{𝐱,\omega }^{(h)+}\psi _{𝐲,\omega }^{(h)}=d𝐱d𝐲\psi _{𝐱,\omega }^{(h)+}\psi _{𝐲,\omega }^{(h)}\{W_{2,\underset{¯}{\sigma },\underset{¯}{\omega }}^{(h)}(𝐱𝐲)\hfill & \\ & \delta (𝐲𝐱)𝑑𝐳W_{2,\underset{¯}{\sigma },\underset{¯}{\omega }}^{(h)}(𝐱𝐳)c_\beta (z_0x_0)[c_L(zx)+b_Ld_L(zx)]\hfill & \\ & [\overline{}_1\delta (𝐲𝐱)+\frac{i\mathrm{cos}p_F}{2v_0}\overline{}_1^2\delta (𝐲𝐱)]𝑑𝐳W_{2,\underset{¯}{\sigma },\underset{¯}{\omega }}^{(h)}(𝐱𝐳)c_\beta (z_0x_0)a_Ld_L(zx)\hfill & \\ & +_0\delta (𝐲𝐱)d𝐳W_{2,\underset{¯}{\sigma },\underset{¯}{\omega }}^{(h)}(𝐱𝐳)d_\beta (z_0x_0)c_L(zx)\}.\hfill & (3.21)\hfill \end{array}$$ 3)If $`2n=2`$ and, possibly after a suitable permutation of the fields, $`\underset{¯}{\sigma }=(+,)`$, $`\omega _1=\omega _2=\omega `$, by (2.74), $$\begin{array}{cc}& 𝑑𝐱_1𝑑𝐱_2W_{2,\underset{¯}{\sigma },\underset{¯}{\omega }}^{(h)}(𝐱_1𝐱_2)\psi _{𝐱_1,\omega }^{(h)+}\psi _{𝐱_2,\omega }^{(h)}=\hfill \\ & =𝑑𝐱_1𝑑𝐱_2W_{2,\underset{¯}{\sigma },\underset{¯}{\omega }}^{(h)}(𝐱_1𝐱_2)\psi _{𝐱_1,\omega }^{(h)+}T_{𝐱_2,𝐱_1,\omega }^{0(h)}\hfill \\ & =𝑑𝐱_1𝑑𝐱_2W_{2,\underset{¯}{\sigma },\underset{¯}{\omega }}^{(h)}(𝐱_1𝐱_2)T_{𝐱_1,𝐱_2,\omega }^{0(h)+}\psi _{𝐱_2,\omega }^{(h)},\hfill \end{array}$$ $`(3.22)`$ where $$T_{𝐲,𝐱,\omega }^{0(h)\sigma }=c_\beta (y_0x_0)c_L(yx)\psi _{𝐱,\omega }^{(h)\sigma }.$$ $`(3.23)`$ Therefore $$\begin{array}{cc}& 𝑑𝐱_1𝑑𝐱_2W_{2,\underset{¯}{\sigma },\underset{¯}{\omega }}^{(h)}(𝐱_1𝐱_2)\psi _{𝐱_1,\omega }^{(h)+}\psi _{𝐱_2,\omega }^{(h)}=\hfill \\ & =𝑑𝐱_1𝑑𝐱_2W_{2,\underset{¯}{\sigma },\underset{¯}{\omega }}^{(h)}(𝐱_1𝐱_2)\psi _{𝐱_1,\omega }^{(h)+}D_{𝐱_2,𝐱_1,\omega }^{1,2(h)}\hfill \\ & =𝑑𝐱_1𝑑𝐱_2W_{2,\underset{¯}{\sigma },\underset{¯}{\omega }}^{(h)}(𝐱_1𝐱_2)D_{𝐱_1,𝐱_2,\omega }^{1,2(h)+}\psi _{𝐱_2,\omega }^{(h)},\hfill \end{array}$$ $`(3.24)`$ where $$D_{𝐲,𝐱,\omega }^{1,2(h)\sigma }=\psi _{𝐲,\omega }^{(h)\sigma }T_{𝐲,𝐱,\omega }^{0(h)\sigma }.$$ $`(3.25)`$ Hence the effect of $``$ can be described as the replacement of a $`\psi ^{(h)\sigma }`$ field with a $`D^{1,2(h)\sigma }`$ field, with a gain in the bounds (see discussion in item 1) above) of a factor $`\gamma ^{(hh^{})}`$. As before, we can also write $$\begin{array}{cc}& d𝐱d𝐲W_{2,\underset{¯}{\sigma },\underset{¯}{\omega }}^{(h)}(𝐱𝐲)\psi _{𝐱,\omega }^{(h)+}\psi _{𝐲,\omega }^{(h)}=d𝐱d𝐲\psi _{𝐱,\omega }^{(h)+}\psi _{𝐲,\omega }^{(h)}\hfill \\ & \left\{W_{2,\underset{¯}{\sigma },\underset{¯}{\omega }}^{(h)}(𝐱𝐲)\delta (𝐲𝐱)𝑑𝐳W_{2,\underset{¯}{\sigma },\underset{¯}{\omega }}^{(h)}(𝐱𝐳)c_\beta (z_0x_0)c_L(zx)\right\}.\hfill \end{array}$$ $`(3.26)`$ 3.2 By using iteratively the “single scale expansion” (2.112), starting from $`\widehat{𝒱}^{(1)}=𝒱^{(1)}`$, we can write the effective potential $`𝒱^{(h)}(\sqrt{Z_h}\psi ^{(h)})`$, for $`h0`$, in terms of a tree expansion, similar to that described, for example, in \[BGPS\]. $`r`$ $`v_0`$ $`v`$ $`h`$ $`h+1`$ $`h_v`$ $`0`$ $`+1`$ $`+2`$ Fig. 1 We need some definitions and notations. 1) Let us consider the family of all trees which can be constructed by joining a point $`r`$, the root, with an ordered set of $`n1`$ points, the endpoints of the unlabeled tree (see Fig. 1), so that $`r`$ is not a branching point. $`n`$ will be called the order of the unlabeled tree and the branching points will be called the non trivial vertices. The unlabeled trees are partially ordered from the root to the endpoints in the natural way; we shall use the symbol $`<`$ to denote the partial order. Two unlabeled trees are identified if they can be superposed by a suitable continuous deformation, so that the endpoints with the same index coincide. It is then easy to see that the number of unlabeled trees with $`n`$ end-points is bounded by $`4^n`$. We shall consider also the labeled trees (to be called simply trees in the following); they are defined by associating some labels with the unlabeled trees, as explained in the following items. 2) We associate a label $`h0`$ with the root and we denote $`𝒯_{h,n}`$ the corresponding set of labeled trees with $`n`$ endpoints. Moreover, we introduce a family of vertical lines, labeled by an an integer taking values in $`[h,2]`$, and we represent any tree $`\tau 𝒯_{h,n}`$ so that, if $`v`$ is an endpoint or a non trivial vertex, it is contained in a vertical line with index $`h_v>h`$, to be called the scale of $`v`$, while the root is on the line with index $`h`$. There is the constraint that, if $`v`$ is an endpoint, $`h_v>h+1`$. The tree will intersect in general the vertical lines in set of points different from the root, the endpoints and the non trivial vertices; these points will be called trivial vertices. The set of the vertices of $`\tau `$ will be the union of the endpoints, the trivial vertices and the non trivial vertices. Note that, if $`v_1`$ and $`v_2`$ are two vertices and $`v_1<v_2`$, then $`h_{v_1}<h_{v_2}`$. Moreover, there is only one vertex immediately following the root, which will be denoted $`v_0`$ and can not be an endpoint; its scale is $`h+1`$. Finally, if there is only one endpoint, its scale must be equal to $`+2`$ or $`h+2`$. 3) With each endpoint $`v`$ of scale $`h_v=+2`$ we associate one of the two contributions to $`𝒱^{(1)}(\psi ^{(1)})`$, written as in (3.1) and a set $`𝐱_v`$ of space-time points (the corresponding integration variables), two for $`\lambda V_\lambda (\psi ^{(1)})`$, one for $`\nu N(\psi ^{(1)})`$; we shall say that the endpoint is of type $`\lambda `$ or $`\nu `$, respectively. With each endpoint $`v`$ of scale $`h_v1`$ we associate one of the four local terms that we obtain if we write $`V^{(h_v1)}`$ (see (2.108)) by using the expressions (3.5) (there are four terms since $`F_\alpha `$ is the sum of two different local terms), and one space-time point $`𝐱_v`$; we shall say that the endpoint is of type $`\nu `$, $`\delta _1`$, $`\delta _2`$, $`\lambda `$, with an obvious correspondence with the different terms. Given a vertex $`v`$, which is not an endpoint, $`𝐱_v`$ will denote the family of all space-time points associated with one of the endpoints following $`v`$. Moreover, we impose the constraint that, if $`v`$ is an endpoint and $`𝐱_v`$ is a single space-time point (that is the corresponding term is local), $`h_v=h_v^{}+1`$, if $`v^{}`$ is the non trivial vertex immediately preceding $`v`$. 4) If $`v`$ is not an endpoint, the cluster $`L_v`$ with frequency $`h_v`$ is the set of endpoints following the vertex $`v`$; if $`v`$ is an endpoint, it is itself a (trivial) cluster. The tree provides an organization of endpoints into a hierarchy of clusters. 5) The trees containing only the root and an endpoint of scale $`h+1`$ will be called the trivial trees; note that they do not belong to $`𝒯_{h,1}`$, if $`h0`$, and can be associated with the four terms in the local part of $`\widehat{𝒱}^{(h)}`$. 6) We introduce a field label $`f`$ to distinguish the field variables appearing in the terms associated with the endpoints as in item 3); the set of field labels associated with the endpoint $`v`$ will be called $`I_v`$. Analogously, if $`v`$ is not an endpoint, we shall call $`I_v`$ the set of field labels associated with the endpoints following the vertex $`v`$; $`𝐱(f)`$, $`\sigma (f)`$ and $`\omega (f)`$ will denote the space-time point, the $`\sigma `$ index and the $`\omega `$ index, respectively, of the field variable with label $`f`$. If $`h_v0`$, one of the field variables belonging to $`I_v`$ carries also a discrete derivative $`\overline{}_1^m`$, $`m\{1,2\}`$, if the corresponding local term is of type $`\delta _m`$, see (3.5). Hence we can associate with each field label $`f`$ an integer $`m(f)\{0,1,2\}`$, denoting the order of the discrete derivative. Note that $`m(f)`$ is not uniquely determined, since we are free to use the first or the second representation of $`F_\alpha ^{(h_v1)}`$ in (3.5); we shall use this freedom in the following. By using (2.112), it is not hard to see that, if $`h0`$, the effective potential can be written in the following way: $$𝒱^{(h)}(\sqrt{Z_h}\psi ^{(h)})+L\beta \stackrel{~}{E}_{h+1}=\underset{n=1}{\overset{\mathrm{}}{}}\underset{\tau 𝒯_{h,n}}{}V^{(h)}(\tau ,\sqrt{Z_h}\psi ^{(h)})$$ $`(3.27),`$ where, if $`v_0`$ is the first vertex of $`\tau `$ and $`\tau _1,..,\tau _s`$ ($`s=s_{v_0}`$) are the subtrees of $`\tau `$ with root $`v_0`$, $`V^{(h)}(\tau ,\sqrt{Z_h}\psi ^{(h)})`$ is defined inductively by the relation $$\begin{array}{cc}& V^{(h)}(\tau ,\sqrt{Z_h}\psi ^{(h)})=\hfill \\ & \frac{(1)^{s+1}}{s!}_{h+1}^T[\overline{V}^{(h+1)}(\tau _1,\sqrt{Z_h}\psi ^{(h+1)});..;\overline{V}^{(h+1)}(\tau _s,\sqrt{Z_h}\psi ^{(h+1)})],\hfill \end{array}$$ $`(3.28)`$ and $`\overline{V}^{(h+1)}(\tau _i,\sqrt{Z_h}\psi ^{(h+1)})`$ a) is equal to $`\widehat{𝒱}^{(h+1)}(\tau _i,\sqrt{Z_h}\psi ^{(h+1)})`$ if the subtree $`\tau _i`$ is not trivial (see (2.107) for the definition of $`\widehat{𝒱}^{(h)}`$); b) if $`\tau _i`$ is trivial and $`h1`$, it is equal to one of the terms in the r.h.s. of (2.108) with scale $`h+1`$ or, if $`h=0`$, to one of the terms contributing to $`\widehat{𝒱}^{(1)}(\psi ^1)`$. If $`h=0`$, the r.h.s. of (3.28) can be written more explicitly in the following way. Given $`\tau 𝒯_{0,n}`$, there are $`n`$ endpoints of scale $`2`$ and only another one vertex, $`v_0`$, of scale $`1`$; let us call $`v_1,\mathrm{},v_n`$ the endpoints. We choose, in any set $`I_{v_i}`$, a subset $`Q_{v_i}`$ and we define $`P_{v_0}=_iQ_{v_i}`$; then we can write (recall that $`Z_0=1`$) $$V^{(0)}(\tau ,\sqrt{Z_0}\psi ^{(0)})=\underset{P_{v_0}}{}V^{(0)}(\tau ,P_{v_0}),$$ $`(3.29)`$ $$V^{(0)}(\tau ,P_{v_0})=\sqrt{Z_0}^{\|P_{v_0}\|}𝑑𝐱_{v_0}\stackrel{~}{\psi }^0(P_{v_0})K_{\tau ,P_{v_0}}^{(1)}(𝐱_{v_0}),$$ $`(3.30)`$ $$K_{\tau ,P_{v_0}}^{(1)}(𝐱_{v_0})=\frac{1}{n!}_1^T[\stackrel{~}{\psi }^{(1)}(P_{v_1}\backslash Q_{v_1}),\mathrm{},\stackrel{~}{\psi }^{(1)}(P_{v_n}\backslash Q_{v_n})]\underset{i=1}{\overset{n}{}}K_{v_i}^{(2)}(𝐱_{v_i}),$$ $`(3.31)`$ where we used the definitions $$\stackrel{~}{\psi }^{(h)}(P_v)=\underset{fP_v}{}\overline{}_1^{m(f)}\psi _{𝐱(f),\omega (f)}^{(h)\sigma (f)},$$ $`(3.32)`$ $$K_{v_i}^{(2)}(𝐱_{v_i})=e^{i𝐩_F_{fI_{v_i}}𝐱(f)\sigma (f)}\{\begin{array}{cc}\lambda v_\lambda (𝐱𝐲)\hfill & \text{if }v_i\text{ is of type }\lambda \text{ and }𝐱_{v_i}=(𝐱,𝐲)\text{,}\hfill \\ \nu \hfill & \text{if }v_i\text{ is of type }\nu \text{,}\hfill \end{array}$$ $`(3.33)`$ and we suppose that the order of the (anticommuting) field variables in (3.32) is suitable chosen in order to fix the sign as in (3.31). Note that the terms with $`P_{v_0}\mathrm{}`$ in the r.h.s. of (3.29) contribute to $`L\beta \stackrel{~}{E}_1`$, while the others contribute to $`𝒱^{(0)}(\sqrt{Z_0}\psi ^{(0)})`$. The potential $`\widehat{𝒱}^{(0)}(\sqrt{Z_1}\psi ^{(0)})`$, needed to iterate the previous procedure, is obtained, as explained in §2.5 and §2.8, by decomposing $`𝒱^{(0)}`$ in the sum of $`𝒱^{(0)}`$ and $`𝒱^{(0)}`$, by moving afterwards some local terms to the free measure and finally by rescaling the fields variables. The representation we get for $`𝒱^{(1)}(\sqrt{Z_1}\psi ^{(1)})`$ depends on the representation we use for $`V^{(0)}(\tau ,P_{v_0})`$. We choose to use that based on (3.14), (3.21) and (3.26), where the regularization is seen, for each term in the r.h.s. of (3.29) with $`P_{v_0}\mathrm{}`$, as a modification of the kernel $$W_{\tau ,P_{v_0}}^{(0)}(𝐱_{P_{v_0}})=d(𝐱_{v_0}\backslash 𝐱_{P_{v_0}})K_{\tau ,P_{v_0}}^{(1)}(𝐱_{v_0}),$$ $`(3.34)`$ where $`𝐱_{P_{v_0}}=_{fP_{v_0}}𝐱(f)`$. In order to remember this choice, we write $$V^{(0)}(\tau ,P_{v_0})=\sqrt{Z_0}^{\|P_{v_0}\|}𝑑𝐱_{v_0}\stackrel{~}{\psi }^{(0)}(P_{v_0})[K_{\tau ,P_{v_0}}^{(1)}(𝐱_{v_0})].$$ $`(3.35)`$ It is then easy to get, by iteration of the previous procedure, a simple expression for $`V^{(h)}(\tau ,\sqrt{Z_h}\psi ^{(h)})`$, for any $`\tau 𝒯_{h,n}`$. We associate with any vertex $`v`$ of the tree a subset $`P_v`$ of $`I_v`$, the external fields of $`v`$. These subsets must satisfy various constraints. First of all, if $`v`$ is not an endpoint and $`v_1,\mathrm{},v_{s_v}`$ are the vertices immediately following it, then $`P_v_iP_{v_i}`$; if $`v`$ is an endpoint, $`P_v=I_v`$. We shall denote $`Q_{v_i}`$ the intersection of $`P_v`$ and $`P_{v_i}`$; this definition implies that $`P_v=_iQ_{v_i}`$. The subsets $`P_{v_i}\backslash Q_{v_i}`$, whose union will be made, by definition, of the internal fields of $`v`$, have to be non empty, if $`s_v>1`$. Given $`\tau 𝒯_{h,n}`$, there are many possible choices of the subsets $`P_v`$, $`v\tau `$, compatible with all the constraints; we shall denote $`𝒫_\tau `$ the family of all these choices and $`𝐏`$ the elements of $`𝒫_\tau `$. Then we can write $$V^{(h)}(\tau ,\sqrt{Z_h}\psi ^{(h)})=\underset{𝐏𝒫_\tau }{}V^{(h)}(\tau ,𝐏);$$ $`(3.36)`$ $`V^{(h)}(\tau ,𝐏)`$ can be represented as in (3.30), that is as $$V^{(h)}(\tau ,𝐏)=\sqrt{Z_h}^{\|P_{v_0}\|}𝑑𝐱_{v_0}\stackrel{~}{\psi }^{(h)}(P_{v_0})K_{\tau ,𝐏}^{(h+1)}(𝐱_{v_0}),$$ $`(3.37)`$ with $`K_{\tau ,𝐏}^{(h+1)}(𝐱_{v_0})`$ defined inductively (recall that $`h_{v_0}=h+1`$) by the equation, valid for any $`v\tau `$ which is not an endpoint, $$\begin{array}{cc}\hfill K_{\tau ,𝐏}^{(h_v)}(𝐱_v)& =\frac{1}{s_v!}\left(\frac{Z_{h_v}}{Z_{h_v1}}\right)^{\frac{\|P_v\|}{2}}\underset{i=1}{\overset{s_v}{}}[K_{v_i}^{(h_v+1)}(𝐱_{v_i})]\hfill \\ & \stackrel{~}{}_{h_v}^T[\stackrel{~}{\psi }^{(h_v)}(P_{v_1}\backslash Q_{v_1}),\mathrm{},\stackrel{~}{\psi }^{(h_v)}(P_{v_{s_v}}\backslash Q_{v_{s_v}})],\hfill \end{array}$$ $`(3.38)`$ where $`\stackrel{~}{}_h^T`$ denotes the truncated expectation with propagator $`g^{(h)}`$ (without the scaling factor $`Z_{h1}`$, which is present in the definition of $`_h^T`$ used in (2.112)) and $`Z_11`$. Moreover, if $`v`$ is an endpoint and $`h_v=2`$, $`K_v^{(h_v)}(𝐱_v)`$ is defined by (3.33), otherwise $$K_v^{(h_v)}(𝐱_v)=\{\begin{array}{cc}\lambda _{h_v1}\hfill & \text{if }v\text{ is of type }\lambda \text{ ,}\hfill \\ i\omega \delta _{h_v1}\hfill & \text{if }v\text{ is of type }\delta _1\text{}d_2\text{ and }\omega (f)=\omega \text{ for both }fI_v\text{,}\hfill \\ \omega \gamma ^{h_v1}\nu _{h_v1}\hfill & \text{if }v\text{ is of type }\nu \text{ and }\omega (f)=\omega \text{ for both }fI_v\text{.}\hfill \end{array}$$ $`(3.39)`$ If $`v`$ is not an endpoint, $`K_v^{(h_v)}=K_{\tau _i,𝐏_i}^{(h_v)}`$, where $`\tau _1,\mathrm{},\tau _{s_v}`$ are the subtrees of $`\tau `$ with root $`v`$, $`𝐏_i=\{P_v,v\tau _i\}`$ and the action of $``$ is defined using the representation (3.14), (3.21) and (3.26) of the regularization operation, seen as a modification of the kernel $$W_{\tau ,𝐏}^{(h_v)}(𝐱_{P_v})=d(𝐱_v\backslash 𝐱_{P_v})K_{\tau ,𝐏}^{(h_v)}(𝐱_v),$$ $`(3.40)`$ where $`𝐱_{P_v}=_{fP_v}𝐱(f)`$. Finally we suppose again that the order of the (anticommuting) field variables is suitable chosen in order to fix the sign as in (3.37). Remark - The definitions (3.14), (3.21) and (3.26) of $``$ are sufficient, even if they are restricted to external fields with $`m(f)=0`$, because we can use the freedom in the definition of $`m(f)`$, see item 6) above, so that the external fields of $`v`$ have always $`m(f)=0`$, if $`v`$ is a vertex where the $``$ operation is acting on. This last claim follows from the observation that, since the truncated expectation in (3.38) vanishes if $`s_v>1`$ and $`P_{v_i}\backslash Q_{v_1}=\mathrm{}`$ for some $`i`$, at least one of the fields associated with the endpoints of type $`\delta _1`$ or $`\delta _2`$, the only ones which have fields with $`m(f)>0`$, has to be an internal field; hence, if one of the two fields is external, we can put $`m(f)=0`$ for it. If $`s_v=1`$ the previous argument should not work, but in this case the only vertex immediately following $`v`$ can be an endpoint of type $`\delta _1`$ or $`\delta _2`$ only if $`v=v_0`$, see item 2 above; however this is not a problem since the action of $``$ on a local term is equal to $`0`$. Note also that the kernel $`K_{\tau ,𝐏}^{(h_v)}(𝐱_v)`$ is translation invariant, if $`_{fP_v}\sigma (f)`$ $`=0`$; in general, it satisfies the relation $$K_{\tau ,𝐏}^{(h_v)}(𝐱_v+𝐱)=e^{i𝐩_F𝐱_{fP_v}\sigma (f)}K_{\tau ,𝐏}^{(h_v)}(𝐱_v).$$ $`(3.41)`$ There is a simple interpretation of $`V^{(h)}(\tau ,𝐏)`$ as the sum of a family $`𝒢_𝐏`$ of connected Feynman graphs build with single scale propagators of different scales, connecting the space-time points associated with the endpoints of the tree. A graph $`g𝒢_𝐏`$ is build by contracting, for any $`v\tau `$, all the internal fields in couples in all possible ways, by using the propagator $`g^{h_v}`$ , so that we get a connected Feynman graph, if we represent as single points all the clusters associated with the vertices immediately following $`v`$. These graphs have the property that the set of lines connecting the endpoints of the cluster $`L_v`$ and having scale $`h^{}h_v`$ is a connected subgraph; by the way this property is indeed another constraint on the possible choices of $`𝐏`$. We shall call these graphs compatible with $`𝐏`$. 3.3 The representation (3.37) of $`V^{(h)}(\tau ,𝐏)`$ is based on the choice of representing the regularization as acting on the kernels. If we use instead the representation of $``$ based on (3.10), (3.11), (3.19) and (3.24), some field variables have to be substituted with new ones, depending on two space-time points and containing possibly some derivatives. As we shall see, these new variables allow to get the right dimensional bounds, at the price of making much more involved the combinatorics. Hence, it is convenient to introduce a label $`r_v(f)`$ to keep trace of the regularization in the vertices of the tree where $`f`$ is associated with an external field and the action of $``$ turns out to be non trivial, that is $`1`$. There are many vertices, where $`=1`$ by definition, that is the vertices with more than $`4`$ external fields, the endpoints and $`v_0`$. For these vertices all external fields will be associated with a label $`r_v(f)=0`$. Moreover, since $`=0`$, the action of $``$ is trivial even in most trivial vertices $`v`$ with $`\|P_v\|4`$. This happens if the vertex (trivial or not) $`\stackrel{~}{v}`$ immediately following $`v`$ has the same number of external fields as $`v`$, since then the kernels associated with $`v`$ and $`\stackrel{~}{v}`$ are identical, up to a rescaling constant. In particular, this remark implies that, given the non trivial vertex $`v`$ and the non trivial vertex $`v^{}`$ immediately preceding $`v`$ on the tree, there are at most two vertices $`\overline{v}`$, such that $`v^{}<\overline{v}v`$ and the action of $``$ is non trivial. For the same reason, given an endpoint $`v`$ of scale $`h_v=+2`$ of type $`\lambda `$ (hence not local), there are at most two vertices between $`v`$ and the non trivial vertex $`v^{}`$ immediately preceding $`v`$, where the action of $``$ is non trivial. Since the number of endpoints is $`n`$ and the number of non trivial vertices is bounded by $`n1`$, the number of vertices where the action of $``$ is non trivial is bounded by $`2(2n1)`$. Let us now consider one of these vertices, which all have $`4`$ or $`2`$ external fields. If $`\|P_v\|=2`$ and the $`\omega `$ indices of the external fields are equal, we keep trace of the regularization by labeling the field variable, which is substituted with a $`D^2`$ field, see (3.19), with $`r_v(f)=2`$ and the other with $`r_v(f)=0`$. In principle we are free to decide which variable is labeled with $`r_v(f)=2`$, that is how we fix the localization point; we make a choice in the following way. If there is no non trivial vertex $`v^{}`$ such that $`v_0v^{}<v`$, we make an arbitrary choice, otherwise we put $`r_v(f)=2`$ for the field which is an internal field in the nearest non trivial vertex preceding $`v`$. In other words, we try to avoid that a field affected by the regularization stays external in the vertices preceding $`v`$. If $`\|P_v\|=2`$ and the $`\omega `$ indices of the external fields are different, we label the field variable, which is substituted with a $`D^{1,2}`$ field, see (3.24), with $`r_v(f)=1`$ and the other with $`r_v(f)=0`$; which variable is labeled with $`r_v(f)=1`$ is decided as in the previous case. If $`\|P_v\|=4`$, first of all we choose the localization point in the following way. If there is a vertex $`v^{}`$ such that $`v_0v^{}<v`$ and $`P_v^{}`$ contains one and only one $`fP_v`$, we chose $`𝐱(f)`$ as the localization point in $`v`$; in the other cases, we make an arbitrary choice. After that, we split the kernel associated with $`v`$ into three terms as in (3.14); then we distinguish the three terms by putting $`r_v(f)=1`$ for the external field which is substituted with a $`D^{1,1(h)}`$ field, when the delta functions are eliminated, and $`r_v(f)=0`$ for the others. The previous definitions imply that, given $`fI_{v_0}`$, it is possible that there are many different vertices in the tree, such that $`r_v(f)0`$, that is many vertices where the corresponding field variable appears as an external field and the action of $``$ is non trivial. As a consequence, the expressions given in §3.1 for the regularized potentials would not be sufficient and we should consider more general expressions, containing as external fields more general variables. Even worse, there is the risk that field derivatives of arbitrary order have to be considered; this event would produce “bad” factorials in the bounds. Fortunately, we can prove that this phenomenon can be easily controlled, thanks to our choice of the localization point, see above, by a more careful analysis of the regularization procedure, that we shall keep trace of by changing the definition of the $`r_v(f)`$ labels. Let us suppose first that $`\|P_v\|=4`$ and that there is $`fP_v`$, such that $`r_{\overline{v}}(f)0`$ for some $`\overline{v}>v`$. We want to show that the action of $``$ on $`v`$ is indeed trivial; hence we can put $`r_v(f)=0`$ for all $`fP_v`$, in agreement with the fact that the contribution to the effective potential associated with $`v`$ is dimensionally irrelevant. First of all, note that it is not possible that $`\|P_{\overline{v}}\|=2`$, as a consequence of the choice of the localization point in the vertices with two external fields, see above. On the other hand, if $`\|P_{\overline{v}}\|=4`$, the fact that the action of $``$ in the vertex $`v`$ is equal to the identity follows from the observation following (3.13) and the definition (2.72). Let us now consider the vertices $`v`$ with $`P_v=(f_1,f_2)`$. We can exclude as before that $`r_{\overline{v}}(f_i)0`$ for $`i=1`$ or $`i=2`$ or both and $`\|P_{\overline{v}}\|=2`$. The same conclusion can be reached, if there is no vertex $`\overline{v}>v`$, such that $`\|P_{\overline{v}}\|=4`$, the action of $``$ on $`\overline{v}`$ is non trivial and both $`f_1`$ and $`f_2`$ belong to the set of its external fields; this claim easily follows from the criterion for the choice of the localization point in the vertices with $`4`$ external fields. If, on the contrary, $`f_1`$ and $`f_2`$ are both labels of external fields of a vertex $`\overline{v}>v`$, such that $`\|P_{\overline{v}}\|=4`$ and the action of $``$ is non trivial, we have to distinguish two possibilities. If there is a non trivial vertex $`v^{}`$ such that $`v_0v^{}<v`$, and one of the external fields of $`v`$, let us say of label $`f_1`$, is an internal field, our choice of the localization points imply that both $`r_v(f_1)`$ and $`r_{\overline{v}}(f_1)`$ are different from $`0`$, while $`r_v(f_2)=r_{\overline{v}}(f_2)=0`$. If there is no non trivial vertex $`v^{}<v`$ with the previous property, that is if $`f_1`$ and $`f_2`$ are both labels of external fields down to $`v_0`$ (hence all vertices between $`v`$ and $`v_0`$ are trivial) or they become together labels of internal fields in some vertex $`v^{}<v`$, we are still free to choose as we want the localization points in $`v`$ and $`\overline{v}`$; we decide to choose them equal. The previous discussion implies that, as a consequence of our prescriptions, a field variable can be affected by the regularization only once, except in the case considered in the last paragraph. However, also in this case, it is easy to see that everything works as we did not apply to the variable with label $`f_1`$ the regularization in the vertex $`\overline{v}`$. In fact, the first or second order zero (modulo $`(L,\beta )`$) in the difference $`𝐱(f_1)𝐱(f_2)`$, related to the regularization in the vertex $`v`$, see §3.1, cancels the contribution of the term proportional to the delta function, related with the regularization of $`\overline{v}`$, see (3.14). This apparent lack of regularization in $`\overline{v}`$ is compensated by the fact that $`𝐱(f_1)𝐱(f_2)`$ is of order $`\gamma ^{h_{\overline{v}}}`$, hence smaller than the factor $`\gamma ^{h_v}`$ sufficient for the regularization of $`v`$ (together with the improving effect of the field derivative). Hence there is a gain with respect to the usual bound of a factor $`\gamma ^{(h_{\overline{v}}h_v)}`$, sufficient to regularize the vertex $`\overline{v}`$. 3.4 There is in principle another problem. Let us suppose that we decide to represent all the non trivial $``$ operations as acting on the field variables. Let us suppose also that the field variable with label $`f`$ is substituted, by the action of $``$ on the vertex $`v`$, with a $`D_{𝐲,𝐱}^{1,i}`$ or a $`D_{𝐲,𝐱}^2`$ field, where $`𝐲=𝐱(f)`$ and $`𝐱=𝐱(f^{})`$ is the corresponding localization point. At first sight it seems possible that even the variable with label $`f^{}`$ can be substituted with a $`D^{1,i}`$ or a $`D^2`$ field by the action of $``$ on a vertex $`\overline{v}>v`$. If this happens, the point $`𝐱(f^{})`$ can not be considered as fixed and there is an “interference” between the two regularization operations, or even more than two, since this phenomenon could involve an ordered chain of vertices. This interference would not produce bad factorials in the bounds, but would certainly make more involved our expansion. However, we can show that, thanks to our localization prescription, this problem is not really present. Let us suppose first that $`\|P_v\|=2`$. In this case, if the field with label $`f^{}`$ is external in some vertex $`\overline{v}>v`$, with $`\|P_{\overline{v}}\|`$ equal to $`2`$ or $`4`$, we are sure that $`𝐱(f^{})`$ is the localization point in $`\overline{v}`$, see §3.3, hence the corresponding filed can not be affected by the action of $``$ on $`\overline{v}`$. The same conclusion can be reached, if $`\|P_v\|=4`$ and $`\|P_{\overline{v}}\|=2`$ If $`\|P_v\|=\|P_{\overline{v}}\|=4`$ and the field with label $`f^{}`$ is substituted, by the action of $``$ on the vertex $`\overline{v}`$, with a $`D^{1,i}`$ or a <sup>2</sup> field, we know that the same can not be true for the field with label $`f`$, since the action of $``$ on $`v`$ is trivial. The previous discussion implies that the field with label $`f^{}`$ can be affected by the regularization (if $`\|P_v\|=\|P_{\overline{v}}\|=4`$) only by changing its $`𝐱`$ label, but this is not a source of any problem. 3.5 In this section we want to discuss the representation of the fields $`D_{𝐲,𝐱,\omega }^{1,i(h)\sigma }`$, $`i=1,2`$, and $`D_{𝐲,𝐱,\omega }^{2(h)\sigma }`$ introduced in §3.1, which allows to exploit the regularization properties of the $``$ operation. In order to do that, we extend the definition of the fields $`\psi _{𝐱,\omega }^{(h)\sigma }`$ to $`\text{}^2`$, by using (2.49); we get functions with values in the Grassman algebra, antiperiodic in $`x_0`$ and $`x`$ with periods $`\beta `$ and $`L`$, respectively. Let us choose a family of positive functions $`\chi _{\eta ,\eta ^{}}(𝐱)`$, $`\eta ,\eta ^{}\{1,0,+1\}`$, on $`\text{}^2`$, such that $$\begin{array}{cc}\hfill \chi _{\eta ,\eta ^{}}(𝐱)& =\{\begin{array}{cc}1\hfill & \text{if }\|x\eta \|1/4\text{ and }\|x_0\eta ^{}\|1/4\hfill \\ 0\hfill & \text{if }\|x\eta \|3/4\text{ or }\|x_0\eta ^{}\|3/4\hfill \end{array}\hfill \\ \hfill \underset{\eta ,\eta ^{}}{}\chi _{\eta ,\eta ^{}}(𝐱)& =1\text{if }𝐱[1,1]\times [1,1].\hfill \end{array}$$ $`(3.42)`$ Given $`𝐱,𝐲\mathrm{\Lambda }\times [\beta /2,\beta /2]`$, if $`\chi _{\eta ,\eta ^{}}(\stackrel{~}{𝐲}\stackrel{~}{𝐱})>0`$, where $`\stackrel{~}{𝐱}=(x/L,x_0/\beta )`$ and $`\stackrel{~}{𝐲}=(y/L,y_0/\beta )`$, we can define $`\overline{𝐲}=𝐲(\eta L,\eta ^{}\beta )`$, so that $`\|x_0\overline{y}_0\|3\beta /4`$ and $`\|x\overline{y}\|3L/4`$. We see immediately that $`D_{𝐲,𝐱,\omega }^{1,1(h)\sigma }=(1)^{\eta +\eta ^{}}D_{\overline{𝐲},𝐱,\omega }^{1,1(h)\sigma }`$ and we can write $$D_{\overline{𝐲},𝐱,\omega }^{1,1(h)\sigma }=[\psi _{\overline{𝐲},\omega }^{(h)\sigma }\psi _{𝐱,\omega }^{(h)\sigma }]+[1G_\sigma (\overline{𝐲}𝐱)]\psi _{𝐱,\omega }^{(h)\sigma }.$$ $`(3.43)`$ It is easy to see that, if $`\|y_0\|3\beta /4`$ and $`\|y\|3L/4`$, $$1G_\sigma (𝐲)=\frac{1}{L}\overline{h}_1(\stackrel{~}{𝐲})d_L(y)+\frac{1}{\beta }\overline{h}_2(\stackrel{~}{𝐲})d_\beta (y_0),\stackrel{~}{𝐲}=(y/L,y_0/\beta ),$$ $`(3.44)`$ where $`\overline{h}_i(𝐲)`$, $`i=1,2`$, are suitable functions, uniformly smooth in $`L`$ and $`\beta `$. Moreover $$\psi _{\overline{𝐲},\omega }^{(h)\sigma }\psi _{𝐱,\omega }^{(h)\sigma }=(\overline{𝐲}𝐱)_0^1𝑑t\psi _{\text{¸}(t),\omega }^{(h)\sigma },\text{¸}(t)=𝐱+t(\overline{𝐲}𝐱),$$ $`(3.45)`$ where $`=(_1,_0)`$ is the gradient, and it is easy to see that, if $`\|y_0\|3\beta /4`$ and $`\|y\|3L/4`$, $$𝐲=(\overline{h}_3(\stackrel{~}{𝐲})d_L(y),\overline{h}_4(\stackrel{~}{𝐲})d_\beta (y_0)),$$ $`(3.46)`$ where $`\overline{h}_i(𝐲)`$, $`i=3,4`$, are other suitable functions, uniformly smooth in $`L`$ and $`\beta `$. Hence we can write $$\begin{array}{ccc}& D_{𝐲,𝐱,\omega }^{1,1(h)\sigma }=\underset{\eta ,\eta ^{}}{}\{[\frac{1}{L}h_{1,\eta ,\eta ^{}}(\stackrel{~}{𝐲},\stackrel{~}{𝐱})d_L(yx)+\frac{1}{\beta }h_{2,\eta ,\eta ^{}}(\stackrel{~}{𝐲},\stackrel{~}{𝐱})d_\beta (y_0x_0)]\psi _{𝐱,\omega }^{(h)\sigma }+\hfill & (3.47)\hfill \\ & +h_{3,\eta ,\eta ^{}}(\stackrel{~}{𝐲},\stackrel{~}{𝐱})d_L(yx)_0^1dt_1\psi _{\text{¸}(t),\omega }^{(h)\sigma }+h_{4,\eta ,\eta ^{}}(\stackrel{~}{𝐲},\stackrel{~}{𝐱})d_\beta (y_0x_0)_0^1dt_0\psi _{\text{¸}(t),\omega }^{(h)\sigma }\},\hfill & \end{array}$$ where $$h_{i,\eta ,\eta ^{}}(\stackrel{~}{𝐲},\stackrel{~}{𝐱})=(1)^{\eta +\eta ^{}}\chi _{\eta ,\eta ^{}}(\stackrel{~}{𝐲}\stackrel{~}{𝐱})\overline{h}_i((\overline{y}x)/L,(\overline{y}_0x_0)/\beta ),i=1,4,$$ $`(3.48)`$ are smooth functions with support in the region $`\{\|yx\eta L\|3L/4,\|y_0x_0\eta ^{}\beta \|3\beta /4\}`$, such that their derivatives of order $`n`$ are bounded by a constant (depending on $`n`$) times $`\gamma ^{nh_{L,\beta }}`$. A similar expression is valid for $`D_{𝐲,𝐱,\omega }^{1,2(h)\sigma }`$. Let us now consider $`D_{𝐲,𝐱,\omega }^{2(h)\sigma }`$, see (3.20). We can write $$D_{𝐲,𝐱,\omega }^{2(h)\sigma }=(1)^{\eta +\eta ^{}}\stackrel{~}{D}_{\overline{𝐲},𝐱,\omega }^{2(h)\sigma }+h(\stackrel{~}{𝐲}\stackrel{~}{𝐱})d_L(yx)\overline{}_1^2\psi _{𝐱,\omega }^{(h)\sigma },$$ $`(3.49)`$ where $`h(𝐲𝐱)`$ is a uniformly smooth function and $$\begin{array}{cc}& \stackrel{~}{D}_{\overline{𝐲},𝐱,\omega }^{2(h)\sigma }=\psi _{\overline{𝐲},\omega }^{(h)\sigma }\psi _{𝐱,\omega }^{(h)\sigma }(\overline{𝐲}𝐱)\psi _{𝐱,\omega }^{(h)\sigma }\hfill \\ & \psi _{𝐱,\omega }^{(h)\sigma }\{[c_\beta (\overline{y}_0x_0)c_L(\overline{y}x)1]+b_Lc_\beta (\overline{y}_0x_0)d_L(\overline{y}x)\}\hfill \\ & \overline{}_1\psi _{𝐱,\omega }^{(h)\sigma }\{[c_\beta (\overline{y}_0x_0)1]d_L(\overline{y}x)+[d_L(\overline{y}x)(\overline{y}x)]\}\hfill \\ & (\overline{y}x)[\overline{}_1\psi _{𝐱,\omega }^{(h)\sigma }_1\psi _{𝐱,\omega }^{(h)\sigma }]_0\psi _{𝐱,\omega }^{(h)\sigma }[d_\beta (\overline{y}_0x_0)c_L(yx)(\overline{y}_0x_0)].\hfill \end{array}$$ $`(3.50)`$ Note that $$\overline{}_1\psi _{𝐱,\omega }^{(h)\sigma }_1\psi _{𝐱,\omega }^{(h)\sigma }=\frac{i\sigma }{L\beta }\underset{𝐤^{}𝒟_{L,\beta }^{}}{}e^{i\sigma 𝐤^{}𝐱}(\mathrm{sin}k^{}k^{})\widehat{\psi }_{𝐤^{},\omega }^{(h)\sigma }$$ $`(3.51)`$ behaves dimensionally as $`_1^3\psi _{𝐱,\omega }^{(h)\sigma }`$, hence we shall define $$\overline{}_1^3\psi _{𝐱,\omega }^{(h)\sigma }=\overline{}_1\psi _{𝐱,\omega }^{(h)\sigma }_1\psi _{𝐱,\omega }^{(h)\sigma }.$$ $`(3.52)`$ It is now easy to show that there exist functions $`h_{\underset{¯}{n},\eta ,\eta ^{}}(𝐲,𝐱)`$, with $`\underset{¯}{n}=(n_1,\mathrm{},n_6)`$, and $`h_{i,j,\eta ,\eta ^{}}(𝐲,𝐱)`$, $`i,j=0,1`$, smooth uniformly in $`L`$ and $`\beta `$, such that $$\begin{array}{ccc}\hfill D_{𝐲,𝐱,\omega }^{2(h)\sigma }& =\underset{\eta ,\eta ^{}}{}\{\underset{\underset{¯}{n}}{}h_{\underset{¯}{n},\eta ,\eta ^{}}(\stackrel{~}{𝐲},\stackrel{~}{𝐱})d_L(yx)^{n_1}d_\beta (y_0x_0)^{n_2}L^{n_3}\beta ^{n_4}\overline{}_1^{n_5}_0^{n_6}\psi _{𝐱,\omega }^{(h)\sigma }+\hfill & \\ & +\underset{i,j}{}h_{i,j,\eta ,\eta ^{}}(\stackrel{~}{𝐲},\stackrel{~}{𝐱})d_i(𝐲𝐱)d_i(𝐲𝐱)_0^1dt(1t)_i_j\psi _{\text{¸}(t),\omega }^{(h)\sigma }\},\hfill & (3.53)\hfill \end{array}$$ the sum over $`\underset{¯}{n}`$ being constrained by the conditions $$n_1+n_22,3\underset{i=3}{\overset{6}{}}n_i2.$$ $`(3.54)`$ 3.6 In order to exploit the regularization properties of formulas like (3.47) or (3.53), one has to prove that the “zeros” $`d_L(yx)`$ and $`d_\beta (y_0x_0)`$ give a contribution to the bounds of order $`\gamma ^h^{}`$, with $`h^{}h`$, if $`h`$ is the scale at which the zero was produced by the action of $``$. In §$`\mathrm{}`$3.7 we shall realize this task by “distributing” the zeros along a path connecting a family of space-time points associated with a subset of field variables. Let $`𝐱_0=𝐱,𝐱_1,\mathrm{},𝐱_n=𝐲`$ be the family of points connected by the path; it is easy to show that $$d_L(yx)=\underset{r=1}{\overset{n}{}}d_L(x_rx_{r1})e^{i\frac{\pi }{L}(x_r+x_{r1}x_nx_0)}.$$ $`(3.55)`$ A similar expression is valid for $`d_\beta (y_0x_0)`$. It can happen that one of the terms in the r.h.s. of (3.55) or the analogous expansion for $`d_\beta (y_0x_0)`$ depends on the same space-time points as the integration variables in the r.h.s. of a term like (3.21) or (3.26). We want to study the effect of this event. Let us call $`W(𝐱𝐲)`$ the kernel appearing in the l.h.s. of (3.21) or (3.26), $`W_R(𝐱𝐲)`$ its regularization, that is the quantity appearing in braces in the corresponding r.h.s., and let us define $$I_{n_1,n_2}=𝑑𝐱𝑑𝐲\psi _{𝐱,\omega }^{(h)+}\psi _{𝐲,\omega }^{(h)}W_R(𝐱𝐲)[e^{i\pi \frac{y}{L}}d_L(yx)]^{n_1}[e^{i\pi \frac{y_0}{\beta }}d_\beta (y_0x_0)]^{n_2}.$$ $`(3.56)`$ In the following we shall meet such expressions for values of $`n_1`$ and $`n_2`$, such that $`1n_1+n_22`$. If $`W(𝐱𝐲)`$ is the kernel appearing in the l.h.s. of (3.26), it is easy to see that, if $`n_1+n_21`$, $$I_{n_1,n_2}=𝑑𝐱𝑑𝐲\psi _{𝐱,\omega }^{(h)+}\psi _{𝐲,\omega }^{(h)}W(𝐱𝐲)[e^{i\pi \frac{y}{L}}d_L(yx)]^{n_1}[e^{i\pi \frac{y_0}{\beta }}d_\beta (y_0x_0)]^{n_2},$$ $`(3.57)`$ that is the presence of the zeros simply erases the effect of the regularization. Let us now suppose that $`W(𝐱𝐲)`$ is the kernel appearing in the l.h.s. of (3.21) and $`W_R(𝐱𝐲)`$ its regularization. We have $$\begin{array}{cc}\hfill I_{1,0}& =d𝐱d𝐲\psi _{𝐱,\omega }^{(h)+}W(𝐱𝐲)d_L(yx)\{D_{𝐲,𝐱,\omega }^{1,3(h)}\hfill \\ & c_\beta (y_0x_0)[\frac{1}{2}\overline{}_1^2(e^{i\pi \frac{x}{L}}\psi _{𝐱,\omega }^{(h)})+\frac{i\mathrm{cos}p_F}{v_0}\overline{}_1(e^{i\pi \frac{x}{L}}\psi _{𝐱,\omega }^{(h)})]\},\hfill \end{array}$$ $`(3.58)`$ $$I_{0,1}=𝑑𝐱𝑑𝐲\psi _{𝐱,\omega }^{(h)+}W(𝐱𝐲)d_\beta (y_0x_0)D_{𝐲,𝐱,\omega }^{1,4(h)},$$ $`(3.59)`$ where $$D_{𝐲,𝐱,\omega }^{1,3(h)}=e^{i\pi \frac{y}{L}}\psi _{𝐲,\omega }^{(h)}c_\beta (y_0x_0)e^{i\pi \frac{x}{L}}\psi _{𝐱,\omega }^{(h)}.$$ $`(3.60)`$ $$D_{𝐲,𝐱,\omega }^{1,4(h)}=e^{i\pi \frac{y_0}{\beta }}\psi _{𝐲,\omega }^{(h)}c_L(yx)e^{i\pi \frac{x_0}{\beta }}\psi _{𝐱,\omega }^{(h)}.$$ $`(3.61)`$ Moreover $$\begin{array}{ccc}& I_{2,0}=d𝐱d𝐲\psi _{𝐱,\omega }^{(h)+}W(𝐱𝐲)d_L(yx)\{d_L(yx)e^{2i\pi \frac{y}{L}}\psi _{𝐲,\omega }^{(h)}\frac{c_\beta (y_0x_0)}{a_L}\hfill & \\ & [\overline{}_1(e^{2i\pi \frac{x}{L}}\psi _{𝐱,\omega }^{(h)})+\frac{i\mathrm{cos}p_F}{v_0}(e^{2i\pi \frac{x}{L}}\psi _{𝐱,\omega }^{(h)}+\frac{1}{2}\overline{}_1^2(e^{2i\pi \frac{x}{L}}\psi _{𝐱,\omega }^{(h)}))]\},\hfill & (3.62)\hfill \end{array}$$ $$I_{0,2}=𝑑𝐱𝑑𝐲\psi _{𝐱,\omega }^{(h)+}W(𝐱𝐲)d_\beta (y_0x_0)^2e^{2i\pi \frac{y_0}{\beta }}\psi _{𝐲,\omega }^{(h)},$$ $`(3.63)`$ $$I_{1,1}=𝑑𝐱𝑑𝐲\psi _{𝐱,\omega }^{(h)+}W(𝐱𝐲)d_L(yx)d_\beta (y_0x_0)e^{i\pi \frac{y}{L}i\pi \frac{y_0}{\beta }}\psi _{𝐲,\omega }^{(h)}.$$ $`(3.64)`$ Note that no cancellations are possible for $`𝐱=𝐲`$ modulo $`(L,\beta )`$ between the various terms contributing to $`I_{n_1,n_2}`$; hence they will be bounded separately. Note also that the fields $`D_{𝐲,𝐱,\omega }^{1,3(h)}`$ and $`D_{𝐲,𝐱,\omega }^{1,4(h)}`$ have a zero of first order for $`𝐱=𝐲`$ modulo $`(L,\beta )`$ and can be represented by expressions analogous to the r.h.s. of (3.47). Moreover, the terms contributing to $`I_{0,1}`$ and $`I_{1,0}`$ and containing these fields can also be written in a form analogous to (3.26). Finally, we want to stress the fact that the integrands in the previous expressions of $`I_{n_1,n_2}`$, $`1n_1+n_22`$, have a zero of order at most two for $`𝐱=𝐲`$ modulo $`(L,\beta )`$, that is a zero of order not higher of the zero introduced in the r.h.s. of (3.56). As it will be more clear in §$`\mathrm{}`$3.7, this property would be lost if one uses the representation (3.19) of the regularization operation, before performing the ”decomposition of the zeros”; one should get in this case a zero of order four and the iteration of the procedure of decomposition of the zeros would produce zeros of arbitrary order and, as a consequence, bad combinatorial factors in the bounds. 3.7 We are now ready to describe in more detail our expansion. First of all, we insert the decomposition (3.14) of $`V^{(h)}(\tau ,\psi ^{(h)})`$ in the vertices with $`\|P_v\|=4`$, by following the prescription for the choice of the localization point described in §3.3. The discussion of §3.3 allows also to define a new label $`r(f)`$, to be called the $``$-label, for any $`fI_{v_0}`$, by putting (i) $`r(f)=0`$, if $`r_v(f)=0`$ for any $`v`$ such that $`fP_v`$; (ii) $`r(f)=(i,v)`$, if there exists one and only one vertex $`v`$, such that $`fP_v`$ and $`r_v(f)=i0`$; (iii) $`r(f)=(2,v,\overline{v})`$, if there are two vertices $`v`$ and $`\overline{v}`$, such that $`v<\overline{v}`$, $`fP_vP_{\overline{v}}`$, $`\|P_v\|=2`$, $`\|P_{\overline{v}}\|=4`$, $`r_v(f)=2`$, $`r_{\overline{v}}(f)=1`$; see discussion in the last two paragraphs of §3.3. Then, we can write $$V^{(h)}(\tau ,\sqrt{Z_h}\psi ^{(h)})=\underset{𝐏𝒫_\tau ,𝐫}{}V^{(h)}(\tau ,𝐏,𝐫),$$ $`(3.65)`$ where $`𝐫=\{r(f),fI_{v_0}\}`$ and the sum over $`𝐫`$ must be understood as the sum over the possible choices of $`𝐫`$ compatible with $`𝐏`$. We can also write $$V^{(h)}(\tau ,𝐏,𝐫)=\sqrt{Z_h}^{\|P_{v_0}\|}𝑑𝐱_{v_0}K_{\tau ,𝐏,𝐫}^{(h)}(𝐱_{v_0})\stackrel{~}{\psi }^{(h)}(P_{v_0}),$$ $`(3.66)`$ with $`K_{\tau ,𝐏,𝐫}^{(h)}(𝐱_{v_0})`$ defined inductively as in (3.38). Let us consider first the action of $``$ on $`V^{(h)}(\tau ,𝐏,𝐫)`$. We can write for $`V^{(h)}(\tau ,𝐏,𝐫)`$ an expression similar to (3.66), if we continue to use for the $``$ operation the representation based on (3.14),(3.21) and (3.26), which affects the kernels leaving the fields unchanged. We shall use the notation $$V^{(h)}(\tau ,𝐏,𝐫)=𝑑𝐱_{v_0}\stackrel{~}{\psi }^{(h)}(P_{v_0})[K_{\tau ,𝐏,𝐫}^{(h)}(𝐱_{v_0})].$$ $`(3.67)`$ Moreover, we define $`𝐫^{}`$ so that $`r^{}(f)=r(f)`$ except for the field labels $`fP_{v_0}`$, for which $`r^{}(f)`$ takes into account also the regularization acting on $`v_0`$. However, we can use for the $``$ operation also the representation based on (3.10), (3.11), (3.19) and (3.24), which can be derived from the previous one by integrating the $`\delta `$-functions; the effect is to replace one of the external fields with one of the fields $`D^{1,i(h)\sigma }`$, $`i=1,2`$ or $`D^{2(h)\sigma }`$. We can describe the result by writing $$V^{(h)}(\tau ,𝐏,𝐫)=𝑑𝐱_{v_0}[\stackrel{~}{\psi }^{(h)}(P_{v_0})]K_{\tau ,𝐏,𝐫}^{(h)}(𝐱_{v_0}).$$ $`(3.68)`$ The discussion in §3.3 and §3.5 implies that there is a finite set $`A_{v_0}`$, such that $$[\stackrel{~}{\psi }^{(h)}(P_{v_0})]=\underset{\alpha A_{v_0}}{}h_\alpha (\stackrel{~}{𝐱}_{P_{v_0}})d_L^{n_1(\alpha )}d_\beta ^{n_2(\alpha )}\underset{fP_{v_0}}{}[\widehat{}_{j_\alpha (f)}^{q_\alpha (f)}\psi ]_{𝐱_\alpha (f),\omega (f)}^{(h)\sigma (f)},$$ $`(3.69)`$ where $`\stackrel{~}{𝐱}_{P_{v_0}}=(L^1x_{P_{v_0}},\beta ^1x_{0P_{v_0}})`$, $`d_L^{n_1(\alpha )}`$ and $`d_\beta ^{n_2(\alpha )}`$ are powers of the functions (2.96), with argument the difference of two points belonging to $`𝐱_{P_{v_0}}`$, and $`\widehat{}_j^q`$, $`q=0,1,2`$, $`j=1,\mathrm{},m_q`$, is a family of operators acting on the field variables, which are dimensionally equivalent to derivatives of order $`q`$. In particular $`m_0=1`$, $`\widehat{}_1^0`$ is the identity and the action of $``$ is trivial, that is $`\|A_{v_0}\|=1`$, $`h_\alpha =1`$, $`n_1(\alpha )=n_2(\alpha )=0`$ and $`q_\alpha (f)=0`$ for any $`fP_{v_0}`$, except in the following cases. 1) If $`\|P_{v_0}\|=4`$ and $`r(f)=0`$ for any $`fP_{v_0}`$, there is $`\overline{f}P_{v_0}`$, such that the action of $``$ over the fields consists in replacing one of the field variables with a $`D_{𝐲,𝐱,\omega }^{1,1(h)\sigma }`$ field, where $`𝐲=𝐱(\overline{f})`$ and $`𝐱=𝐱(f)`$ for some other $`fP_{v_0}`$, see (3.11); moreover, one or two of the other fields change their space-time point. We write $`D_{𝐲,𝐱,\omega }^{1,1(h)\sigma }`$ in the representation (3.47); the resulting expression is of the form (3.69), with $`A_{v_0}`$ consisting of four different terms, such that $`d_L=d_L(yx)`$, $`d_\beta =d_\beta (y_0x_0)`$, $`n_1(\alpha )+n_2(\alpha )=1`$ and, for all $`f\overline{f}`$, $`q_\alpha (f)=0`$, while $`q_\alpha (\overline{f})=1`$. Moreover, if $`f\overline{f}`$, $`𝐱_\alpha (f)`$ is a single point belonging to $`𝐱_{P_{v_0}}`$, not necessarily coinciding with $`𝐱(f)`$, while, if $`f=\overline{f}`$, $`𝐱_\alpha (f)`$ is equal to $`𝐱`$ or to the couple $`(𝐱,𝐲)`$ (using the previous definitions). The precise values of $`𝐱_\alpha (\overline{f})`$ and $`[\widehat{}_{j_\alpha (\overline{f})}^1\psi ]_{𝐱_\alpha (\overline{f}),\omega (\overline{f})}^{(h)\sigma (\overline{f})}`$, together with the functions $`h_\alpha `$, can be deduced from (3.47). 2) If $`P_{v_0}=(f_1,f_2)`$ and $`\omega (f_1)=\omega (f_2)`$, the action of $``$ consists in replacing one of the external fields, of label, let us say, $`f_1`$, with a $`D_{𝐲,𝐱,\omega }^{2(h)\sigma }`$ field, where $`𝐲=𝐱(f_1)`$ and $`𝐱=𝐱(f_2)`$, if $`f_2`$ is the second field label. By using the representation (3.53) of $`D_{𝐲,𝐱,\omega }^{2(h)\sigma }`$, we get an expression of the form (3.69) consisting of many different terms, such that $`d_L=d_L(yx)`$, $`d_\beta =d_\beta (y_0x_0)`$, $`n_1(\alpha )+n_2(\alpha )2`$, $`q_\alpha (f_1)=2`$, $`q_\alpha (f_2)=0`$, $`𝐱_\alpha (f_2)=𝐱(f_2)`$. The values of $`𝐱_\alpha (f_1)`$ and $`[\widehat{}_{j_\alpha (f_1)}^2\psi ]_{𝐱_\alpha (f_1),\omega (f_1)}^{(h)\sigma (f_1)}`$, together with the functions $`h_\alpha `$, can be deduced from (3.53). 3) If $`P_{v_0}=(f_1,f_2)`$ and $`\omega (f_1)=\omega (f_2)`$, the action of $``$ consists in replacing one of the external fields, of label, let us say, $`f_1`$, with a $`D_{𝐲,𝐱,\omega }^{1,2(h)\sigma }`$ field, where $`𝐲=𝐱(f_1)`$ and $`𝐱=𝐱(f_2)`$, if $`f_2`$ is the second field label. By using the analogous of the representation (3.47) for $`D_{𝐲,𝐱,\omega }^{1,2(h)\sigma }`$, we get an expression of the form (3.69) consisting of four different terms, such that $`n_1(\alpha )+n_2(\alpha )=1`$, $`q_\alpha (f_1)=1`$, $`q_\alpha (f_2)=0`$, $`𝐱_\alpha (f_2)=𝐱(f_2)`$. Let us now consider the action of $``$ on $`V^{(h)}(\tau ,\sqrt{Z_h}\psi ^{(h)})`$. We get an expansion similar to that based on (3.68), that we can write, by using (2.79), (3.65) and translation invariance, in the form $$\begin{array}{cc}\hfill V^{(h)}(\tau ,\sqrt{Z_h}\psi ^{(h)})& =\gamma ^hn_h(\tau )Z_hF_\nu ^{(h)}+s_h(\tau )Z_hF_\sigma ^{(h)}+z_h(\tau )Z_hF_\zeta ^{(h)}+\hfill \\ & +a_h(\tau )Z_hF_\alpha ^{(h)}+l_h(\tau )Z_h^2F_\lambda ^{(h)},\hfill \end{array}$$ $`(3.70)`$ where $$\begin{array}{cc}\hfill n_h(\tau )& =\frac{\gamma ^h}{L\beta }\underset{\genfrac{}{}{0pt}{}{𝐏𝒫_\tau ,𝐫}{P_{v_0}=(f_1,f_2),\omega (f_1)=\omega (f_2)=+1}}{}𝑑𝐱_{v_0}h_1(\stackrel{~}{𝐱}_{P_{v_0}})K_{\tau ,𝐏,𝐫}^{(h)}(𝐱_{v_0}),\hfill \\ \hfill s_h(\tau )& =\frac{1}{L\beta }\underset{\genfrac{}{}{0pt}{}{𝐏𝒫_\tau ,𝐫}{P_{v_0}=(f_1,f_2),\omega (f_1)=\omega (f_2)=+1}}{}𝑑𝐱_{v_0}h_2(\stackrel{~}{𝐱}_{P_{v_0}})K_{\tau ,𝐏,𝐫}^{(h)}(𝐱_{v_0}),\hfill \\ \hfill z_h(\tau )& =\frac{1}{L\beta }\underset{\genfrac{}{}{0pt}{}{𝐏𝒫_\tau ,𝐫}{P_{v_0}=(f_1,f_2),\omega (f_1)=\omega (f_2)=+1}}{}𝑑𝐱_{v_0}h_3(\stackrel{~}{𝐱}_{P_{v_0}})d_\beta (x(f_2)x(f_1))K_{\tau ,𝐏,𝐫}^{(h)}(𝐱_{v_0}),\hfill \\ \hfill a_h(\tau )& =\frac{1}{L\beta }\underset{\genfrac{}{}{0pt}{}{𝐏𝒫_\tau ,𝐫}{P_{v_0}=(f_1,f_2),\omega (f_1)=\omega (f_2)=+1}}{}𝑑𝐱_{v_0}h_4(\stackrel{~}{𝐱}_{P_{v_0}})d_L(x(f_2)x(f_1))K_{\tau ,𝐏,𝐫}^{(h)}(𝐱_{v_0}),\hfill \\ \hfill l_h(\tau )& =\frac{1}{L\beta }\underset{\genfrac{}{}{0pt}{}{𝐏𝒫_\tau ,𝐫}{\|P_{v_0}\|=4,\underset{¯}{\sigma }=(+,,+,),\underset{¯}{\omega }=(+1,1,1,+1)}}{}𝑑𝐱_{v_0}h_5(\stackrel{~}{𝐱}_{P_{v_0}})K_{\tau ,𝐏,𝐫}^{(h)}(𝐱_{v_0}),\hfill \end{array}$$ $`(3.71)`$ $`h_i(\stackrel{~}{𝐱}_{P_{v_0}})`$, $`i=1,\mathrm{},5`$, being bounded functions, whose expressions can be deduced from (3.8), (3.16) and (3.22), also taking into account the permutations needed to order the field variables as in the r.h.s. of (3.70). The constants $`n_h`$, $`s_h`$, $`z_h`$, $`a_h`$ and $`l_h`$, which characterize the local part of the effective potential, can be obtained from (3.71) by summing over $`n1`$ and $`\tau 𝒯_{h,n}`$. Finally, the constant $`\stackrel{~}{E}_{h+1}`$ appearing in the l.h.s. of (3.27) can be written in the form $$\stackrel{~}{E}_{h+1}=\underset{n=1}{\overset{\mathrm{}}{}}\underset{\tau 𝒯_{h,n}}{}\stackrel{~}{E}_{h+1}(\tau ),$$ $`(3.72)`$ where $$\stackrel{~}{E}_{h+1}(\tau )=\frac{1}{L\beta }\underset{\genfrac{}{}{0pt}{}{𝐏𝒫_\tau ,𝐫}{P_{v_0}=\mathrm{}}}{}𝑑𝐱_{v_0}K_{\tau ,𝐏,𝐫}^{(h)}(𝐱_{v_0}).$$ $`(3.73)`$ 3.8 We want now to iterate the previous procedure, by using equation (3.38), in order to suitably take into account the non trivial $``$ operations in the vertices $`vv_0`$. We shall focus our discussion on $`V^{(h)}(\tau ,𝐏,𝐫)`$, but the following analysis applies also to $`V^{(h)}(\tau ,𝐏,𝐫)`$ and $`\stackrel{~}{E}_{h+1}(\tau )`$. Let us consider the truncated expectation in the r.h.s. of (3.38) and let us put $`s=s_v`$, $`P_iP_{v_i}\backslash Q_{v_i}`$. Moreover we order in an arbitrary way the sets $`P_i^\pm \{fP_i,\sigma (f)=\pm \}`$, we call $`f_{ij}^\pm `$ their elements and we define $`𝐱^{(i)}=_{fP_i^{}}𝐱(f)`$, $`𝐲^{(i)}=_{fP_i^+}𝐱(f)`$, $`𝐱_{ij}=𝐱(f_{i,j}^{})`$, $`𝐲_{ij}=𝐱(f_{i,j}^+)`$. Note that $`_{i=1}^s\|P_i^{}\|=_{i=1}^s\|P_i^+\|n`$, otherwise the truncated expectation vanishes. A couple $`l(f_{ij}^{},f_{i^{}j^{}}^+)(f_l^{},f_l^+)`$ will be called a line joining the fields with labels $`f_{ij}^{},f_{i^{}j^{}}^+`$ and $`\omega `$ indices $`\omega _l^{},\omega _l^+`$ and connecting the points $`𝐱_l𝐱_{i,j}`$ and $`𝐲_l𝐲_{i^{}j^{}}`$, the endpoints of $`l`$; moreover we shall put $`m_lm(f_l^{})+m(f_l^+)`$. Then, it is well known (see \[Le\], \[BGPS\], for example) that, up to a sign, if $`s>1`$, $$\stackrel{~}{}_h^T(\stackrel{~}{\psi }^{(h)}(P_1),\mathrm{},\stackrel{~}{\psi }^{(h)}(P_s))=\underset{T}{}\underset{lT}{}\overline{}_1^{m_l}g_{\omega _l^{},\omega _l^+}^{(h)}(𝐱_l𝐲_l)𝑑P_T(𝐭)detG^{h,T}(𝐭),$$ $`(3.74)`$ where $`T`$ is a set of lines forming an anchored tree graph between the clusters of points $`𝐱^{(i)}𝐲^{(i)}`$, that is $`T`$ is a set of lines, which becomes a tree graph if one identifies all the points in the same cluster. Moreover $`𝐭=\{t_{i,i^{}}[0,1],1i,i^{}s\}`$, $`dP_T(𝐭)`$ is a probability measure with support on a set of $`𝐭`$ such that $`t_{i,i^{}}=𝐮_i𝐮_i^{}`$ for some family of vectors $`𝐮_i\text{}^s`$ of unit norm. Finally $`G^{h,T}(𝐭)`$ is a $`(ns+1)\times (ns+1)`$ matrix, whose elements are given by $`G_{ij,i^{}j^{}}^{h,T}=t_{i,i^{}}\overline{}_1^{m(f_{ij}^{})+m(f_{i^{}j^{}}^+)}g_{\omega _l^{},\omega _l^+}^{(h)}(𝐱_{ij}𝐲_{i^{}j^{}})`$ with $`(f_{ij}^{},f_{i^{}j^{}}^+)`$ not belonging to $`T`$. If $`s=1`$, the sum over $`T`$ is empty, but we shall still use equation (3.74), by interpreting the r.h.s. as $`1`$, if $`P_1`$ is empty (which is possible, for $`s=1`$), and as $`detG^h(\mathrm{𝟏})`$ otherwise. Inserting (3.74) in the r.h.s. of (3.38) (with $`v=v_0`$) we obtain, up to a sign, $$\begin{array}{cc}& V^{(h)}(\tau ,𝐏,𝐫)=\frac{1}{s_{v_0}!}\sqrt{Z_h}^{\|P_{v_0}\|}\underset{T_{v_0}}{}d𝐱_{v_0}dP_{T_{v_0}}(𝐭)[\stackrel{~}{\psi }^{(h)}(P_{v_0})]\hfill \\ & \left[\underset{lT_{v_0}}{}\overline{}_1^{m_l}g_{\omega _l^{},\omega _l^+}^{(h+1)}(𝐱_l𝐲_l)\right]detG^{h+1,T_{v_0}}(𝐭)\sqrt{\frac{Z_{h+1}}{Z_h}}^{\|P_{v_0}\|}\underset{i=1}{\overset{s_{v_0}}{}}[K_{v_i}^{(h+2)}(𝐱_{v_i})]\hfill \end{array}$$ $`(3.75)`$ Let us now consider the contribution to the r.h.s. of (3.75) of one of the terms in the representation (3.69) of $`\stackrel{~}{\psi }^{(h)}(P_{v_0})`$ with $`n_1(\alpha )+n_2(\alpha )>0`$. For each choice of $`T_{v_0}`$, we decompose the factors $`d_L^{n_1(\alpha )}(yx)`$ and $`d_\beta ^{n_2(\alpha )}(y_0x_0)`$, by using equation (3.55) and the analogous equation for $`d_\beta (y_0x_0)`$, with $`𝐱_0=𝐱`$, $`𝐱_n=𝐲`$ and the other points $`𝐱_r`$, $`r=1,\mathrm{},n1`$, chosen in the following way. Let us consider the unique subset $`(l_1,\mathrm{},l_m)`$ of $`T_{v_0}`$, which selects a path joining the cluster containing $`𝐱_0`$ with the cluster containing $`𝐱_n`$, if one identifies all the points in the same cluster. Let $`(\overline{v}_{i1},\overline{v}_i)`$, $`i=1,m`$, the couple of vertices whose clusters of points are joined by $`l_i`$. We shall put $`𝐱_{2i1}`$, $`i=1,m`$, equal to the endpoint of $`l_i`$ belonging to $`𝐱_{\overline{v}_{i1}}`$ and $`𝐱_{2i}`$ equal to the endpoint of $`l_i`$ belonging to $`𝐱_{\overline{v}_i}`$. This definition implies that there are two points of the sequence $`𝐱_r`$, $`r=0,\mathrm{},n=2m+1`$, possibly coinciding, in any set $`𝐱_{\overline{v}_i}`$, $`i=0,\mathrm{},m`$; these two points are the space-time points of two different fields belonging to $`P_{\overline{v}_i}`$. Since $`n2s_{v_0}1`$, this decomposition will produce a finite number of different terms ($`(2s_{v_0}1)^2`$, since $`n_1(\alpha )+n_2(\alpha )2`$), that we shall distinguish with a label $`\alpha ^{}`$ belonging to a set $`B_{v_0}`$, depending on $`\alpha A_{v_0}`$ and $`T_{v_0}`$. These terms can be described in the following way. Each term is obtained from the one chosen in the r.h.s. of (3.75) by adding a factor $`\mathrm{exp}\{i\pi L^1n_1(\alpha )(x+y)+i\pi \beta ^1n_2(\alpha )(x_0+y_0)\}`$. Moreover each propagator $`g_{\omega _l^{},\omega _l^+}^{(h+1)}(𝐱_l𝐲_l)`$ is multiplied by a factor $`d_{j_\alpha ^{}(l)}^{b_\alpha ^{}(l)}(𝐱_l,𝐲_l)`$, where $`d_j^b`$, $`d=0,1,2`$, $`j=1,\mathrm{},m_b`$ is a family of functions so defined. If $`b=0`$, $`m_0=1`$ and $`d_1^0=1`$. If $`b=1`$, $`m_b=2`$ and $`j`$ distinguishes, up to the sign, the two functions $$e^{i\frac{\pi }{L}(x_l+y_l)}d_L(x_ly_l),e^{i\frac{\pi }{\beta }(x_{0,l}+y_{0,l})}d_\beta (x_{l0}y_{l0}).$$ $`(3.76)`$ If $`b=2`$, $`j`$ distinguishes the three possibilities, obtained by taking the product of two factors equal to one of the terms in (3.76). Finally each one of the vertices $`v_1,\mathrm{},v_{s_{v_0}}`$ is multiplied by a similar factor $`d_{j_\alpha ^{}(v_i)}^{b_\alpha ^{}(v_i)}(𝐱_i,𝐲_i)`$. Note that the definitions were chosen so that $`\|d_j^b(𝐱,𝐲)\|\|𝐝(𝐱𝐲)\|^b`$. Moreover there is the constraint that $$\underset{lT_{v_0}}{}b_\alpha ^{}(l)+\underset{i=1}{\overset{s_{v_0}}{}}b_\alpha ^{}(v_i)=n_1(\alpha )+n_2(\alpha ).$$ $`(3.77)`$ The previous discussion implies that (3.75) can be written in the form $$\begin{array}{cc}& V^{(h)}(\tau ,𝐏,𝐫)=\frac{1}{s_{v_0}!}\sqrt{Z_h}^{\|P_{v_0}\|}\underset{\alpha A_{v_0}}{}\underset{T_{v_0}}{}\underset{\alpha ^{}B_{v_0}}{}d𝐱_{v_0}dP_{T_{v_0}}(𝐭)\hfill \\ & h_\alpha (\stackrel{~}{𝐱}_{P_{v_0}})\left[\underset{fP_{v_0}}{}(\widehat{}_{j_\alpha (f)}^{q_\alpha (f)}\psi )_{𝐱_\alpha (f),\omega (f)}^{(h)\sigma (f)}\right]\left[\underset{lT_{v_0}}{}d_{j_\alpha ^{}(l)}^{b_\alpha ^{}(l)}(𝐱_l,𝐲_l)\overline{}_1^{m_l}g_{\omega _l^{},\omega _l^+}^{(h+1)}(𝐱_l𝐲_l)\right]\hfill \\ & detG^{h+1,T_{v_0}}(𝐭)\sqrt{\frac{Z_{h+1}}{Z_h}}^{\|P_{v_0}\|}\left[\underset{i=1}{\overset{s_{v_0}}{}}d_{j_\alpha ^{}(v_i)}^{b_\alpha ^{}(v_i)}(𝐱_i,𝐲_i)K_{v_i}^{(h+2)}(𝐱_{v_i})\right],\hfill \end{array}$$ $`(3.78)`$ where the function $`h_\alpha (\stackrel{~}{𝐱}_{P_{v_0}})`$ has be redefined in order to absorb the factor $`\mathrm{exp}\{i\pi L^1n_1(\alpha )(x+y)+i\pi \beta ^1n_2(\alpha )(x_0+y_0)\}`$. 3.9 We are now ready to begin the iteration of the previous procedure, by considering those among the vertices $`v_1,\mathrm{},v_{s_{v_0}}`$, where the action of $``$ is non trivial. It turns out that we can not simply repeat the arguments used for $`v_0`$, but we have to consider some new situations and introduce some new prescriptions, which will be however sufficient to complete the iteration up to the endpoints, without any new problem. Let us select a term in the r.h.s. of (3.78) and one of the vertices immediately following $`v_0`$, let us say $`\overline{v}`$, where the action of $``$ is non trivial. We have to consider a few different cases. A) Suppose that $`b(\overline{v})=0`$ (we shall omit the dependence on $`\alpha `$ and $`\alpha ^{}`$). In this case the action of $``$ is exploited following essentially the same procedure as for $`v_0`$. If $``$ is different from the identity, we move its action on the external fields of $`\overline{v}`$, by using the analogous of (3.69), by taking into account that some of the external fields of $`\overline{v}`$ are internal fields of $`v_0`$, hence they are involved in the calculation of the truncated expectation (3.74). This means that, if $`f`$ is the label of an internal field with $`q(f)>0`$, the corresponding (non trivial) $`\widehat{}_{j(f)}^{q(f)}`$ operator acts on the quantities in the r.h.s. of (3.78), which depend on $`f`$, that is $`d_{j(l)}^{b(l)}(𝐱_l,𝐲_l)g_{\omega _l^{},\omega _l^+}^{(h+1)}(𝐱_l𝐲_l)`$ or the matrix elements of $`detG^{h+1,T_{v_0}}`$, which are obtained by contracting the field with label $`f`$ with another internal field. For example, if $`𝐱(f)=𝐱_l`$ and $`\widehat{}_{j(f)}^{q(f)}`$ is the operator associated with the third term in the r.h.s. of (3.47), we must substitute $`d_{j(l)}^{b(l)}(𝐱_l,𝐲_l)\overline{}_1^{m_l}g_{\omega _l^{},\omega _l^+}^{(h+1)}(𝐱_l𝐲_l)`$ with $$_0^1𝑑t_1[d_{j(l)}^{b(l)}(\text{¸}(t)𝐲_l)\overline{}_1^{m_l}g_{\omega _l^{},\omega _l^+}^{(h+1)}(\text{¸}(t)𝐲_l)],$$ $`(3.79)`$ with $`\text{¸}(t)=𝐱^{}+t(\overline{𝐱}_l𝐱^{})`$, for some $`𝐱^{}𝐱_{\overline{v}}`$, $`\overline{𝐱}_l`$ being defined in terms of $`𝐱_l`$ as $`\overline{y}`$ is defined in terms of $`y`$ in §3.5 (that is $`\overline{𝐱}_l`$ and $`𝐱_l`$ are equivalent representation of the same point on the space-time torus). There is apparently another problem, related to the possibilities that the operators $`\widehat{}_{j(f)}^{q(f)}`$ related with the action of $``$ on $`\overline{v}`$ do not commute with the functions $`h_\alpha `$ and the field variables introduced by the action of $``$ on $`v_0`$. However, the discussion in §3.4 implies that this can not happen, because of our prescription for the choice of the localization points. This argument is of general validity, hence we will not consider anymore this problem in the following. B) If $`b(\overline{v})>0`$, we shall proceed in a different way, in order to avoid growing powers of the factors $`d_L`$ and $`d_\beta `$, which should produce at the end bad combinatorial factors in the bounds. We need to distinguish four different cases. B1) If $`\|P_{\overline{v}}\|=4`$, we do not use the decomposition (3.47) for the field changed by the action of $``$ in a $`D^{1,1}`$ field, but we simply write it as the sum of the two terms in the r.h.s. of (3.12) (in some cases the second term does not really contributes, because the argument of the factor $`d_{j(\overline{v})}^{b(\overline{v})}`$ is the same as the argument of the delta function in the representation (3.14) of the $``$ action, but this is not true in general). We still get a representation of the form (3.69) for $`[\stackrel{~}{\psi }^{(h)}(P_{\overline{v}})]`$, but with the property that $`q(f)=0`$ for any $`\alpha A_{\overline{v}}`$ and any $`fP_{\overline{v}}`$. This procedure works, because we do not need to exploit the regularization property of $``$ in this case, as the following analysis will make clear. B2) If $`\|P_{\overline{v}}\|=2`$, and the $`\omega `$-labels of the external fields are different, the action of $``$, after the insertion of the zero, is indeed trivial, as explained in §3.6, see (3.57). Hence we do not make any change in the external fields. B3) If $`\|P_{\overline{v}}\|=2`$, the $`\omega `$-labels of the external fields are equal and $`b(\overline{v})=2`$, the presence of the factor $`d_{j(\overline{v})}^{b(\overline{v})}`$ does not allow to use for the action of $``$ on the external fields the representation (3.69), because that factor depends on the space-time labels of the external fields. However, we can use the representation following from the equations (3.62),(3.63),(3.64), by considering the different terms in the r.h.s. as different contributions (in any case no cancellations among such terms are possible). Note that this representation has the same properties of the representation (3.69) and can be written exactly in the same form, by suitable defining the various quantities. In particular, it is still true that $`n_1(\alpha )+n_2(\alpha )2`$. Of course, we have to take also into account that some of the external fields of $`\overline{v}`$ are internal fields of $`v_0`$, but this can be done exactly as in item A). B4) Finally, if $`\|P_{\overline{v}}\|=2`$, the $`\omega `$-labels of the external fields are equal and $`b_{\overline{v}}=1`$, we use for the action of $``$ on the external fields the representation following from the equations (3.58) and (3.59), after writing for the fields $`D^{1,3}`$ and $`D^{1,4}`$ the analogous of the decomposition (3.47). The above procedure can be iterated, by decomposing the factors $`d_{j(v)}^{b(v)}`$ coming from the previous steps of the iteration along the spanning tree associated with the clusters $`L_v`$, up to the endpoints. The final result can be described in the following way. Let us call a zero each factor equal to one of the two terms in (3.76). Each zero produced by the action of $``$ on the vertex $`v`$ is distributed along a tree graph $`S_v`$ on the set $`x_v`$, obtained by putting together an anchored tree graph $`T_{\overline{v}}`$ for each non trivial vertex $`\overline{v}v`$ and adding a line for the couple of space-time points belonging to the set $`𝐱_{\overline{v}}`$ for each (not local) endpoint $`\overline{v}v`$ with $`h_{\overline{v}}=2`$ of type $`\lambda `$ or $`u`$. At the end we have many terms, which are characterized, for what concerns the zeros, by a tree graph $`T`$ on the set $`x_{v_0}`$ and not more than two zeros on each line $`lT`$; the very important fact that there are at most two zeros on each line follows from the considerations in item B) of §3.9. 3.10 The final result can be written in the following way: $$\begin{array}{cc}\hfill V^{(h)}(\tau ,𝐏,𝐫)& =\sqrt{Z_h}^{\|P_{v_0}\|}\underset{T𝐓}{}\underset{\alpha A_T}{}d𝐱_{v_0}W_{\tau ,𝐏,𝐫,T,\alpha }(𝐱_{v_0})\hfill \\ & \left\{\underset{fP_{v_0}}{}[\widehat{}_{j_\alpha (f)}^{q_\alpha (f)}\psi ]_{𝐱_\alpha (f),\omega (f)}^{(h)\sigma (f)}\right\},\hfill \end{array}$$ $`(3.80)`$ where $$\begin{array}{cc}& W_{\tau ,𝐏,𝐫,T,\alpha }(𝐱_{v_0})=h_\alpha (\stackrel{~}{𝐱}_{v_0})\left[\underset{v\text{not e.p.}}{}(Z_{h_v}/Z_{h_v1})^{\|P_v\|/2}\right]\hfill \\ & \left[\underset{i=1}{\overset{n}{}}d_{j_\alpha (v_i^{})}^{b_\alpha (v_i^{})}(𝐱_i,𝐲_i)K_{v_i^{}}^{h_i}(𝐱_{v_i^{}})\right]\{\underset{v\text{not e.p.}}{}\frac{1}{s_v!}dP_{T_v}(𝐭_v)\hfill \\ & detG_\alpha ^{h_v,T_v}(𝐭_v)\left[\underset{lT_v}{}\widehat{}_{j_\alpha (f_l^{})}^{q_\alpha (f_l^{})}\widehat{}_{j_\alpha (f_l^+)}^{q_\alpha (f_l^+)}[d_{j_\alpha (l)}^{b_\alpha (l)}(𝐱_l,𝐲_l)\overline{}_1^{m_l}g_{\omega _l^{},\omega _l^+}^{(h_v)}(𝐱_l𝐲_l)]\right]\},\hfill \end{array}$$ $`(3.81)`$ $`𝐓`$ is the set of the tree graphs on $`𝐱_{v_0}`$, obtained by putting together an anchored tree graph $`T_v`$ for each non trivial vertex $`v`$ and adding a line (which will be by definition the only element of $`T_v`$) for the couple of space-time points belonging to the set $`𝐱_v`$ for each (not local) endpoint $`v`$ with $`h_v=2`$ of type $`\lambda `$ or $`u`$; $`A_T`$ is a set of indices which allows to distinguish the different terms produced by the non trivial $``$ operations and the iterative decomposition of the zeros; $`v_1^{},\mathrm{},v_n^{}`$ are the endpoints of $`\tau `$, $`f_l^{}`$ and $`f_l^+`$ are the labels of the two fields forming the line $`l`$, “e.p.” is an abbreviation of “endpoint”. Moreover $`G_\alpha ^{h_v,T_v}(𝐭_v)`$ is obtained from the matrix $`G^{h_v,T_v}(𝐭_v)`$, associated with the vertex $`v`$ and $`T_v`$, see (3.74), by substituting $`G_{ij,i^{}j^{}}^{h_v,T_v}=t_{v,i,i^{}}\overline{}_1^{m(f_{ij}^{})+m(f_{i^{}j^{}}^+)}g_{\omega _l^{},\omega _l^+}^{(h_v)}(𝐱_{ij}𝐲_{i^{}j^{}})`$ with $$G_{\alpha ,ij,i^{}j^{}}^{h_v,T_v}=t_{v,i,i^{}}\widehat{}_{j_\alpha (f_{ij}^{})}^{q_\alpha (f_{ij}^{})}\widehat{}_{j_\alpha (f_{ij}^+)}^{q_\alpha (f_{ij}^+)}\overline{}_1^{m(f_{ij}^{})+m(f_{i^{}j^{}}^+)}g_{\omega _l^{},\omega _l^+}^{(h_v)}(𝐱_{ij}𝐲_{i^{}j^{}}).$$ $`(3.82)`$ Finally, $`\widehat{}_j^q`$, $`q=0,1,2,3`$, $`j=1,\mathrm{},m_q`$, is a family of operators, implicitly defined in the previous sections, which are dimensionally equivalent to derivatives of order $`q`$; for each $`\alpha A_T`$, there is an operator $`\widehat{}_{j_\alpha (f)}^{q_\alpha (f)}`$ associated with each $`fI_{v_0}`$. It would be very difficult to give a precise description of the various contributions to the sum over $`A_T`$, but fortunately we only need to know some very general properties, which easily follows from the discussion in the previous sections. 1) There is a constant $`C`$ such that, $`T𝐓_\tau `$, $`\|A_T\|C^n`$ and, $`\alpha A_T`$, $`\|h_a(\stackrel{~}{𝐱}_{v_0})\|C^n`$. 2) For any $`\alpha A_T`$, the following inequality is satisfied $$\left[\underset{fI_{v_0}}{}\gamma ^{h_\alpha (f)q_\alpha (f)}\right]\left[\underset{lT}{}\gamma ^{h_\alpha (l)b_\alpha (l)}\right]\underset{v\text{not e.p.}}{}\gamma ^{z(P_v)},$$ $`(3.83)`$ where $`h_\alpha (f)=h_{v_0}1`$ if $`fP_{v_0}`$, otherwise it is the scale of the vertex where the field with label $`f`$ is contracted; $`h_\alpha (l)=h_v`$, if $`lT_v`$ and $$z(P_v)=\{\begin{array}{cc}1\hfill & \text{if }\|P_v\|=4\text{,}\hfill \\ 1\hfill & \text{if }\|P_v\|=2\text{ and }_{fP_v}\omega (f)0,\hfill \\ 2\hfill & \text{if }\|P_v\|=2\text{ and }_{fP_v}\omega (f)=0,\hfill \\ 0\hfill & \text{otherwise.}\hfill \end{array}$$ $`(3.84)`$ 3.11 In order to prove (3.83), let us suppose first that there is no vertex with two external fields and equal $`\omega `$ indices; hence $`q_\alpha (f)1`$, $`fI_{v_0}`$, and $`b_\alpha (l)1`$, $`lT`$. Let us choose $`fI_{v_0}`$, such that $`q_\alpha (f)=1`$; by analyzing the procedure described in §3.8 and §3.9, one can easily see that there are three vertices $`v^{}<v\overline{v}`$ and a line $`lT_{\overline{v}}`$, such that (i) the field with label $`f`$ is affected by the action of $``$ on the vertex $`v`$; (ii) $`h_v^{}=h_\alpha (f)`$ and $`b_\alpha (l)=1`$; (iii) if $`v\stackrel{~}{v}<\overline{v}`$ and $`\stackrel{~}{l}T_{\stackrel{~}{v}}`$, then $`b_\alpha (\stackrel{~}{l})=0`$; (iv) if $`v^{}<\stackrel{~}{v}\overline{v}`$ and $`f\stackrel{~}{f}P_{\stackrel{~}{v}}`$, then $`q_\alpha (\stackrel{~}{f})=0`$. (ii) follows from the definition of $`h_\alpha (f)`$ and from the remark that the zero produced by the action of $``$ on $`v`$ is moved by the process of distribution of the zeros along $`T`$ in some vertex $`\overline{v}v`$. The property (iii) characterizes $`\overline{v}`$; in fact the procedure described in item B1) and B2) of §3.9 guarantees that no zero can be produced by the action of $``$ in the vertices between $`v`$ and $`\overline{v}`$, if the zero in $`\overline{v}`$ “originated” from the regularization in $`v`$. (iv) follows from the previous remark and from the fact that the action of $``$ is trivial in all the vertices between $`v^{}`$ and $`v`$, see §3.3. The previous considerations imply that we can associate each factor $`\gamma ^{h_\alpha (f)}`$ in the l.h.s. of (3.83) with a factor $`\gamma ^{b_\alpha (l)}`$, by forming disjoint pairs; with each pair we can associate two vertices $`v^{}`$ and $`\overline{v}`$ and the path on $`\tau `$ containing all the vertices $`v^{}<\stackrel{~}{v}\overline{v}`$. Since each vertex with four external fields or two external fields and different $`\omega `$ indices certainly belongs to one of these paths, the inequality (3.83) then follows from the trivial identity $$\gamma ^{(b_\alpha (l)h_\alpha (f))}=\gamma ^{(h_{\overline{v}}h_v^{})}=\underset{v^{}<\stackrel{~}{v}\overline{v}}{}\gamma ^1.$$ $`(3.85)`$ In order to complete the proof, we have now to consider also the possibility that there is some vertex with two external fields and equal $`\omega `$ indices, where the action of $``$ is non trivial. This means that there is some $`fI_{v_0}`$, such that $`q_\alpha (f)=2`$ or even (see B4) in §3.9) $`q_\alpha (f)=1`$, if there is a zero associated with a line of the spanning tree related with the vertex where $`f`$ is affected by the regularization. One can proceed essentially in the same way, but has to consider a few different situations, since the value of $`q_\alpha (f)`$ is not fixed and, if $`q_\alpha (f)=2`$, there are two zeros to associate with a single factor $`\gamma ^{2h_\alpha (f)}`$ in the l.h.s. of (3.83). We shall not give the details, which have essentially to formalize the claim that each order one derivative couples with a order one zero, so that the corresponding factors in the l.h.s. of (3.83) contribute a factor $`\gamma ^1`$ to all vertices between the vertex where the derivative takes its action and the vertex where the zero is “sitting”. Let us now introduce, given any set $`PI_{v_0}`$, the notation $$q_\alpha (P)=\underset{fP}{}q_\alpha (f),m(P)=\underset{fP}{}m(f).$$ $`(3.86)`$ Note that, by the remark at the end of §3.2, $`m(P_v)=0`$ for any $`vv_0`$ which is not an endpoint of type $`\delta _1`$ or $`\delta _2`$ and that also $`m(P_{v_0})=0`$ for all the terms in the r.h.s. of (3.80). We also define $$\begin{array}{cc}\hfill \|\stackrel{}{v}_h\|& =\{\begin{array}{cc}sup\{\|\lambda \|,\|\nu \|\},\hfill & \text{if }h=+1\text{,}\hfill \\ sup\{\|\lambda _h\|,\|\delta _h\|,\|\nu _h\|\},\hfill & \text{if }h=0\text{.}\hfill \end{array}\hfill \\ \hfill \epsilon _h& =\underset{h^{}>h}{sup}\|\stackrel{}{v}_h^{}\|.\hfill \end{array}$$ $`(3.87)`$ Moreover, we suppose that the condition (2.117) is satisfied, so that $`h^{}0`$. We shall prove the following theorem. 3.12 Theorem. Let $`h>h^{}0`$, with $`h^{}`$ defined by (2.116). If the bounds (2.98) are satisfied and, for some constants $`c_1`$, $$\underset{h^{}>h}{sup}\left\|\frac{Z_h^{}}{Z_{h^{}1}}\right\|e^{c_1\epsilon _h^2},\underset{h^{}>h}{sup}\left\|\frac{\sigma _h^{}}{\sigma _{h^{}1}}\right\|e^{c_1\epsilon _h},$$ $`(3.88)`$ there exists a constant $`\overline{\epsilon }`$ (depending on $`c_1`$) such that, if $`\epsilon _h\overline{\epsilon }`$, then, for a suitable constant $`c_0`$, independent of $`c_1`$, as well as of $`u`$, $`L`$ and $`\beta `$, $$\begin{array}{cc}\hfill \underset{\tau 𝒯_{h,n}}{}& \underset{\genfrac{}{}{0pt}{}{𝐏}{\|P_{v_0}\|=2m}}{}\underset{𝐫}{}\underset{T𝐓}{}\underset{\genfrac{}{}{0pt}{}{\alpha A_T}{q_\alpha (P_{v_0})=k}}{}𝑑𝐱_{v_0}\|W_{\tau ,𝐏,𝐫,T,\alpha }(𝐱_{v_0})\|\hfill \\ & L\beta \gamma ^{hD_k(P_{v_0})}(c_0\epsilon _h)^n,\hfill \end{array}$$ $`(3.89)`$ where $$D_k(P_{v_0})=2+m+k.$$ $`(3.90)`$ Moreover $$\underset{\tau 𝒯_{h,n}}{}\left[\|n_h(\tau )\|+\|z_h(\tau )\|+\|a_h(\tau )\|+\|l_h(\tau )\|\right](c_0\epsilon _h)^n,$$ $`(3.91)`$ $$\underset{\tau 𝒯_{h,n}}{}\|s_h(\tau )\|\|\sigma _h\|(c_0\epsilon _h)^n,$$ $`(3.92)`$ $$\underset{\tau 𝒯_{h,n}}{}\|\stackrel{~}{E}_{h+1}(\tau )\|\gamma ^{2h}(c_0\epsilon _h)^n.$$ $`(3.93)`$ 3.13 An important role in the proof of Theorem 3.12 plays the estimation of $`detG_\alpha ^{h_v,T_v}(𝐭_v)`$, that we shall now discuss, by referring to §3.8 and §3.10 for the notation. From now on $`C`$ will denote a generic constant independent of $`u`$, $`L`$ and $`\beta `$. Given a vertex $`v`$ which is not an endpoint and an anchored tree graph $`T_v`$ (empty, if $`v`$ is trivial), we consider the set of internal fields which do not belong to the any line of $`T_v`$ and the corresponding sets $`\stackrel{~}{P}^{\sigma ,\omega }`$ of field labels with $`\sigma (f)=\sigma `$ and $`\omega (f)=\omega `$. The sets $`_\omega \stackrel{~}{P}^{,\omega }`$ and $`_\omega \stackrel{~}{P}^{+,\omega }`$ label the rows and the columns, respectively, of the matrix $`G_\alpha ^{h_v,T_v}(𝐭_v)`$, hence they contain the same number of elements; however, $`\|\stackrel{~}{P}^{,\omega }\|`$ can be different from $`\|\stackrel{~}{P}^{+,\omega }\|`$, if $`h0`$. We introduce an integer $`\rho (T_v)`$, that we put equal to $`1`$, if $`\|\stackrel{~}{P}^{,\omega }\|\|\stackrel{~}{P}^{+,\omega }\|`$, equal to $`0`$ otherwise. We want to prove that $$\begin{array}{cc}& \|detG_\alpha ^{h_v,T_v}(𝐭_v)\|\left(\frac{\|\sigma _{h_v}\|}{\gamma ^{h_v}}\right)^{\rho (T_v)}C^{_{i=1}^{s_v}\|P_{v_i}\|\|P_v\|2(s_v1)}\hfill \\ & \gamma ^{\frac{h_v}{2}\left(_{i=1}^{s_v}\|P_{v_i}\|\|P_v\|2(s_v1)\right)}\gamma ^{h_v_{i=1}^{s_v}\left[q_\alpha (P_{v_i}\backslash Q_{v_i})+m(P_{v_i}\backslash Q_{v_i})\right]}\hfill \\ & \gamma ^{h_v_{lT_v}\left[q_\alpha (f_l^+)+q_\alpha (f_l^{})+m(f_l^+)+m(f_l^{})\right]}.\hfill \end{array}$$ $`(3.94)`$ In order to prove this inequality, we shall suppose, for simplicity, that all the operators $`\widehat{}_{j(f)}^{q(f)}`$ and $`\widehat{}_1^{m(f)}`$ acting on the fields with field label $`f_{\sigma ,\omega }\stackrel{~}{P}^{\sigma ,\omega }`$ are equal to the identity. It is very easy to modify the following argument, in order to prove that each operator $`\widehat{}_{j(f)}^{q(f)}`$ or $`\widehat{}_1^{m(f)}`$ gives a contribution to the bound proportional to $`\gamma ^{h_vq(f)}`$ or $`\gamma ^{h_vm(f)}`$, so proving (3.94) in the general case. The proof is based on the well known Gram-Hadamard inequality, stating that, if $`M`$ is a square matrix with elements $`M_{ij}`$ of the form $`M_{ij}=<A_i,B_j>`$, where $`A_i`$, $`B_j`$ are vectors in a Hilbert space with scalar product $`<,>`$, then $$\|detM\|\underset{i}{}A_iB_i.$$ $`(3.95)`$ where $`\|\|\|\|`$ is the norm induced by the scalar product. Let $`=\text{}^s_0`$, where $`_0`$ is the Hilbert space of complex four dimensional vectors $`F(𝐤^{})=(F_1(𝐤^{}),\mathrm{},F_4(𝐤^{})`$), $`F_i(𝐤^{})`$ being a function on the set $`𝒟_{L,\beta }^{}`$, with scalar product $$<F,G>=\underset{i=1}{\overset{4}{}}\frac{1}{\beta L}\underset{𝐤^{}}{}F_i^{}(𝐤^{})G_i(𝐤^{}).$$ $`(3.96)`$ If $`h_v0`$, it is easy to verify that $$G_{ij,i^{}j^{}}^{h_v,T_v}=t_{i,i^{}}g_{\omega _l^{},\omega _l^+}^{(h_v)}(𝐱_{ij}𝐲_{i^{}j^{}})=<𝐮_iA_{𝐱(f_{ij}^{}),\omega (f_{ij}^{})}^{(h_v)},𝐮_i^{}B_{𝐱(f_{i^{}j^{}}^+),\omega (f_{i^{}j^{}}^+)}^{(h_v)}>,$$ $`(3.97)`$ where $`𝐮_i\text{}^s`$, $`i=1,\mathrm{},s`$, are the vectors such that $`t_{i,i^{}}=𝐮_i𝐮_i^{}`$, and $$\begin{array}{cc}\hfill A_{𝐱,\omega }^{(h)}(𝐤^{})& =e^{i𝐤^{}𝐱}\frac{\sqrt{\stackrel{~}{f}_h(𝐤^{})}}{\sqrt{A_h(𝐤^{})}}\{\begin{array}{cc}(ik_0+E(k^{}),0,i\sigma _{h1}(𝐤^{}),0),\hfill & \text{if }\omega =+1\text{,}\hfill \\ (0,i\sigma _{h1}(𝐤^{}),0,\sigma _{h1}),\hfill & \text{if }\omega =1\text{,}\hfill \end{array}\hfill \\ \hfill B_{𝐱,\omega }^{(h)}& =e^{i𝐤^{}𝐲}\frac{\sqrt{\stackrel{~}{f}_h(𝐤^{})}}{\sqrt{A_h(𝐤^{})}}\{\begin{array}{cc}(1,1,0,0),\hfill & \text{if }\omega =+1\text{,}\hfill \\ (0,0,1,(ik_0E(k^{}))/\sigma _{h1}),\hfill & \text{if }\omega =1\text{.}\hfill \end{array}\hfill \end{array}$$ $`(3.98)`$ Let us now define $`n_+=\|\stackrel{~}{P}^{,+}\|`$, $`m_+=\|\stackrel{~}{P}^{+,+}\|`$, $`m=\|\stackrel{~}{P}^{,+}\|+\|\stackrel{~}{P}^,\|=\|\stackrel{~}{P}^{+,+}\|+\|\stackrel{~}{P}^{+,}\|`$; by using (3.95) and (3.98), it is easy to see, by proceeding as in §2.7, that, if the conditions (2.98) hold, $$\|detG_\alpha ^{h_v,T_v}(𝐭_v)\|C^m\gamma ^{h_vn_+}\|\sigma _{h_v}\|^{mn_+}\left(\frac{\gamma ^{h_v}}{\|\sigma _{h_v}\|}\right)^{mm_+}=C^m\gamma ^{h_vm}\left(\frac{\|\sigma _{h_v}\|}{\gamma ^{h_v}}\right)^{m_+n_+}.$$ $`(3.99)`$ Since $`2m=_{i=1}^{s_v}\|P_{v_i}\|\|P_v\|2(s_v1)`$ and $`_{i=1}^{s_v}q_\alpha (P_{v_i}/Q_{v_i})_{lT_v}[q_\alpha (f_l^+)+q_\alpha (f_l^{})]=0`$, we get the inequality (3.94), if $`m_+n_+`$, by using (2.116). The case $`m_+<n_+`$ can be treated in a similar way, by exchanging the definitions of $`A_{𝐱,\omega }^{(h)}(𝐤^{})`$ and $`B_{𝐱,\omega }^{(h)}(𝐤^{})`$. 3.14 Proof of Theorem 3.12. By using (3.81) and (3.94) we get $$\begin{array}{cc}& d𝐱_{v_0}\|W_{\tau ,𝐏,𝐫,T,\alpha }(𝐱_{v_0})\|C^nJ_{\tau ,𝐏,𝐫,T,\alpha }\underset{v\text{not e.p.}}{}\{(Z_{h_v}/Z_{h_v1})^{\|P_v\|/2}\hfill \\ & C^{_{i=1}^{s_v}\|P_{v_i}\|\|P_v\|2(s_v1)}\left(\frac{\|\sigma _{h_v}\|}{\gamma ^{h_v}}\right)^{\rho (T_v)}\gamma ^{\frac{h_v}{2}\left(_{i=1}^{s_v}\|P_{v_i}\|\|P_v\|2(s_v1)\right)}\hfill \\ & \gamma ^{h_v_{i=1}^{s_v}\left[q_\alpha (P_{v_i}\backslash Q_{v_i})+m(P_{v_i}\backslash Q_{v_i})\right]}\gamma ^{h_v_{lT_v}\left[q_\alpha (f_l^+)+q_\alpha (f_l^{})+m(f_l^+)+m(f_l^{})\right]}\},\hfill \end{array}$$ $`(3.100)`$ where $$\begin{array}{cc}\hfill J_{\tau ,𝐏,𝐫,T,\alpha }& =d𝐱_{v_0}\|\left[\underset{i=1}{\overset{n}{}}d_{j_\alpha (v_i^{})}^{b_\alpha (v_i^{})}(𝐱_i,𝐲_i)K_{v_i^{}}^{h_i}(𝐱_{v_i^{}})\right]\hfill \\ & \left\{\underset{v\text{not e.p.}}{}\frac{1}{s_v!}\left[\underset{lT_v}{}\widehat{}_{j_\alpha (f_l^{})}^{q_\alpha (f_l^{})}\widehat{}_{j_\alpha (f_l^+)}^{q_\alpha (f_l^+)}[d_{j_\alpha (l)}^{b_\alpha (l)}(𝐱_l,𝐲_l)\widehat{}_1^{m_l}g_{\omega _l^{},\omega _l^+}^{(h_v)}(𝐱_l𝐲_l)]\right]\right\}\|.\hfill \end{array}$$ $`(3.101)`$ In §$`\mathrm{}`$3.15 we will prove that $$\begin{array}{cc}\hfill J_{\tau ,𝐏,𝐫,T,\alpha }& C^nL\beta (\epsilon _h)^n\underset{v\text{not e.p.}}{}[\frac{1}{s_v!}C^{2(s_v1)}\gamma ^{h_vn_\nu (v)}\left(\underset{l\overline{T}_v}{}\right\|\frac{\sigma _{h_v}}{\gamma ^{h_v}}\left\|\right)\hfill \\ & \gamma ^{h_v_{lT_v}b_\alpha (l)}\gamma ^{h_v(s_v1)}\gamma ^{h_v_{lT_v}\left[q_\alpha (f_l^+)+q_\alpha (f_l^{})+m(f_l^+)+m(f_l^{})\right]}],\hfill \end{array}$$ $`(3.102)`$ where $`n_\nu (v)`$ is the number of vertices of type $`\nu `$ with scale $`h_v+1`$ and $`\overline{T}_v`$ is the subset of the lines of $`T_v`$ corresponding to non diagonal propagators, that is propagators with different $`\omega `$ indices. It is easy to see that $$\underset{v\text{not e.p.}}{}h_v\underset{i=1}{\overset{s_v}{}}q_\alpha (P_{v_i}\backslash Q_{v_i})+hq_\alpha (P_{v_0})=\underset{fI_{v_0}}{}h_\alpha (f)q_\alpha (f)$$ $`(3.103)`$ and, by using also the remark after (3.86), that $$\begin{array}{cc}& \underset{\overline{v}v}{}\left\{\frac{1}{2}\left(\underset{i=1}{\overset{s_{\overline{v}}}{}}\|P_{\overline{v}_i}\|\|P_{\overline{v}}\|\right)2(s_{\overline{v}1})+n_\nu (\overline{v})+\underset{i=1}{\overset{s_{\overline{v}}}{}}m(P_{\overline{v}_i}\backslash Q_{\overline{v}_i})\right\}=\hfill \\ & =\frac{1}{2}(\|I_v\|\|P_v\|)+m(I_v\backslash P_v)+\underset{\overline{v}v}{}n_\nu (\overline{v})2(n_v1)=\frac{1}{2}\|P_v\|+2.\hfill \end{array}$$ $`(3.104)`$ By inserting (3.102) in (3.100) and using (3.83), (3.103), (3.104), we find $$\begin{array}{cc}& d𝐱_{v_0}\|W_{\tau ,𝐏,𝐫,T,\alpha }(𝐱_{v_0})\|C^nL\beta \epsilon _h^n\gamma ^{hD_k(P_{v_0})}\underset{vV_2}{}\frac{\|\sigma _{h_v}\|}{\gamma ^{h_v}}\hfill \\ & \underset{v\text{not e.p.}}{}\left\{\frac{1}{s_v!}C^{_{i=1}^{s_v}\|P_{v_i}\|\|P_v\|}(Z_{h_v}/Z_{h_v1})^{\|P_v\|/2}\gamma ^{[2+\frac{\|P_v\|}{2}+z(P_v)]}\right\},\hfill \end{array}$$ $`(3.105)`$ where $`V_2`$ is the set of vertices, which are not endpoints, such that $`\rho (T_v)+\|\stackrel{~}{T}_v\|>0`$, while the vertices $`\overline{v}>v`$ do not enjoy this property. Let us now consider a vertex $`v`$, which is not an endpoint, such that $`\|P_v\|=2`$ and $`_{fP_v}\omega (f)0`$. We want to show that there is a vertex $`\overline{v}v`$, such that $`\overline{v}V_2`$. In order to prove this claim, we note that, if $`v^{}`$ is an endpoint, then $`_{fP_v^{}}\sigma (f)\omega (f)=0`$, while $`_{fP_v}\sigma (f)\omega (f)0`$. Since all diagonal propagators join two fields with equal $`\omega `$ indices and opposite $`\sigma `$ indices, given any Feynman graph connecting the endpoints of the cluster $`L_v`$, at least one of its lines has to be a non diagonal propagator, so that at least one of the vertices $`\overline{v}v`$ must belong to $`V_2`$. Moreover, if $`vV_2`$, $$\frac{\|\sigma _{h_v}\|}{\gamma ^{h_v}}=\frac{\|\sigma _h\|}{\gamma ^h}\frac{\|\sigma _{h_v}\|}{\|\sigma _h\|}\gamma ^{hh_v}\frac{\|\sigma _h\|}{\gamma ^h}\gamma ^{(hh_v)(1c_1\epsilon _h)}C\gamma ^{(hh_v)(1/2)},$$ $`(3.106)`$ if $`\epsilon _h\overline{\epsilon }`$ and $`\overline{\epsilon }1/(2c_1)`$. We have used the second inequality in (3.88) and the definition (2.116), implying that $`\|\sigma _h\|\frac{a_0}{4\gamma }\gamma ^h`$, if $`hh^{}`$. It follows that $$\underset{vV_2}{}\frac{\|\sigma _{h_v}\|}{\gamma ^{h_v}}C^n\underset{v\text{not e.p.}}{}\gamma ^{\frac{1}{2}\stackrel{~}{z}(P_v)},$$ $`(3.107)`$ where $$\stackrel{~}{z}(P_v)=\{\begin{array}{cc}1\hfill & \text{if }\|P_v\|=2\text{ and }_{fP_v}\omega (f)0,\hfill \\ 0\hfill & \text{otherwise,}\hfill \end{array}$$ $`(3.108)`$ so that $$2+\frac{\|P_v\|}{2}+z(P_v)+\frac{\stackrel{~}{z}(P_v)}{2}\frac{1}{2},v\text{not e.p.}.$$ $`(3.109)`$ Hence (3.105) can be changed in $$\begin{array}{cc}& d𝐱_{v_0}\|W_{\tau ,𝐏,𝐫,T,\alpha }(𝐱_{v_0})\|C^nL\beta \epsilon _h^n\gamma ^{hD_k(P_{v_0})}\hfill \\ & \underset{v\text{not e.p.}}{}\left\{\frac{1}{s_v!}C^{_{i=1}^{s_v}\|P_{v_i}\|\|P_v\|}(Z_{h_v}/Z_{h_v1})^{\|P_v\|/2}\gamma ^{[2+\frac{\|P_v\|}{2}+z(P_v)+\frac{\stackrel{~}{z}(P_v)}{2}]}\right\},\hfill \end{array}$$ $`(3.110)`$ In order to complete the proof of the bound (3.89), we have to perform the sums in the r.h.s. of (3.89). The number of unlabeled trees is $`4^n`$; fixed an unlabeled tree, the number of terms in the sum over the various labels of the tree is bounded by $`C^n`$, except the sums over the scale labels and the sets $`𝐏`$. The number of addenda in the sums over $`\alpha `$ and $`𝐫`$ is again bounded by $`C^n`$, since the action of $``$ can be non trivial at most two times between two consecutive non trivial vertices (see §3.3) and the number of non trivial vertices is of order $`n`$. Regarding the sum over $`T`$, it is empty if $`s_v=1`$. If $`s_v>1`$ and $`N_{v_i}\|P_{v_i}\|\|Q_{v_i}\|`$, the number of anchored trees with $`d_i`$ lines branching from the vertex $`v_i`$ can be bounded, by using Caley’s formula, by $$\frac{(s_v2)!}{(d_11)!\mathrm{}(d_{s_v}1)!}N_{v_1}^{d_1}\mathrm{}N_{v_{s_v}}^{d_{s_v}};$$ hence the number of addenda in $`_{T𝐓}`$ is bounded by $`_{v\text{not }\text{e.p.}}s_v!C^{_{i=1}^{s_v}\|P_{v_i}\|\|P_v\|}`$. In order to bound the sums over the scale labels and $`𝐏`$ we first use the inequality, following from (3.109) and the first inequality in (3.88), if $`c_1\epsilon _h^21/16`$, $$\begin{array}{cc}& \underset{v\text{not e.p.}}{}(Z_{h_v}/Z_{h_v1})^{\|P_v\|/2}\gamma ^{\frac{1}{2}[2+\frac{\|P_v\|}{2}+z(P_v)+\frac{\stackrel{~}{z}(P_v)}{2}]}]\hfill \\ & [\underset{\stackrel{~}{v}}{}\gamma ^{\frac{1}{40}(h_{\stackrel{~}{v}}h_{\stackrel{~}{v}^{}})}][\underset{v\text{not e.p.}}{}\gamma ^{\frac{\|P_v\|}{40}}],\hfill \end{array}$$ $`(3.111)`$ where $`\stackrel{~}{v}`$ are the non trivial vertices, and $`\stackrel{~}{v}^{}`$ is the non trivial vertex immediately preceding $`\stackrel{~}{v}`$ or the root. The factors $`\gamma ^{\frac{1}{40}(h_{\stackrel{~}{v}}h_{\stackrel{~}{v}^{}})}`$ in the r.h.s. of (3.111) allow to bound the sums over the scale labels by $`C^n`$. Finally the sum over $`𝐏`$ can be bound by using the following combinatorial inequality, trivial for $`\gamma `$ large enough, but valid for any $`\gamma >1`$ (see \[BGPS\], $3). Let $`\{p_v,v\tau \}`$ a set of integers such that $`p_v_{i=1}^{s_v}p_{v_i}`$ for all $`v\tau `$ which are not endpoints; then $$\underset{v\text{not e.p.}}{}\underset{p_v}{}\gamma ^{\frac{p_v}{40}}C^n.$$ $`(3.112)`$ It follows that $$\underset{\genfrac{}{}{0pt}{}{𝐏}{\|P_{v_0}\|=2m}}{}\underset{v\text{not e.p.}}{}\gamma ^{\frac{\|P_v\|}{40}}\underset{v\text{not e.p.}}{}\underset{p_v}{}\gamma ^{\frac{p_v}{40}}C^n.$$ $`(3.113)`$ The proof of the bounds (3.91) and (3.93) is very similar. For the terms contributing to $`n_h`$ one gets a bound like (3.89), with $`m=1`$ and $`k=0`$, but the factor $`\gamma ^{hD_k(P_{v_0})}=\gamma ^h`$ is compensated by the factor $`\gamma ^h`$ appearing in the definition of $`n_h(\tau )`$, see (3.71). For the terms contributing to $`z_h`$ and $`a_h`$ $`D_k(P_{v_0})=0`$ ($`m=k=1`$), as well as for those contributing to $`l_h`$ ($`m=2`$, $`k=0`$). Finally, for the terms contributing to $`\stackrel{~}{E}_{h+1}`$, $`D_k(P_{v_0})=2`$. For the terms contributing to $`s_h`$, $`D_k(P_{v_0})=1`$, but each term has also at least one small factor $`\|\sigma _h\|\gamma ^h`$ in its bound, since $`\|V_2\|1`$, see (3.106); so we get the bound (3.92). 3.15 Proof of (3.102). We shall refer to the definitions and the discussion in §3.7 and §3.9. Let us consider the factor in the r.h.s. of (3.101) associated with the line $`lT_v`$ and let us suppose that $`𝐱_l𝐱^{(i)}`$, $`𝐲_l𝐱^{(i^{})}`$. By using (3.47), (3.53) and the similar expressions for the other difference fields produced by the regularization, we can write $$\begin{array}{cc}& \widehat{}_{j_\alpha (f_l^{})}^{q_\alpha (f_l^{})}\widehat{}_{j_\alpha (f_l^+)}^{q_\alpha (f_l^+)}[d_{j_\alpha (l)}^{b_\alpha (l)}(𝐱_l,𝐲_l)\overline{}_1^{m_l}g_{\omega _l^{},\omega _l^+}^{(h_v)}(𝐱_l𝐲_l)]=\hfill \\ & =_0^1𝑑t_l_0^1𝑑s_l\stackrel{~}{}_{j_\alpha (f_l^{})}^{q_\alpha (f_l^{})}\stackrel{~}{}_{j_\alpha (f_l^+)}^{q_\alpha (f_l^+)}[d_{j_\alpha (l)}^{b_\alpha (l)}(𝐱_l^{}(t_l),𝐲_l^{}(s_l))\overline{}_1^{m_l}g_{\omega _l^{},\omega _l^+}^{(h_v)}(𝐱_l^{}(t_l)𝐲_l^{}(s_l))],\hfill \end{array}$$ $`(3.114)`$ where, depending on $`\alpha `$, there are essentially two different possibilities for the operators $`\stackrel{~}{}_{j_\alpha }^{q_\alpha }`$ and the space-time points $`𝐱_l^{}(t_l)`$, $`𝐲_l^{}(s_l)`$. Let us consider, for example, $`f_l^{}`$; then the first possibility is that $`\stackrel{~}{}_{j_\alpha }^{q_\alpha }`$ is a derivative of order $`q_\alpha `$ and $$𝐱_l^{}(t_l)=\stackrel{~}{𝐱}_l+t_l(\overline{𝐱}_l\stackrel{~}{𝐱}_l),\text{for some}\stackrel{~}{𝐱}_l𝐱^{(i)},$$ $`(3.115)`$ $`\overline{𝐱}_l`$ being defined in terms of $`𝐱_l`$ as $`\overline{y}`$ is defined in terms of $`y`$ in §3.5 (that is $`\overline{𝐱}_l`$ and $`𝐱_l`$ are equivalent representation of the same point on the space-time torus). The second possibility is that $`\stackrel{~}{}_{j_\alpha }^{q_\alpha }`$ is a local operator of the form $`L^{n_1}\beta ^{n_2}\overline{}_1^{n_3}_0^{n_4}`$, with $`q_\alpha _{i=1}^4n_iq_\alpha +1`$, and $`𝐱_l^{}(t_l)=\stackrel{~}{𝐱}_l𝐱^{(i)}`$. Note that, by (2.40), $`L^{n_1}\beta ^{n_2}\gamma ^{h_{L,\beta }(n_1+n_2)}\gamma ^{h_v(n_1+n_2)}`$. By proceeding as in the proof of lemma (2.6) and using (2.105) it is very easy to show that, for any $`N>1`$, $$\begin{array}{cc}& \left\|\stackrel{~}{}_{j_\alpha (f_l^{})}^{q_\alpha (f_l^{})}\stackrel{~}{}_{j_\alpha (f_l^+)}^{q_\alpha (f_l^+)}[d_{j_\alpha (l)}^{b_\alpha (l)}(𝐱_l^{}(t_l),𝐲_l^{}(s_l))\overline{}_1^{m_l}g_{\omega _l^{},\omega _l^+}^{(h_v)}(𝐱_l^{}(t_l)𝐲_l^{}(s_l))]\right\|\hfill \\ & C\frac{\gamma ^{h_v[1+q_\alpha (f_l^+)+q_\alpha (f_l^{})+m(f_l^{})+m(f_l^+)b_\alpha (l)]}}{1+[\gamma ^{h_v}\|𝐝(𝐱_l^{}(t_l)𝐲_l^{}(s_l))\|]^N}\left(\frac{\|\sigma _{h_v}\|}{\gamma ^{h_v}}\right)^{\rho _l},\hfill \end{array}$$ $`(3.116)`$ where $`𝐝(𝐱)`$ is defined in (2.97) and $`\rho _l=1`$ if $`\omega (f_l^{})\omega (f_l^+)`$, $`\rho _l=0`$ otherwise. We used here the fact that, if $`h_v=+1`$, then $`q_\alpha (f_l^{})=q_\alpha (f_l^+)=0`$, which allows to avoid the problems connected with the singularity of the time derivatives of the scale $`1`$ propagator at $`x_{l,0}^{}(t_l)y_{l,0}^{}(s_l)=0`$. Let us now consider the contribution of the endpoints to the r.h.s. of (3.101) and recall (see §3.10) that $`T_{v_i^{}}`$ is empty, if $`\|𝐱_{v_i^{}}\|=1`$, hence $`b_\alpha (v_i^{})=0`$, while, if $`𝐱_{v_i^{}}=(𝐱_i,𝐲_i)`$, $`T_{v_i^{}}`$ contains the line $`l_i`$ connecting $`𝐱_i`$ with $`𝐲_i`$ and $`h_{v_i^{}}=2`$. By using (3.33) and (3.39), we get, if $`h_ih_{v_i^{}}`$ and $`S_\nu \{i:v_i^{}\text{is of type}\nu \}`$, $$\begin{array}{cc}& \left\|\left[\underset{i=1}{\overset{n}{}}d_{j_\alpha (v_i^{})}^{b_\alpha (v_i^{})}(𝐱_i,𝐲_i)K_{v_i^{}}^{h_i}(𝐱_{v_i^{}})\right]\right\|\hfill \\ & C^n\epsilon _h^n\underset{i:\|𝐱_{v_i^{}}\|=2}{}\frac{1}{[1+\|𝐝(𝐱_i𝐲_i)\|]^N}\underset{iS_\nu }{}\gamma ^{(h_i1)}.\hfill \end{array}$$ $`(3.117)`$ Let us now remark that, after the insertion of the bounds (3.116) and (3.117) in the r.h.s. of (3.101), by possibly changing the constant $`C`$, we can substitute $`𝑑𝐱_{v_0}`$, which is there a shorthand for $`_{𝐱𝐱_{v_0}}_{x\mathrm{\Lambda }}𝑑x_0`$, with the real integral over $`(\text{𝕋}_{L,\beta })^{\|𝐱_{v_0}\|}`$, where $`\text{𝕋}_{L,\beta }`$ is the space-time torus $`[L/2,L/2]\times [\beta /2,\beta /2]`$. Moreover, equation (3.115) can be thought, and we shall do that, as defining an interval on $`\text{𝕋}_{L,\beta }`$, when $`t_l`$ spans the interval $`[0,1]`$; this is possible thanks to the introduction of the partition (3.42) in §3.5. Hence, in order to complete the proof of (3.102), we have to show that, fixed a point $`\overline{𝐱}𝐱_{v_0}`$, the interpolation parameters associated with the regularization operations and an integer $`N3`$, $$_\mathrm{\Xi }d(𝐱_{v_0}\backslash \overline{𝐱})\underset{v\tau }{}\underset{lT_v}{}\frac{1}{1+[\gamma ^{h_v}\|𝐝(𝐱_l^{}(t_l)𝐲_l^{}(s_l))\|]^N}\underset{v\tau }{}C\gamma ^{h_v(s_v1)},$$ $`(3.118)`$ where $`\mathrm{\Xi }`$ denotes the subset of $`(\text{𝕋}_{L,\beta })^{\|𝐱_{v_0}\backslash \overline{𝐱}\|}`$ satisfying all the constraints associated with the interpolated points of the form (3.115). Let us call $`\stackrel{~}{T}=_v\stackrel{~}{T}_v`$, where $`\stackrel{~}{T}_v`$ is the set of lines connecting $`𝐱_l^{}(t_l)`$ with $`𝐲_l^{}(s_l)`$, for any $`lT_v`$. $`\stackrel{~}{T}`$ is not a tree in general; however, for any $`v`$, $`\stackrel{~}{T}_v`$ is still an anchored tree graph between the clusters of points $`𝐱^{(i)}`$, $`i=1,\mathrm{},s_v`$. Hence, the proof of (3.118) becomes trivial, if we can show that $$d(𝐱_{v_0}\backslash \overline{𝐱})=\underset{l\stackrel{~}{T}}{}d𝐫_l,$$ $`(3.119)`$ where $`𝐫_l=𝐱_l^{}(t_l)𝐲_l^{}(s_l)`$. In order to prove (3.119), we can proceed, for example, as in \[BM1\]. Let us consider first a vertex $`v`$ with $`\|T_v\|>0`$, which is maximal with respect to the tree order; hence either $`v`$ is a non local endpoint with $`h_v=2`$ or it is a non trivial vertex with no vertex $`v^{}`$ with $`\|T_v^{}\|>0`$ following it. In this case $`\stackrel{~}{T}_v=T_v`$, that is no line depends on the interpolation parameters, and $`\stackrel{~}{T}_v`$ is a tree on the set $`𝐱_v`$, so that we get immediately the identity $$d𝐱_v=d\overline{𝐱}^{(v)}\underset{l\stackrel{~}{T}_v}{}d𝐫_l,$$ $`(3.120)`$ where $`\overline{𝐱}^{(v)}`$ is an arbitrary point of $`𝐱_v`$. If we use (3.120) for the family $`S_0`$ of all maximal vertices with $`\|T_v\|>0`$, we get $$d𝐱_{v_0}=\underset{vS_0}{}\left[d\overline{𝐱}^{(v)}\underset{l\stackrel{~}{T}_v}{}d𝐫_l\right].$$ $`(3.121)`$ Let us now consider a line $`\overline{l}\stackrel{~}{T}`$, which connects two clusters of points $`𝐱_{v_1}`$ and $`𝐱_{v_2}`$, with $`v_iS_0`$, $`i=1,2`$. By (3.115) $$𝐫_{\overline{l}}=𝐱_{\overline{l}}^{}(t_{\overline{l}})𝐲_l^{}(s_{\overline{l}})=t_{\overline{l}}𝐱_{\overline{l}}+(1t_{\overline{l}})\overline{𝐱}_{\overline{l}}𝐲_{\overline{l}}^{}(s_{\overline{l}}),$$ $`(3.122)`$ implying that $$\overline{𝐱}^{(v_1)}=𝐫_{\overline{l}}+\overline{𝐱}^{(v_1)}𝐫_{\overline{l}}=𝐫_{\overline{l}}+t_{\overline{l}}(\overline{𝐱}^{(v_1)}𝐱_{\overline{l}})+(1t_{\overline{l}})(\overline{𝐱}^{(v_1)}\overline{𝐱}_{\overline{l}})+𝐲_{\overline{l}}^{}(s_{\overline{l}}).$$ $`(3.123)`$ Since $`𝐲_{\overline{l}}^{}(s_{\overline{l}})`$ depends only on the variables $`𝐱_{v_2}`$ and $`(\overline{𝐱}^{(v_1)}𝐱_{\overline{l}})`$ and $`(\overline{𝐱}^{(v_1)}\overline{𝐱}_{\overline{l}})`$ both depend only on $`\{𝐫_l,l\stackrel{~}{T}_{v_1}\}`$, we get $$\underset{i=1}{\overset{2}{}}\left[d\overline{𝐱}^{(v_i)}\underset{l\stackrel{~}{T}_{v_i}}{}d𝐫_l\right]=d𝐫_{\overline{l}}d\overline{𝐱}^{(v_2)}\underset{i=1}{\overset{2}{}}\underset{l\stackrel{~}{T}_{v_i}}{}d𝐫_l.$$ $`(3.124)`$ By iterating this procedure, one gets (3.119). 3.16 As we have discussed in §2.13, it is not necessary to perform the scale decomposition of the Grassmanian integration up to the last scale $`h_{L,\beta }`$, but we can stop it to the scale $`h^{}`$, defined in (2.116). Hence, we redefine $`\stackrel{~}{E}_h^{}`$, so that $$e^{L\beta \stackrel{~}{E}_h^{}}=P_{Z_{h^{}1},\sigma _{h^{}1},C_h^{}}(d\psi ^{(h^{})})e^{\widehat{𝒱}^{(h^{})}(\sqrt{Z_{h^{}1}}\psi ^{(h^{})})},$$ $`(3.125)`$ implying that $$E_{L,\beta }=\underset{h=h^{}}{\overset{1}{}}[\stackrel{~}{E}_h+t_h].$$ $`(3.126)`$ Thanks to Lemma 2.12, we can proceed as in the proof of Theorem 3.12 to prove the following Theorem. 3.17 Theorem. There exists a constant $`\overline{\epsilon }`$ such that, if $`\epsilon _h^{}\overline{\epsilon }`$ and, for $`h=h^{}`$, (2.98) holds and the bounds (3.88) are satisfied, then $$\underset{\tau 𝒯_{h^{}1,n}}{}\|\stackrel{~}{E}_h^{}(\tau )\|\gamma ^{2h^{}}(C\epsilon _h^{})^n.$$ $`(3.127)`$ 3.18 Theorems 3.12 and 3.17, together with (3.126) and (2.118), imply that the expansion defining $`E_{L,\beta }`$ is convergent, uniformly in $`L,\beta `$. 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warning/0004/hep-ph0004265.html	ar5iv	text	# The Radiative Leptonic 𝐵_𝑐→𝜏⁢(𝜈_𝜏)̄⁢𝛾 Decay in Two Higgs Doublet Model ## 1 Introduction The observation of the $`B_c`$-meson by the CDF collaboration in the channel $`B_cJ/\psi \mathrm{}\nu `$, with ground state mass $`B_c=6.4\pm 0.39\pm 0.13`$ GeV, and lifetime $`\tau (B_c)=0.46_{0.16}^{+0.18}`$$`\pm 0.03`$ ps , has stimulated up the investigation of the properties of $`B_c`$-mesons theoretically, and experimentally on a new footing. The particular interest of this observation is related to the fact that, the meson ground-state with $`\overline{b}c(b\overline{c})`$ can decay only weakly, thus providing the rather unique opportunity of investigating weak decays in a heavy quarkonium system. Moreover, studying this meson could offer an unique probe to check the perturbative QCD predictions more precisely, and one can get essential new information about the confinement scale inside hadrons. The theoretical study of the pure-leptonic decays of $`B_c`$-meson, such as $`B_c\mathrm{}\overline{\nu _{\mathrm{}}}`$ ($`\mathrm{}=e,\mu ,\tau `$) can be used to determine the leptonic decay constant $`f_{B_c}`$ , as well as the fundamental parameters in the Standard Model (SM), such as the Cabibbo- Kobayashi-Maskawa (CKM) matrix elements which are poorly known at present. Nevertheless, the well-known ”helicity suppression” effect make an experimental difficulty in the measurement of purely leptonic decays of $`B_c`$. Although, the $`B_c\tau \overline{\nu _\tau }`$ channel is free of the helicity suppression, the observation of this decay is experimentally difficult due to the efficiency problem for detecting the $`\tau `$ lepton. Recently, the radiative leptonic $`B_c\mathrm{}\overline{\nu _{\mathrm{}}}\gamma `$ decays ($`\mathrm{}=e,\mu ,\tau `$), received considerable attention as a testing ground of SM and ”new physics”, where no helicity suppression exists any more \[3-5\]. Among various radiative leptonic decays, the $`B_c\tau \overline{\nu _\tau }\gamma `$ decay provokes special interest, since the SM predictions has been exploited to establish a bound on the branching ratio of the above mentioned decay of order $`10^5`$ , and therefore can be potentially measurable in the up coming LHC B-factories, where the number of $`B_c`$-mesons that will be produced are estimated to be $`2.0\times 10^{12}`$ . This will provide an alternative way to determine the decay constant $`f_{B_c}`$ and the CKM matrix elements. The decay $`B_c\tau \overline{\nu _\tau }\gamma `$ receive two types of contributions: internal bremsstrahlung (IB), and structure-dependent (SD) parts. The IB contributions are still helicity suppressed, while the SD ones contain the electromagnetic coupling constant $`\alpha `$ but they are free of helicity suppression. Therefore, the radiative decay rates of $`B_c\mathrm{}\overline{\nu _{\mathrm{}}}\gamma `$ could have an enhancement with respect to the purely leptonic modes of $`B_c\mathrm{}\overline{\nu _{\mathrm{}}}`$ due to the SD contributions, thus it enable to establish ”new physics” beyond the standard model. In this work, we will study the radiative leptonic $`B_c\tau \overline{\nu _\tau }\gamma `$ decay in the framework of the two-Higgs doublet model (2HDM) \[9-11\] at large $`tan\beta `$. The so-called Model I and Model II are considered, which are differ only in the couplings of the charged Higgs bosons to fermions. Subsequently, this paper is organized as follows: in section 2, the theoretical formalism relevance for the $`B_c\tau \overline{\nu _\tau }\gamma `$ decay in 2HDM is presented. Section 3, is devoted to the numerical analysis and the discussion of the results. ## 2 Formalism for the $`B_c\tau \overline{\nu _\tau }\gamma `$ decay In the Standard Model (SM), the decay $`B_c\tau \overline{\nu _\tau }\gamma `$ can be studied to a very good approximation in terms of four-fermion interactions. The effective Hamiltonian relevant to the process $`B_c\tau \overline{\nu _\tau }`$ is: $`H_{eff}`$ $`=`$ $`{\displaystyle \frac{G_F}{\sqrt{2}}}a_1V_{cb}\overline{c}\gamma _\mu (1\gamma _5)b\overline{\tau }\gamma _\mu (1\gamma _5)\nu _\tau ,`$ (1) where $`G_F`$ is the Fermi constant, $`V_{cb}`$ is the CKM mixing element, $`a_1`$ is a QCD corrected factor, which is equal $`a_11.13`$. However, in the next discussion we will put $`a_11`$ The emission of a real photon in leptonic decays of heavy mesons $`B_c\tau \overline{\nu _\tau }\gamma `$ can proceed via the two mechanisms mentioned in Sec. I. For the IB amplitude, the charged $`B_c`$-meson emits leptons via the axial-vector current, and the photon is radiated from the external charged particles. On the other hand the SD amplitude is governed by the vector and axial vector form factors, in which the photon is emitted from intermediate states. Gauge invariance leaves only two form factors $`f_{1,2}(p^2)`$ undetermined in the SD part. The possible diagrams for this two mechanisms are shown in Fig. 1. Following this framework, the general form of the gauge invariant amplitude corresponding to Fig. 1 can be written as : $`M(B_c\tau \overline{\nu _\tau }\gamma )=M_1+M_2,`$ (2) where $`M_1`$ and $`M_2`$ represent the contributions of ”inner bremsstrahlung” (IB), and ”structure-dependent” (SD) parts, given by: $`M_1`$ $`=`$ $`ie{\displaystyle \frac{G_F}{\sqrt{2}}}V_{cb}f_{B_c}m_\tau \epsilon _\alpha \overline{u}(p_1)\left\{{\displaystyle \frac{\gamma _\alpha \overline{)}q+2p_{1\alpha }}{2p_1.q}}{\displaystyle \frac{p_\alpha }{p.q}}\right\}\left(1\gamma _5\right)v(p_2),`$ (3) $`M_2`$ $`=`$ $`ie{\displaystyle \frac{G_F}{\sqrt{2}}}V_{cb}\epsilon _\alpha \{{\displaystyle \frac{if_1(p^2)}{m_{B_c}^2}}ϵ_{\alpha \beta \rho \gamma }p_\rho q_\gamma `$ (4) $`+`$ $`{\displaystyle \frac{f_2(p^2)}{m_{B_c}^2}}[p_\alpha q_\beta g_{\alpha \beta }p.q]\}\overline{u}(p_1)\gamma _\beta (1\gamma _5)v(p_2),`$ where $`\epsilon _\alpha `$ is the photon polarization vector, $`p_1`$, $`p_2`$, and q are the four momenta of $`\tau `$, $`\nu _\tau `$, and $`\gamma `$, respectively. $`f_{B_c}`$ is the $`B_c`$\- meson leptonic decay constant, $`f_{1,2}(p^2)`$ corresponding to parity conserving and a parity violating formfactors, $`p=P_{B_c}=p_1+p_2`$ is the momentum transfer to lepton pair. The necessary matrix elements related to the $`f_{B_c}`$, and to the hadron transition form factors $`f_{1,2}(p^2)`$ are defined as follows : $`0\overline{c}\gamma _\mu \gamma _5bB_c=if_{B_c}P_{B\mu },`$ (5) $`\gamma (q)\overline{c}\gamma _\alpha bB_c(p+q)=e{\displaystyle \frac{f_1(p^2)}{m_{B_c}^2}}ϵ_{\alpha \beta \rho \gamma }\epsilon ^\beta p^\rho q^\gamma ,`$ (6) $`\gamma (q)\overline{c}\gamma _\alpha \gamma _5bB_c(p+q)=ie{\displaystyle \frac{f_2(p^2)}{m_{B_c}^2}}\epsilon _\alpha [g_{\alpha \beta }(p.q)p_\alpha q_\beta ].`$ (7) We want now to consider the $`B_c\tau \overline{\nu }_\tau \gamma `$ decay in the context of a 2HDM with no flavor changing neutral currents (FCNC) allowed at the tree level, i.e. Model I and Model II. The interaction lagrangian of fermions with the charged Higgs fields in both models is given by : $`L`$ $`=`$ $`{\displaystyle \frac{g_W}{2\sqrt{2}M_W}}\{V_{ij}m_{u_i}X\overline{u_i}(1\gamma _5)d_j+V_{ij}m_{d_j}Y\overline{u_i}(1+\gamma _5)d_j`$ (8) $`+`$ $`m_{\mathrm{}}Z\overline{\nu }(1+\gamma _5)\mathrm{}\}H^\pm +h.c.,`$ where $`g_W`$ is the weak coupling constant, $`M_W`$ is the $`W`$\- boson mass, $`H^\pm `$ is the charged physical field, and $`V_{ij}`$ is the relevant elements of CKM matrix. In model I, $`X=cot\beta `$, $`Y=Z=cot\beta `$ and in model II, $`X=cot\beta `$, and $`Y=Z=tan\beta `$. The decay $`B_c\tau \overline{\nu }_\tau \gamma `$ in 2HDM proceeds through the same Feynman diagrams (which are displayed in Fig. 1) that mediate the process in SM, except the $`W`$-boson is replaced by the charged scalar Higgs boson $`H^\pm `$, i.e. $`(WH^\pm )`$. Fig. 1 The matrix element corresponding to the diagram $`(WH^\pm )`$ where the photon is radiated from $`\tau `$ lepton, (see fig.1) is: $`M_1^{2HDM}`$ $`=`$ $`ie{\displaystyle \frac{G_F}{\sqrt{2}}}V_{cb}f_{B_c}\epsilon _\alpha {\displaystyle \frac{m_\tau m_{B_c}^2}{M_H^2(m_b+m_c)}}(m_bZYm_cZX)\times `$ (9) $`\overline{u}(p_1)\left\{{\displaystyle \frac{2p_{1\alpha }+\gamma _\alpha \overline{)}q}{2p_1q}}\right\}(1\gamma _5)v(p_2),`$ where $`f_{B_c}`$ is the leptonic decay constant of $`B_c`$ meson, defined as: $`0\overline{c}\gamma _5bB_c(p+q)=if_{B_c}{\displaystyle \frac{m_{B_c}^2}{(m_b+m_c)}}.`$ (10) While the contribution of the structure dependent part to the $`B_c\tau \overline{\nu }_\tau \gamma `$ decay, i.e., when photon is radiated from initial quark lines, due to the charged Higgs exchange can be obtained by considering the following correlation function: $`M_\alpha ^{SD}`$ $`=`$ $`ie{\displaystyle \frac{G_F}{\sqrt{2}}}V_{cb}\epsilon _\alpha Z{\displaystyle \frac{m_\tau }{M_H^2}}{\displaystyle }d^4xe^{iqx}\times `$ (11) $`<0T\left\{\left[\overline{c}(0)\left(m_bY(1+\gamma _5)+m_cX(1\gamma _5)\right)b(0)\right]J_\alpha ^{el}(x)\right\}B_c>\times `$ $`\overline{u}(p_1)\gamma _\beta \left(1\gamma _5\right)v(p_2),`$ where $`J_\alpha ^{el}(x)`$ is the electromagnetic current for $`b`$ or $`c`$ quarks. The hadronic matrix elements involving the scalar and pseudoscalar currents in Eq.(11) are parameterized such that: $`{\displaystyle d^4xe^{iqx}<0T(\overline{c}(0)\gamma _5b(0)J_\alpha ^{el}(x))B_c(p+q)>}=f_{B_c}{\displaystyle \frac{m_{B_c}^2}{(m_b+m_c)}}{\displaystyle \frac{p_\alpha }{p.q}},`$ (12) $`{\displaystyle d^4xe^{iqx}<0T(\overline{c}(0)b(0)J_\alpha ^{el}(x))B_c(p+q)>}=0.`$ (13) The parameterization of the hadronic matrix elements given in Eqs. (12,13) are particularly well suited for our purposes since in the 2HDM, the vertex $`b(c)(1\gamma _5)`$ or $`b(c)(1+\gamma _5)`$, hence, the vector part of this correlator is zero, and the active part of this correlator is the axial-part given by: $`M_\alpha ^{A(SD)}`$ $`=`$ $`ie{\displaystyle \frac{G_F}{\sqrt{2}}}V_{cb}f_{B_c}\epsilon _\alpha {\displaystyle \frac{m_\tau m_{B_c}^2}{M_H^2}}{\displaystyle \frac{(m_bZYm_cZX)}{(m_b+m_c)}}\times `$ (14) $`\overline{u}(p_1){\displaystyle \frac{p_\alpha }{pq}}(1\gamma _5)v(p_2),`$ and the total matrix element in 2HDM becomes: $`M_2^{2HDM}(B_c\tau \overline{\nu }_\tau \gamma )`$ $`=`$ $`ie{\displaystyle \frac{G_F}{\sqrt{2}}}V_{cb}f_{B_c}\epsilon _\alpha {\displaystyle \frac{m_\tau m_{B_c}^2}{M_H^2}}{\displaystyle \frac{(m_bZYm_cZX)}{(m_b+m_c)}}\times `$ (15) $`\overline{u}(p_1)\left\{{\displaystyle \frac{\gamma _\alpha \overline{)}q+2p_{1\alpha }}{2p_1.q}}{\displaystyle \frac{p_\alpha }{p.q}}\right\}\left(1\gamma _5\right)v(p_2).`$ At this accuracy it is easy to check that the modified total amplitude for the radiative leptonic B-decays of $`B_c\tau \overline{\nu }_\tau \gamma `$ is gauge invariant: $`M_{(total)}(B_c\tau \overline{\nu _\tau }\gamma )=M_1^{new}+M_2,`$ (16) where $`M_1^{new}`$ $`=`$ $`ie{\displaystyle \frac{G_F}{\sqrt{2}}}V_{cb}f_{B_c}m_\tau \epsilon _\alpha C^{2HDM}\overline{u}(p_1)\left\{{\displaystyle \frac{\gamma _\alpha \overline{)}q+2p_{1\alpha }}{2p_1.q}}{\displaystyle \frac{p_\alpha }{p.q}}\right\}\left(1\gamma _5\right)v(p_2).`$ (17) Therefore, in this model the charged Higgs contribution modifies only the so-called $`M_1`$ part of the SM, and it does not induce any new contribution to the so-called $`M_2`$ (see Eq.3): $`C^{2HDM}=\left\{{\displaystyle \frac{m_{B_c}^2}{M_H^2}}{\displaystyle \frac{(m_bZYm_cZX)}{(m_b+m_c)}}+1\right\}.`$ (18) The 2HDM is sensitive to two basic free parameters, namely $`tan\beta `$, and the charged Higgs mass $`M_H`$. If we formally set $`Z0`$ in Eq. (18), the resulting expression is expected to coincide with the $`B_c\tau \overline{\nu }_\tau \gamma `$ decay, which was investigated in the framework of SM . After lengthy, but straightforward calculation for the squared matrix element, we get: $`M_{total}^2=M_1^{new}^2+2\text{Re}\left[M_1^{new}M_2^{}\right]+M_2^2,`$ (19) where $`M_1^{new}^2`$ $`=`$ $`{\displaystyle \frac{G^2}{2}}V_{cb}^2e^2(4f_{B_c}^2m_\tau ^2)C^{2HDM}^2{\displaystyle \frac{1}{(p_1q)^2(pq)^2}}`$ (20) $`\times `$ $`\{2p^2(p_1p_2)(p_1q)^2+(pq)^2[(p_1p_2)(2m_\tau ^2p_1q)+(p_2q)(m_\tau ^22p_1q)]`$ $`+`$ $`(pq)(p_1q)[(pp_2)(p_1q)(pp_1)(4p_1p_2+p_2q)]\},`$ $`2\text{Re}\left[M_1^{new}M_2^{}\right]`$ $`=`$ $`{\displaystyle \frac{G^2}{2}}V_{cb}^2e^2(16f_{B_c}m_\tau ^2)C^{2HDM}{\displaystyle \frac{1}{(p_1q)(pq)}}`$ (21) $`\times `$ $`\{{\displaystyle \frac{f_2(p^2)}{m_{B_c}^2}}p^2(p_1q)(p_2q)+(pq)^2[{\displaystyle \frac{f_2(p^2)}{m_{B_c}^2}}(p_1p_2+p_2q){\displaystyle \frac{f_1(p^2)}{m_{B_c}^2}}(p_2q)]`$ $``$ $`{\displaystyle \frac{f_2(p^2)}{m_{B_c}^2}}[(pp_2)(p_1q)+(pp_1)(p_2q)]\},`$ $`M_2^2`$ $`=`$ $`{\displaystyle \frac{G^2}{2}}V_{cb}^2e^216\left[{\displaystyle \frac{f_1(p^2)^2}{m_{B_c}^4}}+{\displaystyle \frac{f_2(p^2)^2}{m_{B_c}^4}}\right]`$ (22) $`\times `$ $`\left\{(pp_2)(pq)(p_1q)+(p_2q)\left[(pp_1)(pq)p^2(p_1q)\right]\right\}.`$ All calculations have been performed in the rest frame of the $`B_c`$ meson. The dot products of the four–vectors are defined if the photon and neutrino (or electron) energies are specified. The Dalitz boundary for the photon energy $`E_\gamma `$ and neutrino energy $`E_2`$ are defined as follows: $`{\displaystyle \frac{m_{B_c}^22m_{B_c}E_\gamma m_\tau ^2}{2m_{B_c}}}`$ $`E_2`$ $`{\displaystyle \frac{m_{B_c}^22m_{B_c}E_\gamma m_\tau ^2}{2(m_{B_c}2E_\gamma )}},`$ $`0`$ $`E_\gamma `$ $`{\displaystyle \frac{m_{B_c}^2m_\tau ^2}{2m_{B_c}}}.`$ (23) It is now straightforward to work out the expression for the differential decay rate in the lepton and photon energies: $`{\displaystyle \frac{d\mathrm{\Gamma }}{dE_2dE_\gamma }}={\displaystyle \frac{1}{64\pi ^3m_{B_c}}}M_{total}^2.`$ (24) The differential $`(d\mathrm{\Gamma }/dE_\gamma )`$ and total decay width are singular at the lower limit of the photon energy, and this singularity which is present only in the $`M_1^{new}^2`$ contribution is due to the soft photon emission from charged lepton line. On the other hand, $`M_2^2`$ and $`\text{Re}\left[M_1^{new}M_2^{}\right]`$ terms are free from this singularity. In this limit the $`B_c\tau \overline{\nu }_\tau \gamma `$ decay can not be distinguished from the $`B_c\tau \overline{\nu }_\tau `$ decay. Therefore, in order to obtain a finite result for the decay width, we must consider both decays together. The infrared singularity arising in the $`M_1^{new}^2`$ contribution must be canceled with $`O(\alpha )`$ virtual correction to the $`B_c\tau \overline{\nu }_\tau `$ decay. In the SM this cancellation explicitly was shown in . In this work, the $`B_c\tau \overline{\nu }_\tau \gamma `$ process is not considered as a $`O(\alpha )`$ correction to the $`B_c\tau \overline{\nu }_\tau `$ decay, but rather a separate decay channel with hard photon radiation. Therefore we impose a cut off value on the photon energy, which will set an experimental limit on the minimum detectable photon energy. We consider the case for which the photon energy threshold is larger than $`50`$, MeV i.e., $`E_\gamma am_{B_c}`$, where $`a0.01`$. An integration over all the possible values of the lepton energy $`E_2`$ gives the total decay width as a function of the photon energy: $`\mathrm{\Gamma }`$ $`=`$ $`{\displaystyle \frac{G^2\alpha m_{B_c}^3}{64\pi ^2}}V_{cb}^2\{4f_{B_c}^2C^{2HDM}^2{\displaystyle _\delta ^{1r}}dx{\displaystyle \frac{r}{x(1x)}}[4+8r4r^2+10x14rx`$ (25) $`+`$ $`4r^2x9x^2+7rx^2+3x^3+(1x)(22r^23x+rx+2x^2)\mathrm{}n({\displaystyle \frac{1x}{r}})]`$ $``$ $`4f_{B_c}C^{2HDM}{\displaystyle _\delta ^{1r}}dx{\displaystyle \frac{rx}{1x}}[(1rx)(f_1(x)x+f_2(x)(1+r2x))`$ $``$ $`(1x)(f_1(x)x+f_2(x)(2rx))\mathrm{}n({\displaystyle \frac{1x}{r}})]`$ $`+`$ $`{\displaystyle \frac{1}{3}}{\displaystyle _\delta ^{1r}}dx[f_1(x)^2+f_2(x)^2]{\displaystyle \frac{1}{(1x)^2}}x^3(2+r2x)(1rx)^2\},`$ where $`x=2E_\gamma /m_{B_c}`$is the dimensionless photon energy, $`r=m_\tau ^2/m_{B_c}^2`$ and $`\delta =2a`$. ## 3 NUMERICAL ANALYSIS To calculate the decay width, explicit forms of the form factors $`f_1`$ and $`f_2`$ are needed. These form factors were calculated in the framework of the light-front quark model in , and in the light-cone QCD sum rules , where in it was found that, the best agreement is achieved by the following pole forms for the form factors: $`f_1(p^2)={\displaystyle \frac{f_1(0)}{1p^2/m_1^2}},f_2(p^2)={\displaystyle \frac{f_2(0)}{1p^2/m_2^2}},`$ (26) where $`f_1(0)=0.44\pm 0.04\text{GeV},m_1^2=43.1\text{GeV}^2,`$ $`f_2(0)=0.21\pm 0.02\text{GeV},m_2^2=48.0\text{GeV}^2.`$ On the other hand, in evaluating the decay width, we have used the following set of parameters: $`G_F=1.17.10^5GeV^2`$, $`\alpha =1/137`$, $`m_b=4.8GeV`$, $`m_c=1.4GeV`$, $`m_{B_c}=6.3GeV`$, $`m_\tau =1.78GeV`$, $`f_{B_c}=0.36`$ GeV , $`V_{cb}=0.04`$ and, $`\tau (B_c)=0.46\times 10^{12}s`$ . In this regard we should also recall that the free parameters of the 2HDM model namely $`tan\beta `$, and $`M_H`$ are not arbitrary, but there are some semiquantitative restrictions on them using the existing experimental data. The most direct bound on the charged Higgs boson mass comes from top quark decays, which yield the bound $`M_H>\mathrm{\hspace{0.17em}147}GeV`$ for large $`tan\beta `$ . For pure type-II 2HDM’s one finds $`M_H>\mathrm{\hspace{0.17em}300}GeV`$, coming from the virtual Higgs boson contributions to $`bs\gamma `$ . Furthermore, there are no experimental upper bounds on the mass of the charged Higgs boson, but one generally expects to have $`M_H<1TeV`$ in order that perturbation theory remain valid . For large $`tan\beta `$ the most stringent constraints on $`tan\beta `$ and $`M_H`$ are actually on their ratio, $`tan\beta /M_H`$. The current limits come from the measured branching ratio for the inclusive decay $`BX\tau \overline{\nu }`$, giving $`tan\beta /M_H<0.46GeV^1`$ , and from the upper limit on the branching ratio for $`B\tau \overline{\nu }`$, giving $`tan\beta /M_H<0.38GeV^1`$ . For illustrative purposes we consider four values of $`tan\beta `$, namely $`tan\beta =`$5, 10, 30 and 50 and let $`m_{H^\pm }=150GeV`$. Then we consider two values of $`m_{H^\pm }`$, namely $`M_{H^\pm }`$=200, 400 GeV, and we allow $`tan\beta `$ to range between 0 to 60. The results of this numerical analysis are graphically shown in figures 2-4. In these figures the differences between the 2HDM’s and the SM are shown for two different fixed cut off values, i.e., $`\delta =0.016`$ and $`\delta =0.032`$ both for Model I and Model II. The results for the differential decay branching ratio dBR($`B_c\tau \overline{\nu }\gamma )/dx`$ as a function of $`x=2E_\gamma /m_{B_c}`$ for different values of $`tan\beta `$, $`M_H=150GeV`$ are presented in Fig.2, while the branching ratio (BR) for $`B\tau \overline{\nu }\gamma `$ decay is shown in figure 3-4 as a function of $`M_H`$ for various values of $`tan\beta `$, and as a function of $`tan\beta `$ for different values of $`M_H`$. Results are shown for Model I, and Mode lI. It is observed that Model II gives both a bigger differential decay branching ratio, and a bigger branching ratio than the SM rates of (up to three orders of magnitude \[$`M_H=150GeV`$\]) for large values of $`tan\beta >\mathrm{\hspace{0.17em}20}`$, while for small values of $`tan\beta <\mathrm{\hspace{0.17em}10}`$ results approaches its SM value. In model I the situation is somewhat totally different. Curves overlap with the SM results all the way. This behavior obviously reflects the $`H^\pm `$ fermion couplings, which are proportional to $`cot\beta `$ in this model. In conclusion, this study shows that the branching ratios for $`B_c\tau \overline{\nu }\gamma `$ could be at the level of $`10^4`$ in the 2HDM, which may be detectable at the ongoing LHC. When enough $`B_c`$ events are collected, this decay will be able to provide alternative channel to extract new restrictions for the free parameters $`tan\beta `$ and $`M_{H^+}`$ of the 2HDM model. Figure Captions Figure 1. : The relevant Feynman diagrams, responsible for $`B_c\tau \overline{\nu }\gamma `$ decay. Figure 2. : The dependence of the differential Branching ratio of dBR($`B_c\tau \overline{\nu }\gamma )/dx`$ as a function of the photon energy $`x=2E_\gamma /m_{B_c}`$ for both models Model I, and Model II. Figure 3. : The dependence of the Branching ratio on the charged Higgs boson mass at different values of $`tan\beta `$ for both models Model I, and Model II. Figure 4. : The dependence of the Branching ratio on $`tan\beta `$ at different values of the charged Higgs boson mass for both models Model I, and Model II.
warning/0004/cond-mat0004181.html	ar5iv	text	# 1 Introduction ## 1 Introduction The $`q`$-state Potts model has served as a valuable model for the study of phase transitions and critical phenomena . On a lattice, or, more generally, on a (connected) graph $`G`$, at temperature $`T`$, this model is defined by the partition function $$Z(G,q,v)=\underset{\{\sigma _n\}}{}e^\beta $$ (1.1) with the (zero-field) Hamiltonian $$=J\underset{ij}{}\delta _{\sigma _i\sigma _j}$$ (1.2) where $`\sigma _i=1,\mathrm{},q`$ are the spin variables on each vertex $`iG`$; $`\beta =(k_BT)^1`$; and $`ij`$ denotes pairs of adjacent vertices. The graph $`G=G(V,E)`$ is defined by its vertex set $`V`$ and its edge set $`E`$; we denote the number of vertices of $`G`$ as $`n=n(G)=\|V\|`$ and the number of edges of $`G`$ as $`e(G)=\|E\|`$. We use the notation $$K=\beta J,a=u^1=e^K,v=a1$$ (1.3) so that the physical ranges are (i) $`a1`$, i.e., $`v0`$ corresponding to $`\mathrm{}T0`$ for the Potts ferromagnet, and (ii) $`0a1`$, i.e., $`1v0`$, corresponding to $`0T\mathrm{}`$ for the Potts antiferromagnet. One defines the (reduced) free energy per site $`f=\beta F`$, where $`F`$ is the actual free energy, via $$f(\{G\},q,v)=\underset{n\mathrm{}}{lim}\mathrm{ln}[Z(G,q,v)^{1/n}].$$ (1.4) where we use the symbol $`\{G\}`$ to denote $`lim_n\mathrm{}G`$ for a given family of graphs. Let $`G^{}=(V,E^{})`$ be a spanning subgraph of $`G`$, i.e. a subgraph having the same vertex set $`V`$ and an edge set $`E^{}E`$. Then $`Z(G,q,v)`$ can be written as the sum - $`Z(G,q,v)`$ $`=`$ $`{\displaystyle \underset{G^{}G}{}}q^{k(G^{})}v^{e(G^{})}`$ (1.5) $`=`$ $`{\displaystyle \underset{r=k(G)}{\overset{n(G)}{}}}{\displaystyle \underset{s=0}{\overset{e(G)}{}}}z_{rs}q^rv^s`$ (1.8) where $`k(G^{})`$ denotes the number of connected components of $`G^{}`$ and $`z_{rs}0`$. Since we only consider connected graphs $`G`$, we have $`k(G)=1`$. The formula (1.5) enables one to generalize $`q`$ from $`_+`$ to $`_+`$ (keeping $`v`$ in its physical range). This generalization is sometimes denoted the random cluster model ; here we shall use the term “Potts model” to include both positive integral $`q`$ as in the original formulation in eqs. (1.1) and (1.2), and the generalization to real (or complex) $`q`$, via eq. (1.5). The formula (1.5) shows that $`Z(G,q,v)`$ is a polynomial in $`q`$ and $`v`$ (equivalently, $`a`$) with maximum and minimum degrees indicated in eq. (1.8). The Potts model partition function on a graph $`G`$ is essentially equivalent to the Tutte polynomial - and Whitney rank polynomial , , - for this graph, as discussed in the appendix. In this paper we shall present exact calculations of the Potts model partition function for strips of the triangular lattice of width $`L_y=2`$ vertices (or equivalently edges) and arbitrary length $`L_x`$ with various boundary conditions. This is a natural continuation of a previous study by one of us of the Potts model model on the analogous strips of the square lattice , and the reader is referred to that paper for background and further references. We envision the strip of the triangular lattice as being formed by starting with a ladder graph, i.e. $`2\times L_x`$ strip of the square lattice, and then adding edges joining, say, the lower left and upper right vertices of each square. The longitudinal (transverse) direction is taken as the horizontal, $`x`$ (vertical, $`y`$) direction. We use free transverse boundary conditions and consider free, periodic (= cyclic), and Möbius longitudinal boundary conditions. These families of graphs are denoted, respectively, as $`S_m`$ (for open strip), $`L_m`$ (for ladder), and $`ML_m`$ (for Möbius ladder), where $`L_x=m+1`$ (edges) for $`S_m`$ (following our labelling convention in ) and $`L_x=m`$ for $`L_m`$ and $`ML_m`$. One has $`n(S_m)=2(m+2)`$ and $`n(L_m)=n(ML_m)=2m`$. Each vertex on the cyclic strip $`L_m`$ has degree (coordination number) $`\mathrm{\Delta }=4`$; this is also true of the interior vertices on the open strip $`S_m`$. The Möbius strip can involve a seam, as discussed in . The motivations for these exact calculations of Potts model partition functions for strips of various lattices were discussed in . Clearly, new exact calculations of Potts model partition functions are of value in their own right. In addition, it was shown that these calculations can give insight into the complex-temperature phase diagram of the 2D Potts model on the given lattice. This is useful, since the 2D Potts model has never been solved except in the $`q=2`$ Ising case. Furthermore, with these exact results one can study both the $`T=0`$ critical point of the $`q`$-state Potts ferromagnet and, for certain $`q`$ values ($`q=2`$ for the square strip; $`q=2,3`$ for the present strip of the triangular lattice) the $`T=0`$ critical point of the Potts antiferromagnet. In addition, via the correspondence with the Tutte polynomial, our calculations yield several quantities of relevance to mathematical graph theory. Various special cases of the Potts model partition function are of interest. One special case is the zero-temperature limit of the Potts antiferromagnet (AF). For sufficiently large $`q`$, on a given lattice or graph $`G`$, this exhibits nonzero ground state entropy (without frustration). This is equivalent to a ground state degeneracy per site (vertex), $`W>1`$, since $`S_0=k_B\mathrm{ln}W`$. The $`T=0`$ (i.e., $`v=1`$) partition function of the above-mentioned $`q`$-state Potts antiferromagnet (PAF) on $`G`$ satisfies $$Z(G,q,1)=P(G,q)$$ (1.9) where $`P(G,q)`$ is the chromatic polynomial (in $`q`$) expressing the number of ways of coloring the vertices of the graph $`G`$ with $`q`$ colors such that no two adjacent vertices have the same color . The minimum number of colors necessary for this coloring is the chromatic number of $`G`$, denoted $`\chi (G)`$. We have $$W(\{G\},q)=\underset{n\mathrm{}}{lim}P(G,q)^{1/n}$$ (1.10) Using the formula (1.5) for $`Z(G,q,v)`$, one can generalize $`q`$ from $`_+`$ not just to $`_+`$ but to $``$ and $`a`$ from its physical ferromagnetic and antiferromagnetic ranges $`1a\mathrm{}`$ and $`0a1`$ to $`a`$. A subset of the zeros of $`Z`$ in the two-complex dimensional space $`^2`$ defined by the pair of variables $`(q,a)`$ can form an accumulation set in the $`n\mathrm{}`$ limit, denoted $``$, which is the continuous locus of points where the free energy is nonanalytic. This locus is determined as the solution to a certain $`\{G\}`$-dependent equation . For a given value of $`a`$, one can consider this locus in the $`q`$ plane, and we denote it as $`_q(\{G\},a)`$. In the special case $`a=0`$ (i.e., $`v=1`$) where the partition function is equal to the chromatic polynomial, the zeros in $`q`$ are the chromatic zeros, and $`_q(\{G\},a=0)`$ is their continuous accumulation set in the $`n\mathrm{}`$ limit -. In a series of papers we have given exact calculations of the chromatic polynomials and nonanalytic loci $`_q`$ for various families of graphs (for further references on this $`a=0`$ special case, see ). With the exact Potts partition function for arbitrary temperature, one can study $`_q`$ for $`a0`$ and, for a given value of $`q`$, one can study the continuous accumulation set of the zeros of $`Z(G,q,v)`$ in the $`a`$ plane; this will be denoted $`_a(\{G\},q)`$. It will often be convenient to consider the equivalent locus in the $`u=1/a`$ plane, namely $`_u(\{G\},q)`$. We shall sometimes write $`_q(\{G\},a)`$ simply as $`_q`$ when $`\{G\}`$ and $`a`$ are clear from the context, and similarly with $`_a`$ and $`_u`$. One gains a unified understanding of the separate loci $`_q(\{G\})`$ for fixed $`a`$ and $`_a(\{G\})`$ for fixed $`q`$ by relating these as different slices of the locus $``$ in the $`^2`$ space defined by $`(q,a)`$ as we shall do here. Following the notation in and our other earlier works on $`_q(\{G\})`$ for $`a=0`$, we denote the maximal region in the complex $`q`$ plane to which one can analytically continue the function $`W(\{G\},q)`$ from physical values where there is nonzero ground state entropy as $`R_1`$ . The maximal value of $`q`$ where $`_q`$ intersects the (positive) real axis was labelled $`q_c(\{G\})`$. Thus, region $`R_1`$ includes the positive real axis for $`q>q_c(\{G\})`$. Correspondingly, in our works on complex-temperature properties of spin models, we had labelled the complex-temperature extension (CTE) of the physical paramagnetic phase as (CTE)PM, which will simply be denoted PM here, the extension being understood, and similarly with ferromagnetic (FM) and antiferromagnetic (AFM); other complex-temperature phases, having no overlap with any physical phase, were denoted $`O_j`$ (for “other”), with $`j`$ indexing the particular phase . Here we shall continue to use this notation for the respective slices of $``$ in the $`q`$ and $`a`$ or $`u`$ planes. We record some special values of $`Z(G,q,v)`$ below, beginning with the $`q=0`$ special case $$Z(G,0,v)=0$$ (1.11) which implies that $`Z(G,q,v)`$ has an overall factor of $`q`$. In general (and for all the graphs considered here), this is the only overall factor that it has. We also have $$Z(G,1,v)=\underset{G^{}G}{}v^{e(G^{})}=a^{e(G)}.$$ (1.12) For temperature $`T=\mathrm{}`$, i.e., $`v=0`$, $$Z(G,q,0)=q^{n(G)}.$$ (1.13) Further, $$Z(G,q,1)=P(G,q)=\left[\underset{s=0}{\overset{\chi (G)1}{}}(qs)\right]U(G,q)$$ (1.14) where $`U(G,q)`$ is a polynomial in $`q`$ of degree $`n(G)\chi (G)`$. We note some values of chromatic numbers for the strip graphs considered here: $$\chi (S_m)=3$$ (1.15) $$\chi (L_m)=\{\begin{array}{cc}3\hfill & \text{if }m=0\text{ mod 3}\hfill \\ 4\hfill & \text{if }m=1\text{ or }m=2\text{ mod 3}\hfill \end{array}$$ (1.16) and $$\chi (ML_m)=4.$$ (1.17) Hence, for $`q=1,2,3`$ we have $$Z(G,q,1)=P(G,q)=0\mathrm{for}G=S_m,L_m,ML_m\mathrm{and}q=1,2$$ (1.18) $$Z(G,3,1)=P(G,3)=3!\mathrm{for}G=S_m\mathrm{or}G=L_{m=0mod3}$$ (1.19) where $`G=L_{m=0mod3}`$ means that this applies to $`L_m`$ for $`m=0`$ mod 3. For the graphs $`G=S_m,L_m`$, and $`ML_m`$, (i) the result (1.18) implies that for $`q=2`$, $`Z(G,2,v)`$ contains at least one power of the factor $`(v+1)=a`$; for $`q=1`$, one already knows the form of $`Z(G,1,v)`$ from (1.12); (ii) the result (1.19) implies that for the cases where $`\chi (G)=4`$, $`Z(G,3,v)`$ contains at least one power of $`(v+1)`$ as a factor. Another basic property, evident from eq. (1.5), is that (i) the zeros of $`Z(G,q,v)`$ in $`q`$ for real $`v`$ and hence also the continuous accumulation set $`_q`$ are invariant under the complex conjugation $`qq^{}`$; (ii) the zeros of $`Z(G,q,v)`$ in $`v`$ or equivalently $`a`$ for real $`q`$ and hence also the continuous accumulation set $`_a`$ are invariant under the complex conjugation $`aa^{}`$. Just as the importance of noncommutative limits was shown in (eq. (1.9) of) on chromatic polynomials, so also one encounters an analogous noncommutativity here for the general partition function (1.5) of the Potts model for nonintegral $`q`$: at certain special points $`q_s`$ (typically $`q_s=0,1\mathrm{},\chi (G)`$) one has $$\underset{n\mathrm{}}{lim}\underset{qq_s}{lim}Z(G,q,v)^{1/n}\underset{qq_s}{lim}\underset{n\mathrm{}}{lim}Z(G,q,v)^{1/n}.$$ (1.20) Because of this noncommutativity, the formal definition (1.4) is, in general, insufficient to define the free energy $`f`$ at these special points $`q_s`$; it is necessary to specify the order of the limits that one uses in eq. (1.20). We denote the two definitions using different orders of limits as $`f_{qn}`$ and $`f_{nq}`$: $$f_{nq}(\{G\},q,v)=\underset{n\mathrm{}}{lim}\underset{qq_s}{lim}n^1\mathrm{ln}Z(G,q,v)$$ (1.21) $$f_{qn}(\{G\},q,v)=\underset{qq_s}{lim}\underset{n\mathrm{}}{lim}n^1\mathrm{ln}Z(G,q,v).$$ (1.22) In Ref. and our subsequent works on chromatic polynomials and the above-mentioned zero-temperature antiferromagnetic limit, it was convenient to use the ordering $`W(\{G\},q_s)=lim_{qq_s}lim_n\mathrm{}P(G,q)^{1/n}`$ since this avoids certain discontinuities in $`W`$ that would be present with the opposite order of limits. In the present work on the full temperature-dependent Potts model partition function, we shall consider both orders of limits and comment on the differences where appropriate. Of course in discussions of the usual $`q`$-state Potts model (with positive integer $`q`$), one automatically uses the definition in eq. (1.1) with (1.2) and no issue of orders of limits arises, as it does in the Potts model with real $`q`$. As a consequence of the noncommutativity (1.20), it follows that for the special set of points $`q=q_s`$ one must distinguish between (i) $`(_a(\{G\},q_s))_{nq}`$, the continuous accumulation set of the zeros of $`Z(G,q,v)`$ obtained by first setting $`q=q_s`$ and then taking $`n\mathrm{}`$, and (ii) $`(_a(\{G\},q_s))_{qn}`$, the continuous accumulation set of the zeros of $`Z(G,q,v)`$ obtained by first taking $`n\mathrm{}`$, and then taking $`qq_s`$. For these special points, $$(_a(\{G\},q_s))_{nq}(_a(\{G\},q_s))_{qn}.$$ (1.23) From eq. (1.11), it follows that for any $`G`$, $$\mathrm{exp}(f_{nq})=0\mathrm{for}q=0$$ (1.24) and thus $$(_a)_{nq}=\mathrm{}\mathrm{for}q=0.$$ (1.25) However, for many families of graphs, including the circuit graph $`C_n`$, and cyclic and Möbius strips of the square or triangular lattice, if we take $`n\mathrm{}`$ first and then $`q0`$, we find that $`(_u)_{qn}`$ is nontrivial. Similarly, from (1.12) we have, for any $`G`$, $$(_a)_{nq}=\mathrm{}\mathrm{for}q=1$$ (1.26) since all of the zeros of $`Z`$ occur at the single discrete point $`a=0`$ (and in the case of a graph $`G`$ with no edges, $`Z=1`$ with no zeros). However, as the simple case of the circuit graph shows , $`(_u)_{qn}`$ is, in general, nontrivial. As derived in , a general form for the Potts model partition function for the strip graphs considered here, or more generally, for recursively defined families of graphs comprised of $`m`$ repeated subunits (e.g. the columns of squares of height $`L_y`$ vertices that are repeated $`L_x`$ times to form an $`L_x\times L_y`$ strip of a regular lattice with some specified boundary conditions), is $$Z(G,q,v)=\underset{j=1}{\overset{N_\lambda }{}}c_{G,j}(\lambda _{G,j}(q,v))^m$$ (1.27) where $`N_\lambda `$ depends on $`G`$. The Potts ferromagnet has a zero-temperature phase transition in the $`L_x\mathrm{}`$ limit of the strip graphs considered here, and this has the consequence that for cyclic and Möbius boundary conditions, $``$ passes through the $`T=0`$ point $`u=0`$. It follows that $``$ is noncompact in the $`a`$ plane. Hence, it is usually more convenient to study the slice of $``$ in the $`u=1/a`$ plane rather than the $`a`$ plane. Since $`a\mathrm{}`$ as $`T0`$ and $`Z`$ diverges like $`a^{e(G)}`$ in this limit, we shall use the reduced partition function $`Z_r`$ defined by $$Z_r(G,q,v)=a^{e(G)}Z(G,q,v)=u^{e(G)}Z(G,q,v)$$ (1.28) which has the finite limit $`Z_r1`$ as $`T0`$. For a general strip graph $`(G_s)_m`$ of type $`G_s`$ and length $`L_x=m`$, we can write $`Z_r((G_s)_m,q,v)`$ $`=`$ $`u^{e((G_s)_m)}{\displaystyle \underset{j=1}{\overset{N_\lambda }{}}}c_{G_s,j}(\lambda _{G_s,j})^m{\displaystyle \underset{j=1}{\overset{N_\lambda }{}}}c_{G_s,j}(\lambda _{G_s,j,u})^m`$ (1.29) with $$\lambda _{G_s,j,u}=u^{e((G_s)_m)/m}\lambda _{G_s,j}.$$ (1.30) For example, for the $`L_y=2`$ strips of the triangular lattice of interest here, with free transverse boundary conditions and either free or periodic longitudinal boundary conditions, we have $`e(S_m)=4m+5`$ and $`e(L_m)=4m`$. (In the case of the cyclic strip with $`m=2`$, $`L_m`$ degenerates in the sense that it contains multiple edges connecting some pairs of vertices, and similarly, for $`m=1`$, $`L_m`$ contains both multiple edges and loops; the above formula applies to these cases if one takes care to count these multiple edges and loops, as is discussed further below.) ## 2 Strip of Triangular Lattice with Free Longitudinal Boundary Conditions In this section we present the Potts model partition function $`Z(S_m,q,v)`$ for the $`L_y=2`$ strip of the triangular lattice $`S_m`$ with arbitrary length $`L_x=m+1`$ (i.e., containing $`m+1`$ squares) and free transverse and longitudinal boundary conditions. One convenient way to express the results is in terms of a generating function: $$\mathrm{\Gamma }(S,q,v,z)=\underset{m=0}{\overset{\mathrm{}}{}}Z(S_m,q,v)z^m.$$ (2.1) We have calculated this generating function using the deletion-contraction theorem for the corresponding Tutte polynomial $`T(S_m,x,y)`$ and then expressing the result in terms of the variables $`q`$ and $`v`$. We find $$\mathrm{\Gamma }(S,q,v,z)=\frac{𝒩(S,q,v,z)}{𝒟(S,q,v,z)}$$ (2.2) where $$𝒩(S,q,v,z)=A_{S,0}+A_{S,1}z$$ (2.3) with $$A_{S,0}=q(v^5+5v^4+8v^3+2v^3q+10v^2q+5vq^2+q^3)$$ (2.4) $$A_{S,1}=q(v+1)^2(v+q)^3v^2$$ (2.5) and $$𝒟(S,q,v,z)=1(v^4+4v^3+7v^2+4qv+q^2)z+(v+1)^2(v+q)^2v^2z^2.$$ (2.6) (The generating function for the Tutte polynomial $`T(S_m,x,y)`$ is given in the appendix.) Writing $$𝒟(S,q,v,z)=\underset{j=1}{\overset{2}{}}(1\lambda _{S,j}z)$$ (2.7) we have $$\lambda _{S,(1,2)}=\frac{1}{2}\left[T_{S12}\pm (3v+v^2+q)\sqrt{R_{S12}}\right]$$ (2.8) where $$T_{S12}=v^4+4v^3+7v^2+4qv+q^2$$ (2.9) and $$R_{S12}=q^2+2qv2qv^2+5v^2+2v^3+v^4.$$ (2.10) Ref. presented a formula to obtain the chromatic polynomial for a recursive family of graphs in the form of sums of powers of $`\lambda _j`$’s starting from the generating function, and the generalization of this to the full Potts model partition function was given in . Using this, we have $$Z(S_m,q,v)=\frac{(A_{S,0}\lambda _{S,1}+A_{S,1})}{(\lambda _{S,1}\lambda _{S,2})}\lambda _{S,1}^m+\frac{(A_{S,0}\lambda _{S,2}+A_{S,1})}{(\lambda _{S,2}\lambda _{S,1})}\lambda _{S,2}^m$$ (2.11) (which is symmetric under $`\lambda _{S,1}\lambda _{S,2}`$). Although both the $`\lambda _{S,j}`$’s and the coefficient functions involve the square root $`\sqrt{R_{S12}}`$ and are not polynomials in $`q`$ and $`v`$, the theorem on symmetric functions of the roots of an algebraic equation guarantees that $`Z(S_m,q,v)`$ is a polynomial in $`q`$ and $`v`$ (as it must be by (1.5) since the coefficients of the powers of $`z`$ in the equation (2.6) defining these $`\lambda _{S,j}`$’s are polynomials in these variables $`q`$ and $`v`$. As will be shown below, for $`q2`$ (and for $`0q<1`$, with the $`f_{qn}`$ and $`(_u)_{qn}`$ definitions) the singular locus $`_u`$ for this strip consists of arcs that do not separate the $`u`$ plane into different regions, so that the PM phase and its complex-temperature extension occupy all of this plane, except for these arcs. For these ranges of $`q`$, the (reduced) free energy is given by $$f=\frac{1}{2}\mathrm{ln}\lambda _{S,1}.$$ (2.12) This is analytic for all finite temperature, for both the ferromagnetic and antiferromagnetic sign of the spin-spin coupling $`J`$. The internal energy and specific heat can be calculated in a straightforward manner from the free energy (2.12); since the resultant expressions are somewhat cumbersome, we do not list them here. In contrast, for the range $`1<q<2`$, the Potts antiferromagnet on the $`L_x\mathrm{}`$ limit of this open strip does have a phase transition at finite temperature (see eq. (2.2.6)), but this has unphysical features, including a specific heat and partition function $`Z`$ that are negative for some range of temperature, and non-existence of a thermodynamic limit independent of boundary conditions. This will be discussed further below. We next consider the limits of zero temperature for the antiferromagnetic and ferromagnetic cases. Let us define $$D_k(q)=\frac{P(C_k,q)}{q(q1)}=\underset{s=0}{\overset{k2}{}}(1)^s\left(\genfrac{}{}{0pt}{}{k1}{s}\right)q^{k2s}$$ (2.13) and $`P(C_k,q)`$ is the chromatic polynomial for the circuit (cyclic) graph $`C_k`$ with $`k`$ vertices, $$P(C_k,q)=(q1)^k+(q1)(1)^k$$ (2.14) so that $`D_2=1`$, $`D_3=q2`$, $`D_4=q^23q+3`$, and so forth for other $`D_k`$’s. In the $`T=0`$ Potts antiferromagnet limit $`v=1`$, $`\lambda _{S,1}=D_3^2=(q2)^2`$ and $`\lambda _{S,2}=0`$, so that eq. (2.2) reduces to the generating function for the chromatic polynomial for this open strip $$\mathrm{\Gamma }(S,q,1;z)=\frac{q(q1)(q2)^2}{1(q2)^2z}$$ (2.15) Equivalently, the chromatic polynomial is $$P(S_m,q)=q(q1)(q2)^{2(m+1)}.$$ (2.16) For the ferromagnetic case with general $`q`$, in the low-temperature limit $`v\mathrm{}`$, $$\lambda _{S,1}=v^4+4v^3+O(v^2),\lambda _{S,2}=v^2+2(q1)v+O(1)\mathrm{as}v\mathrm{}$$ (2.17) so that $`\|\lambda _{S,1}\|`$ is never equal to $`\|\lambda _{S,2}\|`$ in this limit, and hence $`_u`$ does not pass through the origin of the $`u`$ plane for the $`n\mathrm{}`$ limit of the open square strip: $$u=0_u(\{S\}).$$ (2.18) In contrast, as will be shown below, $`_u`$ does pass through $`u=0`$ for this strip with cyclic or Möbius boundary conditions. ### 2.1 $`_q(\{S\})`$ for fixed $`a`$ We start with the value $`a=0`$ corresponding to the Potts antiferromagnet at zero temperature. In this case, $`Z(S_m,q,v=1)=P(S_m,q)`$, where this chromatic polynomial was given in eq. (2.16), and the continuous locus $`_q=\mathrm{}`$. For $`a`$ in the finite-temperature antiferromagnetic range $`0<a<1`$, $`_q`$ consists of a single self-conjugate arc crossing the positive real $`q`$ axis at $$q_c(\{S\})=v(3+v)=(1a)(2+a)$$ (2.1.1) where the factor $`(3v+v^2+q)`$ multiplying the square root in eq. (2.8) vanishes, so that the two roots coincide. Parenthetically, we observe that this is same as the value of $`q_c`$ for the Potts model on the $`n\mathrm{}`$ limit of the circuit graph, $`C_n`$ . The arc endpoints occur at the branch points where the square root is zero, viz., $$q_{S,endpt.}=(a1)[a2\pm 2i\sqrt{a}].$$ (2.1.2) As $`a`$ increases from 0 to 1, the crossing point (2.1.1) decreases monotonically from 2 to 0, and as $`a`$ reaches the infinite-temperature value 1, $`_q`$ shrinks to a point at the origin. In the ferromagnetic range $`a>1`$ the self-conjugate arc crosses the negative real axis, since $`q_c(\{S\})<0`$, and has support in the $`Re(q)<0`$ half-plane. In Figs. 1 and 2 we show $`_q`$ and associated zeros of $`Z`$ in the $`q`$ plane for the antiferromagnetic value $`a=0.1`$ and the ferromagnetic value $`a=2`$. ### 2.2 $`_u`$ for fixed $`q`$ We show several plots of the locus $`_u`$ for various values of $`q`$ in Figs. 3 \- 8. For values of $`q`$ where noncommutativity occurs, we consider $`_{qn}`$. Given the algebraic structure of $`\lambda _{S,j}`$, $`j=1,2`$, the degeneracy of magnitudes $`\|\lambda _{S,1}\|=\|\lambda _{S,2}\|`$ and hence the locus $`_u`$ occurs where (i) $`T_{S12}=0`$, (ii) the prefactor of the square root vanishes: $`3v+v^2+q=0`$, (iii) $`R_{S12}=0`$, (iv) if $`q`$ is real, where $`R_{S12}<0`$ so that the square root is pure imaginary, and (v) elsewhere for complex $`u`$, where the degeneracy condition is satisfied. All of these five possibilities are realized here, although some, such as (i) can yield solutions already subsumed by other conditions; for example, for the case $`q=1.8`$, shown in Fig. 5, $`T_{S12}`$ vanishes at two real points, $`u29.1,1.35`$, both of which are subsumed within the solution of condition (iv), which is a line segment $`69.486<u<1.291`$, and so forth for various other solutions. In cases where $`_u`$ does not enclose regions, $`\lambda _{S,1}`$ is dominant everywhere in the $`u`$ plane, and degenerate in magnitude with $`\lambda _{S,2}`$ on $`_u`$; the cases where $`_u`$ does enclose regions will be discussed individually. For large values of $`q`$, we find that $`_u`$ consists of the union of (i) a circular arc centered at $`u1/q`$ that crosses the negative real $`u`$ axis, and (ii) a line segment on the negative $`u`$ axis. For example, for $`q=10`$, the arc has its endpoints at two of the branch point zeros of $`\sqrt{R_{S12}}`$, at $`u=(7\pm \sqrt{15}i)/320.21875\pm 0.12103i`$, and the line segment, which is the solution to the condition (iv) above, that $`R_{S12}<0`$ for real $`q`$, has left and right endpoints at $`u_1=1`$ and $`u_2=1/4`$. As $`q`$ decreases, $`u_1`$ and $`u_2`$ move toward more negative values, and as $`q`$ decreases to 2, $`u_1\mathrm{}`$, i.e. $`a=0`$, while $`u_21`$. For $`q=2`$, $`_u`$ is the union of the above line segment $`\mathrm{}<u<1`$ and the circular arc $$u=1+\sqrt{2}e^{\pm i\theta },\theta _e\theta \pi $$ (2.2.1) where $$\theta _e=arctan\left(\frac{\sqrt{7}}{11}\right)13.5^{}.$$ (2.2.2) That is, the endpoints of the circular arc occur at $$u_e=\frac{1}{8}(3\pm \sqrt{7}i)$$ (2.2.3) Proceeding, we observe that for $`1<q<2`$, the locus $`_u`$ separates the $`u`$ plane into different regions, while for $`0<q<1`$, as was true for the range $`q>2`$, $`_u`$ does not separate the $`u`$ plane into different regions. In the range $`q_m<q<2`$, where $$q=q_m=\left(\frac{5}{4}\right)^2=1.5625,$$ (2.2.4) $`_u`$ crosses the real $`u`$ axis vertically at three different points and separates the $`q`$ plane into three regions: (i) the paramagnetic phase, which includes the infinite-temperature point $`u=1`$, where $`\lambda _{S,1}`$ is dominant; (ii) the surrounding O phase that extends to the circle at infinity in the $`u`$ plane, i.e. to the point $`a=0`$, in which phase $`\lambda _{S,2}`$ is dominant, and (iii) a third phase, of O type, located slightly to the left of $`u=0`$, in which $`\lambda _{S,2}`$ is dominant. It is interesting that although $`\lambda _{S,2}`$ is dominant in both of the O phases, there is still a boundary that separates them completely; this is a result of the fact that on this boundary, the other $`\lambda `$, namely $`\lambda _{S,1}`$, becomes degenerate in magnitude with $`\lambda _{S,2}`$. Two of the points where $`_u`$ crosses the real axis occur where the condition (ii) holds, i.e., where the prefactor $`3v+v^2+q`$ multiplying the square root in eq. (2.8), vanishes, at $`v=(1/2)[3\pm \sqrt{94q}]`$, i.e., $$u=\frac{\left[1\pm \sqrt{94q}\right]}{2(2q)}.$$ (2.2.5) For $`q=2ϵ`$ with $`0<ϵ<<1`$, these crossings are at $`u=1+O(ϵ)`$ and $`u1/ϵ`$. As $`q`$ decreases, the smaller (larger) solution moves to the right (left). Specifically, as $`q`$ decreases from 2 to 0, the larger solution decreases monotonically from infinity to 1 while the smaller one increases monotonically from $`1`$ to $`1/2`$. For example, for $`q=1.8`$, the crossings in eq. (2.2.5) are at $`u=(1/2)(5\pm 3\sqrt{5})0.854,5.854`$. The existence of the right crossing on the positive real $`u`$ axis for $`u>1`$ means that the free energy of the Potts antiferromagnet is nonanalytic at the temperature $$T_{p,S}=\frac{J}{k_B\mathrm{ln}\left[(1/2)(1+\sqrt{94q})\right]}\mathrm{for}0<q<2$$ (2.2.6) (where both $`J`$ and the log are negative). However, just as found in , this nonanalyticity has associated unphysical features, including negative $`Z`$, negative specific heat in the low-temperature phase, and non-existence of a thermodynamic limit independent of boundary conditions in the low-temperature phase. As $`q`$ decreases from 2 to the value $`q_m`$ given above, the left and right endpoints of the line segment merge, and it shrinks to a point at $`u_1=u_2=4`$. For this value $`q=q_m`$, the crossings in eq. (2.2.5) are $`u=(4/7)(2\pm \sqrt{11})3.038,0.7524`$. In Fig. 6 we show $`_u`$ for $`q=q_m`$. The arc endpoints occur at $`u=(20\pm 8\sqrt{6}i)/490.4082\pm 0.3999i`$. For $`q_m<q<2`$, there is a multiple point on the negative real axis where a branch of $`_u`$ crosses this axis vertically and intersects with the line segment. There are also multiple points at complex-conjugate values of $`u`$ where the circular arc intersects the closed curve. As $`q`$ decreases from 2 to $`q_m`$, the circular arc in the left-hand complex plane enlarges while the closed curve extending into the right-hand plane shrinks and becomes convex. As $`q`$ decreases below $`q_m`$, the closed curve evident to the left in Fig. 6 breaks apart, forming two complex-conjugate arcs, and in the interval $`1<q<q_m`$, $`_u`$ no longer contains any line segment but instead consists of these complex-conjugate arcs and the closed curve to the right; an illustrative example is given in Fig. 7, where the arc endpoints occur at $`u=\pm 2i`$ and $`u=(1/9)(4\pm 2\sqrt{5}i)0.444\pm 0.497i`$. As $`q`$ decreases toward 1, the arcs shrink to points at $`u=e^{\pm i\pi /3}`$. At $`q=1`$, $`(_u)_{qn}`$ is an oval (the solution to the degeneracy equation $`\|u(1u)\|=1`$) that (i) crosses the real axis at $`u=(1/2)(1\pm \sqrt{5})`$, i.e. at approximately, 1.618 and $`0.618`$, and (ii) crosses the imaginary axis at $`u=\pm [(1/2)(1+\sqrt{5})]^{1/2}i\pm 0.7862i`$. For $`0<q<1`$, $``$ consists only of two disjoint arcs, as illustrated for the value $`q=1/2`$ in Fig. 8. For $`q=0`$, the locus $`(_u)_{qn}`$ is the circular arc $$u=\frac{1}{2}e^{i\theta },\frac{\pi }{2}\theta \frac{3\pi }{2}$$ (2.2.7) that crosses the real axis at $`u=1/2`$ and has endpoints at $`u=\pm i/2`$. This is qualitatively similar to $`(_u)_{qn}`$ for the $`L_y=2`$ open strip of the square lattice at $`q=0`$, which was another circular arc crossing the negative real axis at $`u=1/3`$ with endpoints at $`u=(1\pm 2\sqrt{2}i)/9`$ . Note that at $`q=0,1`$, one encounters the noncommutativity (1.20); if one sets $`q`$ to either of these values first and then takes $`n\mathrm{}`$, the resultant $`_u=\mathrm{}`$. Our findings for this $`L_y=2`$ open strip of the triangular lattice may be contrasted with those for the $`L_y=2`$ open strip of the square lattice in . In the latter case, $`_u`$ consisted of arcs that never enclosed any regions for real $`q0`$. ## 3 Cyclic and Möbius Strips of the Triangular Lattice ### 3.1 Results for $`Z`$ By either using an iterative application of the deletion-contraction theorem for Tutte polynomials and converting the result to $`Z`$, or by using a transfer matrix method (in which one starts with a $`q^2\times q^2`$ transfer matrix and generalizes to arbitrary $`q`$), one can calculate the partition function for the cyclic and Möbius ladder graphs of arbitrary length, $`Z(G,q,v)`$, $`G=L_m,ML_m`$. We have used both methods as checks on the calculation. Our results have the general form (1.27) with $`N_\lambda =6`$: $$Z(G_m,q,v)=\underset{j=1}{\overset{6}{}}c_{G,j}(\lambda _{G,j}(q,v))^m,G_m=L_m,ML_m$$ (3.1.1) where $$\lambda _{L,j}=\lambda _{ML,j},j=1,\mathrm{},6$$ (3.1.2) We find $$\lambda _{L,1}=v^2$$ (3.1.3) and $$\lambda _{L,5}=\lambda _{S,1},\lambda _{L,6}=\lambda _{S,2}.$$ (3.1.4) The $`\lambda _j`$ for $`j=2,3,4`$, are the solutions of the cubic equation $`\eta ^3v(v^3+8v+4v^2+2q)\eta ^2+v^2(2v^3q+6v^3+8qv+q^2+8v^2+2v^4+6qv^2)\eta `$ (3.1.5) $``$ $`v^4(v+1)^2(v+q)^2=0.`$ (3.1.7) For the cyclic strip the coefficients in eq. (3.1.1) are $$c_{L,1}=q^23q+1$$ (3.1.8) $$c_{L,j}=q1\mathrm{for}j=2,3,4$$ (3.1.9) $$c_{L,j}=1\mathrm{for}j=5,6$$ (3.1.10) For the Möbius strip, we have $$c_{ML,1}=1$$ (3.1.11) and $$c_{ML,j}=1\mathrm{for}j=5,6.$$ (3.1.12) The coefficients $`c_{ML,j}`$, $`j=2,3,4`$ are more complicated and are given by the generating function (see the appendix) by the formula given in . Our exact calculations (cf. eq. (3.1.2)) yield the following general result $$(\{L\})=(\{ML\}).$$ (3.1.13) This is the same result that one of us found for the analogous strip of the square lattice and is in accord with the conclusion that the singular locus is the same for an infinite-length finite-width strip graph for given transverse boundary conditions, independent of the longitudinal boundary condition. Owing to the equality (3.1.13), we shall henceforth, for brevity of notation, refer to both $`(\{L\})`$ and $`(\{ML\})`$ as $`(\{L\})`$ and similarly for specific points on $``$, such as $`q_c(\{L\})=q_c(\{ML\})`$, etc. Our main interest here is in large $`m`$ and the $`m\mathrm{}`$ limit. However, for completeness, we make the following remark. If $`m3`$, then $`L_m`$ is a (proper) graph, but the $`m=1`$ and $`m=2`$ cases requires special consideration; in these cases, $`L_m`$ degenerates and is not a proper graph<sup>1</sup><sup>1</sup>1A proper graph has no multiple edges or loops, where a loop is an edge that connects a vertex to itself. A multigraph may contain multiple edges, but no loops, while a pseudograph may contain both multiple edges and loops .. $`L_2`$ is the multigraph obtained from the complete graph $`K_4`$ <sup>2</sup><sup>2</sup>2The complete graph $`K_p`$ is the graph with $`p`$ vertices each of which is adjacent to all of the other vertices. by doubling two non-adjacent edges (i.e., edges that do not connect to any common vertex). $`L_1`$ is the pseudograph obtained by connecting two vertices with a double edge and adding a loop to each vertex. Our calculation of $`Z(L_m,q,v)`$ and the corresponding Tutte polynomial $`T(L_m,x,y)`$ apply not just for the uniform cases $`m3`$ but also for the special cases $`m=1,2`$ if for $`m=2`$ one includes the multiple edges and for $`m=1`$ the multiple edges and loops in the evaluation of (1.1), (1.2), and (1.5). Note that in the $`T=0`$ case for the antiferromagnet, the resulting partition function, or equivalently, the chromatic polynomial, is not sensitive to multiple edges, i.e. is the same for a graph in which two vertices are connected by one edge or multiple edges; however, the general partition function (Tutte polynomial) is sensitive to multiple edges. The chromatic polynomial is sensitive to loops and vanishes identically when a pseudograph has any loops. ### 3.2 Special values and expansions of $`\lambda `$’s We discuss some special cases. First, for the zero-temperature Potts antiferromagnet, i.e. the case $`a=0`$ ($`v=1`$), the partition functions $`Z(L_m,q,v)`$ and $`Z(ML_m,q,v)`$ reduce, in accordance with the general result (1.9), to the respective chromatic polynomials $`P(L_m,q)`$ and $`P(ML_m,q)`$ with $`\lambda _j`$’s comprised of the four terms 1, $`(q2)^2`$, and $`(1/2)[52q\pm \sqrt{94q}]`$. The two remaining $`\lambda _{L,j}`$’s vanish in this limit. (Since only four $`\lambda _{L,j}`$’s occur in the chromatic polynomial, a different numbering convention was used in than here, where, in general, six occur.) For the infinite-temperature value $`a=1`$, we have $`\lambda _{L,j}=0`$ for $`j=1,2,3,4,6`$, while $`\lambda _{L,5}=q^2`$, so that $`Z(G,q,a=1)=q^{2m}=q^n`$ for $`G=L_m,ML_m`$, in accord with the general result (1.13). For the real interval $`q3`$ and the region $`R_1`$ to which one can analytically continue from this interval (see Fig. 9 below), $`W=q2`$. Hence, $`W=1`$ at $`q=3`$. A technical remark is the following: one can, and it is convenient to, take the $`n\mathrm{}`$ limit with $`m=0`$ mod 3, so that, by eqs. (1.16) and (1.19), $`P=3!`$, so that the limit for $`W`$ exists at $`a=0`$ as well as at $`a0`$. In contrast, if one took $`n\mathrm{}`$ using all positive integer values of $`m`$, then, strictly speaking, the limit for $`W`$ in eq. (1.10) would not exist for $`a=0`$, since $`P`$ would have the nonconvergent values $`6,0,0,6,,..`$ for $`m=3k,3k+1,3k+2,3k+3,..`$ For the Möbius longitudinal boundary condition, no such convenient choice is possible, since $`\chi =4`$ for all $`m`$; here there is a noncommutativity: if one starts with $`q`$ slightly larger than 3, takes $`n\mathrm{}`$ to calculate $`W`$, and then lets $`q3`$, one gets $`W(q=3)=1`$, but if one sets $`q=3`$ first and then lets $`n\mathrm{}`$, one gets $`W(q=3)=0`$. At $`q=0`$ (with appropriate choices of branch cuts) we find that $$\lambda _{L,1}=\lambda _{L,3}=v^2$$ (3.2.1) $$\lambda _{L,2}=\lambda _{L,5}=\frac{v^2}{2}\left[v^2+4v+7+(v+3)\sqrt{v^2+2v+5}\right]$$ (3.2.2) and $$\lambda _{L,4}=\lambda _{L,6}=\frac{v^2}{2}\left[v^2+4v+7(v+3)\sqrt{v^2+2v+5}\right].$$ (3.2.3) Since there are dominant terms that are degenerate, namely $`\lambda _{L,2}=\lambda _{L,5}`$, it follows that $$q=0\mathrm{is}\mathrm{on}_q(\{L\})a.$$ (3.2.4) This was also true of the circuit graph and cyclic and Möbius square strips with $`L_y=2`$ for which the general Potts model partition function (Tutte polynomial) was calculated in . For $`q=0`$, the coefficients $`c_j=1`$ for $`j=1,5,6`$ and $`c_j=1`$ for $`j=2,3,4`$ so that the equal terms cancel each other pairwise, yielding $`Z(L_m,q=0,v)=0`$, in accordance with the general result (1.11). The noncommutativity (1.20) occurs here: $`\mathrm{exp}(f_{nq})=0`$, while $`\|\mathrm{exp}(f_{qn})\|=\|\lambda _{L,5}\|^{1/2}`$. At $`q=1,2`$ we again encounter noncommutativity in the calculation of the free energy. For $`q=1`$, $`f_{nq}=2\mathrm{ln}a=2K`$, while $`f_{qn}`$ depends on which phase one is in for a given value of $`a`$. For $`0<a<a_c(q=1)`$, where, from eq. (3.3.4), $`a_c(q=1)=(1/2)(1+\sqrt{5})0.6180`$, $`f_{qn}=(1/2)\mathrm{ln}\lambda _{L,c}`$, where $`\lambda _{L,c}`$ is the cube root that is dominant in this phase (corresponding to $`(1/2)[52q+\sqrt{94q}]`$ for $`a0`$), while for $`a>a_c(q=1)`$, $`f_{qn}=(1/2)\mathrm{ln}\lambda _{L,5}`$. Similarly, for $`q=2`$, again with an appropriate choice of branch cuts, $$\lambda _{L,1}=\lambda _{L,3}=v^2$$ (3.2.5) $$\lambda _{L,(2,4)}=\frac{v(v+1)}{2}\left[v^2+3v+4\pm \left[v(v+1)(v^2+5v+8)\right]^{1/2}\right]$$ (3.2.6) and $$\lambda _{L,(5,6)}=\frac{(v+1)(v+2)}{2}\left[(v^2+v+2)\pm \left[(v+1)(v+2)(v^2v+2)\right]^{1/2}\right].$$ (3.2.7) For this value, $`q=2`$, the coefficients are $`c_{L,1}=1`$, while $`c_{L,j}=1`$, $`2j6`$; hence, the $`(\lambda _{L,1})^m`$ and $`(\lambda _{L,3})^m`$ terms cancel each other and make no contribution to $`Z`$, which reduces to $$Z(L_m,q=2,v)=\underset{j=2,4,5,6}{}(\lambda _{L,j})^m$$ (3.2.8) Hence also, $`f_{qn}f_{nq}`$ at $`q=2`$. We observe that the $`\lambda _j`$’s have a more symmetric form for $`q=2`$ when expressed in terms of the variable $`u`$: $$\lambda _{L,1,u}=\lambda _{L,3,u}=u^2(1u)^2$$ (3.2.9) $$\lambda _{L,(2,4),u}=\frac{(1u)}{2}\left[1+u+2u^2\pm \left[(1u)(1+3u+4u^2)\right]^{1/2}\right]$$ (3.2.10) and $$\lambda _{L,(5,6),u}=\frac{(1+u)}{2}\left[1u+2u^2\pm \left[(1+u)(13u+4u^2)\right]^{1/2}\right].$$ (3.2.11) One sees that each member of the pair $`\lambda _{L,(2,4),u}`$ is equal to the respective member of the pair $`\lambda _{L,(5,6),u}`$ with the replacement $`uu`$. It follows that $`\|\lambda _{L,2,u}\|=\|\lambda _{L,5,u}\|`$ and $`\|\lambda _{L,4,u}\|=\|\lambda _{L,6,u}\|`$ for $`u`$ pure imaginary. In general, for $`q_+`$, the partition function $`Z(L_m,q,v)`$ for the cyclic width $`L_y=2`$ strip of the triangular lattice is identical to the partition function for the 1D $`q`$-state Potts model with nearest-neighbor and next-nearest-neighbor spin-spin couplings that are equal in magnitude. This equality can be seen easily by redrawing the strip of the triangular lattice as a line with additional couplings between next-nearest-neighbor vertices on this line. In Ref. Tsai and one of us calculated $`Z`$ for the latter model (with, in general, unequal nearest and next-nearest-neighbor spin-spin couplings). Hence, in particular, the $`q=2`$ and $`q=3`$ of eq. (3.1.1) coincide with the results in (section IX of) . In order to study the zero-temperature critical point in the ferromagnetic case and also the properties of the complex-temperature phase diagram, we calculate the $`\lambda _{G,j,u}`$’s corresponding to the $`\lambda _{G,j}`$’s, using eq. (1.30). In the vicinity of the point $`u=0`$ we have $$\lambda _{L,1,u}=u^2(1u)^2$$ (3.2.12) and the Taylor series expansions $$\lambda _{L,2,u}=12u^3+2(q1)u^4+O(u^5)$$ (3.2.13) $$\lambda _{L,3,u}=u^2+O(u^3)$$ (3.2.14) $$\lambda _{L,4,u}=u^2+O(u^3)$$ (3.2.15) $$\lambda _{L,5,u}=1+2(q1)u^3\left[1+u+O(u^2)\right]$$ (3.2.16) $$\lambda _{L,6,u}=u^2+2(q2)u^3+O(u^4).$$ (3.2.17) Hence, at $`u=0`$, $`\lambda _{L,2,u}`$ and $`\lambda _{L,5,u}`$ are dominant and $`\|\lambda _{L,2,u}\|=\|\lambda _{L,5,u}\|`$, so that the point $`u=0`$ is on $`_u`$ for any $`q0,1`$, where the noncommutativity (1.20) occurs. To determine the angles at which the branches of $`_u`$ cross each other at $`u=0`$, we write $`u`$ in polar coordinates as $`u=re^{i\theta }`$, expand the degeneracy equation $`\|\lambda _{L,2,u}\|=\|\lambda _{L,5,u}\|`$, for small $`r`$, and obtain $`qr^3\mathrm{cos}(3\theta )=0`$, which implies that (for $`q0,1`$) in the limit as $`r=\|u\|0`$, $$\theta =\frac{(2j+1)\pi }{6},j=0,1,\mathrm{},5$$ (3.2.18) or equivalently, $`\theta =\pm \pi /6`$, $`\pm \pi /2`$, and $`\theta =\pm 5\pi /6`$. Hence there are six curves forming three branches of $`_u`$ intersecting at $`u=0`$ and successive branches cross at an angle of $`\pi /3`$. The point $`u=0`$ is thus a multiple point on the algebraic curve $`_u`$, in the technical terminology of algebraic geometry (i.e., a point where several branches of an algebraic curve cross ). In the vicinity of the origin, $`u=0`$, these branches define six complex-temperature phases: the paramagnetic (PM) phase for $`\pi /6\theta \pi /6`$, together with the phases $`O_j`$ for $`1j5`$, with $`O_j`$ occupying the sector $`(2j1)\pi /6\theta (2j+1)\pi /6`$. Note that $`O_3=O_3^{}`$, $`O_4=O_2^{}`$, and $`O_5=O_1^{}`$. For the case of interest here, namely, $`q>0`$, $`\lambda _{L,5,u}`$ is dominant in the PM phase and in the $`O_2`$ and $`O_2^{}`$ phases, while $`\lambda _{L,2,u}`$ is dominant in the $`O_1`$, $`O_1^{}`$, and $`O_3`$ phases. One of our interesting findings is that for $`q=2`$ and for $`q=3`$ the Potts antiferromagnet on the infinite-length, width $`L_y=2`$ strip of the triangular lattice has a zero-temperature critical point. (As must be the case for this to be physical, this is independent of the longitudinal boundary conditions.) In the $`q=2`$ Ising case, this involves frustration, and the 2D Ising antiferromagnet on the triangular lattice also has a $`T=0`$ critical point . In contrast, for the $`q=3`$ case, the $`T=0`$ critical point does not involve any frustration, and the $`q`$ value at which this occurs for our $`L_y=2`$ strip is different than the value, $`q=4`$, where the Potts antiferromagnet has a $`T=0`$ critical point on the full 2D triangular lattice. In order to study the $`T=0`$ critical point for the $`L_y=2`$ strip for these two values of $`q`$, it is useful to calculate expansions of the $`\lambda _j`$’s; only $`\lambda _{L,5}`$ and $`\lambda _{L,2}`$ are necessary for physical thermodynamic properties, while the full set of $`\lambda _{L,j}`$, $`j=1,2,..,6`$ is necessary for the study of the singular locus $``$. For $`q=2`$, besides the exact expressions $`\lambda _{L,1}=\lambda _{L,3}=(1a)^2`$ from eq. (3.2.5), we have the expansions $$\lambda _{L,(2,4)}=a+\frac{a^2}{2}(1+a^2)\pm i\sqrt{a}\left[a+\frac{9}{8}a^2+O(a^3)\right]$$ (3.2.19) and $$\lambda _{L,(5,6)}=a+\frac{a^2}{2}(1+a^2)\pm \sqrt{a}\left[a+\frac{9}{8}a^2+O(a^3)\right].$$ (3.2.20) Note that (i) these are not Taylor series expansions in $`a`$, but rather in the variable $`\sqrt{a}`$, and (ii) each member of the pair $`\lambda _{L,(2,4)}`$ is equal to the respective member of the pair $`\lambda _{L,(5,6)}`$ with the replacement $`aa`$. As shown above, for $`f_{nq}`$ and $`_{nq}`$, where one sets $`q=2`$ first and then takes $`n\mathrm{}`$, $`\lambda _{L,j}`$, $`j=1,3`$, make no contribution, and $`_{nq}`$ is determined by the degeneracy in magnitude of $`\lambda _{L,j}`$, $`j=2,4,5,6`$. From the expansions (3.2.19) and (3.2.20), it follows that in the neighborhood of the point $`a=0`$, $`(_a)_{nq}`$ is determined by the equation $`\|\lambda _{L,2}\|=\|\lambda _{L,5}\|`$ and is a vertical line segment along the imaginary $`a`$ axis. This will be discussed further below (see Fig. 16). For $`q=3`$, besides the exact expression $`\lambda _{L,1}=(1a)^2`$, we calculate the expansions $$\lambda _{L,(2,3)}=e^{\pm 2\pi i/3}2e^{\pm \pi i/3}a+\left(1\pm \frac{2i\sqrt{3}}{3}\right)a^2+O(a^3)$$ (3.2.21) (where the upper (lower) sign applies for $`j=2`$ ($`j=3`$)); $$\lambda _{L,4}=4a^24a^3+O(a^4)$$ (3.2.22) $$\lambda _{L,5}=1+6a3a^2+O(a^3)$$ (3.2.23) and $$\lambda _{L,6}=4a^228a^3+O(a^4).$$ (3.2.24) There are thus four dominant $`\lambda _{L,j}`$’s at $`a=0`$, viz., those with $`j=1,2,3,5`$, which satisfy $`\|\lambda _{L,j}\|=1`$. Writing $`a=\rho e^{i\varphi }`$ and expanding the dominant $`\lambda _{L,j}`$’s in the neighborhood of $`a=0`$, we obtain $$\|\lambda _{L,1}\|^2=14\rho \mathrm{cos}\varphi +O(\rho ^2)$$ (3.2.25) $$\|\lambda _{L,(2,3)}\|^2=14\rho \mathrm{cos}(\varphi \pi /3)+O(\rho ^2)$$ (3.2.26) and $$\|\lambda _{L,5}\|^2=1+12\rho \mathrm{cos}\varphi +O(\rho ^2).$$ (3.2.27) Equating dominant $`\|\lambda _{L,j}\|`$’s, it follows that in the neighborhood of $`a=0`$, $`_a`$ consists of four curves, forming complex-conjugate pairs, passing through $`a=0`$ at the angles $$\varphi =\pm arctan\left(\frac{7\sqrt{3}}{3}\right)76.10^{}.$$ (3.2.28) and $$\varphi =\pm \frac{5\pi }{6}=\pm 150^{}$$ (3.2.29) In Fig. 18 below we show $`_a`$ for this case. ### 3.3 $`_q(\{L\})`$ for fixed $`a`$ #### 3.3.1 Antiferromagnetic case $`0a1`$ We begin with the case $`a=0`$, i.e. the $`T=0`$ limit of the Potts antiferromagnet. The locus $`_q`$ for this case was stated in and given as Fig. 3 in . This locus $`_q`$, shown as Fig. 9 here, separates the complex $`q`$ plane into three different regions: (i) $`R_1`$, including the segments $`q>3`$ and $`q<0`$ on the real axis, where $`W=(q2)`$; (ii) $`R_2`$, including the real interval $`2<q<3`$, in which $`\|W\|=1`$; and (iii) $`R_3`$, including the real segment $`0<q<2`$, in which $`\|W\|=[(1/2)(52q+\sqrt{94q})]^{1/2}`$. At $`q=q_c`$ there are actually four terms that are degenerate in magnitude, $$\lambda _{L,1}=\lambda _{L,5}=\|\lambda _{L,2}\|=\|\lambda _{L,3}\|=1\mathrm{at}q=3\mathrm{for}a=0$$ (3.3.1) corresponding to the property that this is a multiple point on $`_q`$ (in the terminology of algebraic geometry) where four curves intersect. In contrast, the other two points at which $`_q`$ crosses the real axis, at $`q=0`$ and $`q=2`$, involve only the minimal number (two) of degenerate magnitudes and are hence not multiple points. Evidently, $`q_c(\{L\})=3`$ for this $`a=0`$ case. We show calculations of $`_q(\{L\})`$ in Figs. 10, 11, and 12 for finite temperature. As $`a`$ increases from 0, rather than having a fourfold degeneracy of $`\lambda `$’s at $`q=3`$, as in the $`a=0`$ case, eq. (3.3.1), one only has a twofold degeneracy, $$\|\lambda _{L,5}\|=\|\lambda _{L,1}\|$$ (3.3.2) This occurs at the point $$q_c(\{L\})=q_c(\{ML\})=\frac{(1a)(3+2a)}{1+a}=\frac{v(2v+5)}{v+2}.$$ (3.3.3) Corresponding to this value of $`q`$ is the pair $$a_c(\{L\})_\pm =a_c(\{ML\})_\pm =\frac{1}{4}\left[(q+1)\pm \sqrt{q^26q+25}\right].$$ (3.3.4) As $`a`$ increases from 0 to 1, $`q_c(\{L\})`$ decreases monotonically from 3 to 0. Over the interval $`0<a<a_d`$, where $$a_d=\frac{1}{4}(3+\sqrt{17})=0.280776..$$ (3.3.5) the innermost region, $`R_2`$ continues to exist, but shrinks in size. Its left-hand boundary on the real $`q`$ axis, where it is contiguous with region $`R_3`$, continues to lie at $$q_{R_2R_3}=2\mathrm{for}0<a<a_d.$$ (3.3.6) This follows from the fact that this point is the real solution to the degeneracy equation of leading terms $`\|\lambda _{L,1}\|=\|\lambda _{L,3}\|`$, and these are both equal (to $`v^2`$) at $`q=2`$. As $`a`$ increases through the value $`a=a_d`$, $`q_c(\{L\})_{a=a_d}=2`$, so this innermost region $`R_2`$ disappears; this can be seen from the fact that its right-hand boundary point $`q_c`$ becomes equal to its left-hand boundary point, given by $`q_{R_2R_3}`$. Thus, for $`a_d<a<1`$, the locus $`_q`$ separates the $`q`$ plane into only two regions, $`R_1`$ and $`R_3`$. The Potts antiferromagnet has a phase transition at the temperature given by the upper choice of the sign in eq. (3.3.4), i.e., $$T_{p,L}=\frac{J}{k_B\mathrm{ln}(a_c(\{L\})_+)},0<q<3$$ (3.3.7) (where both $`J`$ and the log are negative). However, just as was found for the analogous phase transition of the Potts model on the infinite-length $`L_y=2`$ strip of the square lattice (which occurred for the range $`0<q<2`$ except for $`f_{nq}`$ for the integral value $`q=1`$) , the present phase transition involves unphysical properties, including negative $`Z`$, negative specific heat in the low-temperature phase, and non-existence of a thermodynamic limit independent of boundary conditions in the low-temperature phase. The general existence of such pathologies was noted in . As $`a`$ increases further to 1, $`_q`$ contracts in to a point at the origin, $`q=0`$. #### 3.3.2 Ferromagnetic range $`a1`$ For the Potts ferromagnet, as $`T`$ decreases from infinity, i.e. $`a`$ increases above 1, the locus $`_q`$ forms a lima-bean shaped curve shown for a typical value, $`a=2`$, in Fig. 13. Besides the generally present crossing at $`q=0`$, the point $`q_c(\{L\})`$ at which $`_q`$ crosses the real $`q`$ axis now occurs at negative $`q`$ values. As was true of the model on the analogous width $`L_y=2`$ cyclic and Möbius strips of the square lattice , for physical temperatures, the locus $`_q`$ for the Potts ferromagnet does not cross the positive real $`q`$ axis. Note that this locus does have some support in the $`Re(q)>0`$ half plane, away from the real axis, which was also true of the analogous locus for the $`L_y=2`$ cyclic and Möbius square strip. ### 3.4 $`_u(\{L\})`$ for Fixed $`q`$ We next proceed to the slices of $``$ in the plane defined by the temperature Boltzmann variable $`u`$, for given values of $`q`$, starting with large $`q`$. In the limit $`q\mathrm{}`$, the locus $`_u(\{L\})`$ is reduced to $`\mathrm{}`$. This follows because for large $`q`$, there is only a single dominant $`\lambda _j`$, namely $$\lambda _{L,5}q^2+4qv+O(1)\mathrm{as}q\mathrm{}.$$ (3.4.1) Note that in this case, one gets the same result whether one takes $`q\mathrm{}`$ first and then $`n=2m\mathrm{}`$, or $`n\mathrm{}`$ and then $`q\mathrm{}`$, so that these limits commute as regards the determination of $`_u`$. We first consider values of $`q0,1,2`$, so that no noncommutativity occurs, and $`(_u)_{nq}=(_u)_{qn}_u`$. As discussed above, it is convenient to use the $`u`$ plane since $`_u`$ is compact in this plane, except for the cases $`q=2`$ and $`q=3`$, whereas $`_u`$ is noncompact because of the antiferromagnetic zero-temperature critical point at $`a=1/u=0`$. Extending the discussion in to the case of the strip of the triangular lattice, we observe that the property that the singular locus $`_u`$ passes through the $`T=0`$ point $`u=0`$ for the Potts model with periodic boundary but not with free boundary conditions means that the use of periodic boundary conditions yields a singular locus that manifestly incorporates the zero-temperature critical point, while this is not manifest in $`_u`$ when calculated using free boundary conditions. For $`q=10`$, the locus $`_u`$ is shown in Fig. 14. Six curves on $`_u`$ run into the origin, $`u=0`$, at the angles given in general in eq. (3.2.18). The six corresponding complex-temperature phases contiguous to the origin, $`u=0`$ exhaust the totality of such phases; i.e., there are no such phases that are disjoint from $`u=0`$. The $`\lambda _{L,j}`$’s that are dominant in these phases were given above, after eq. (3.2.18). As is evident in Fig. 14, part of $`_u`$ forms a approximately circular curve. The locus $`_u`$ also includes a line segment on the negative real $`u`$ axis along which $`\lambda _{L,5}`$ and $`\lambda _{L,6}`$ are dominant and are equal in magnitude as complex conjugates of each other. For $`q=2`$, the locus $`(_u)_{nq}`$ is shown in Fig. 15. One sees that, in addition to the six curves intersecting at the ferromagnetic zero-temperature critical point $`u=0`$, $`_u`$ has multiple points at $`u=\pm i`$ where four curves intersect. The complex-temperature phases in the vicinity of $`u=0`$ were determined above after eq. (3.2.18), and these exhaust all complex-temperature phases, i.e. there is none that does not extend in to $`u=0`$. For $`q=3`$, the locus $`_u`$ was given as Fig. 10 in Ref. and is shown with associated partition function zeros in Fig. 17. In this case, in addition to the six phases that are contiguous at $`u=0`$, there is evidently another $`O`$ phase that includes the negative real axis for $`u<1`$, and a complex conjugate pair of $`O`$ phases extending outward from the intersection points at $`u=e^{\pm 2\pi i/3}`$ toward the upper and lower left. At these intersection (multiple) points, six curves on $`_u`$ intersect, just as was true at $`u=0`$. The existence of the intersection points on $`_u`$ at $`u=e^{\pm 2\pi i/3}`$ for the strip would suggest that such points could also be present on the locus $`_u`$ for the $`q=3`$ Potts model on the full triangular lattice. Since this model has not been exactly solved, one can only try to infer $`_u`$ from complex-temperature (Fisher) zeros of the partition function calculated for large triangular lattices . These are consistent with this possibility (see, e.g., Fig. 1 of or Figs. 5-7 of ) although the zeros show such a high degree of scatter in the $`Re(a)<0`$ half-plane that one cannot draw a very decisive conclusion from them. For $`q=4`$ we show the locus $`_u`$ in Fig. 19. In this case, we remark, in For the study of the zero-temperature critical point of the Potts antiferromagnet on the $`L_y=2`$ cyclic and Möbius strips of the triangular lattice, it is useful to display the singular locus $`_a`$ in the $`a`$ plane, since the critical point occurs at $`a=0`$. We show these plots for $`q=2`$ and 3 in Figs. 16 and 18. For $`q=2`$, the boundary $`_a`$ has a multiple point at $`a=\pm i`$ where four branches intersect. For $`q=3`$, four curves on $`_a`$ pass through $`a=0`$ at the angles given in eqs. (3.2.28) and (3.2.29). Note that that the curves that leave the origin at the angles $`\varphi `$ given in eq. (3.2.28) rapidly bend toward the vertical and then back toward the left, so that the complex-temperature phase that is contiguous with $`u=0`$ and lies in the angular wedges between $`76^{}`$ and $`150^{}`$, and its complex-conjugate, are rather narrow. The multiple points at $`a=e^{\pm 2\pi i/3}`$ have been described above. In Figs. 21 and 20 we show the analogous loci $`_u`$ for the $`L_y=2`$ cyclic or Möbius strip of the square lattice studied in for comparison. A particularly striking difference is that, since the $`q=3`$ Potts antiferromagnet is (is not) critical on the $`L_y=2`$ strip of the triangular (square) lattice, the resultant locus $`_u`$ passes through (does not pass through) $`1/u=a=0`$, respectively. The Ising model, $`q=2`$ has both ferromagnetic and antiferromagnetic $`T=0`$ critical points on both the $`L_y=2`$ square and triangular lattice strips, and hence for both strips, $`_u`$ passes through $`1/u=a=0`$. ### 3.5 Connections Between $``$ for Strips and 2D Lattices In earlier work , it was shown that although the physical thermodynamic properties of a discrete spin model are, in general, different in 1D or infinite-length, finite-width strips, which are quasi-1D systems, and in higher dimensions, nevertheless, exact solutions for $`_u`$ in 1D and quasi-1D systems can give insight into $`_u`$ in 2D. This was shown, in particular, for the $`q=2`$ Ising special case of the Potts model, where the comparison can be made rigorously since the model is exactly solvable in 2D. It will often be convenient below to use the equivalent locus in the $`a=1/u`$ plane, $`_a`$. In it was also noted how exact calculations of $`_q`$ and $`W`$ on infinite-length, finite-width strips can give information about the behavior of these quantities on 2D lattices. In Ref. , the connection between $`_a`$ on strips and in on the full 2D square lattice was studied; here we shall do this for the triangular lattice. Again, it is natural to start with the $`q=2`$ Ising case, where one has exact results both for the strip and the full triangular lattice. The comparisons made in and here enable one to formulate a reasonably systematic procedure for making transformations on an exactly calculated locus $`_a`$ for the Potts model on an infinite-length, finite-width strip of a given lattice, in order to construct at least a qualitatively correct locus $`_a`$ on the corresponding 2D lattice. As emphasized in , this represents a new and powerful way of obtaining information about complex-temperature phase diagrams of spin models in 2D (indeed perhaps also higher dimensions) from exact results on strips. It will be recalled that the conventional approach has been the rather laborious procedure of performing exact calculations of the partition function on finite 2D lattices of various sizes with various boundary conditions -. Although some features of $`_a`$ for the Potts model, such as the circle $`\|a1\|=\sqrt{q}`$ for $`Re(a)>0`$ are evident with this conventional procedure, the comlex-temperature (Fisher) zeros often exhibit such considerable scatter for $`Re(a)<0`$ that it is difficult to infer the structure of the locus $`_a`$ in this region. One method was to combine calculations of the zeros with dlog Padé and differential approximant analyses of low-temperature series expansions so as to localize accurately certain points on the complex-temperature boundary -. This method was motivated by the fact that complex-temperature singularities in such quantities as specific heat and magnetization for the 2D Ising model can be calculated exactly and they occur at known points on the complex-temperature boundaries for respective 2D lattices; it was also shown that this was true for the susceptibility, where no exact calculation exists . However, while this method avoids the problem of the scatter in Fisher zeros, it only localizes certain points on $`_a`$. The present method is complementary in that it enables one to gain, at least qualitatively, an idea of the global structure of $`_a`$. We proceed with our exact comparison for the $`q=2`$ Ising case and start by recalling the situation for the square lattice . The locus $`_u`$ for the square lattice consists of the well-known union of circles $`\|u1\|=\sqrt{2}`$ and $`\|u+1\|=\sqrt{2}`$.<sup>3</sup><sup>3</sup>3In Refs. the variable $`z=e^{2K_{Ising}}e^{K_{Potts}}`$ was used for what we denote as $`u`$ here and the variable $`u`$ was used for $`e^{4K_{Ising}}e^{2K_{Potts}}`$, i.e. for what would be denoted $`u^2`$ here. Thus, the circles $`\|u\pm 1\|=\sqrt{2}`$ map to the limaçon of Pascal in the $`u^2`$ plane . The procedure for transforming the locus $`_u`$ found for the infinite-length, width $`L_y=2`$ strip in (see Fig. 20) involves two main steps. First, one retracts each of the curves going through the origin $`u=0`$ so that they no longer pass through this point. It is clear that this retracting is necessary not just for $`q=2`$ but for the general $`q`$-state Potts ferromagnet since the (reduced) free energy and magnetization of this model are analytic in the neighborhood of $`u=0`$, i.e., they have low-temperature series expansions a finite radius of convergence. Given the inversion symmetry $`_u=_a`$ that holds for the square lattice and strips of this lattice (but not for the triangular lattice), this retracting means that the curves are also pulled back from the origin of the $`a`$ plane. The second step in the procedure is to incorporate the feature that the 2D Potts ferromagnet has a finite-temperature phase transition. To build in this feature, one takes the two complex-conjugate ends of the curves with $`Re(u)`$ small and positive that have been pulled back from the origin $`u=0`$ and connects them so that they cross the positive real $`u`$ axis at a point $`u_{PMFM}`$ in the interval $`0<u<1`$; by the inversion symmetry this has the effect of connecting the two other ends in the $`Re(u)>0`$ half plane and having them cross the real axis at $`u_{PMAFM}=u_{PMFM}^1`$. The third step is to build in the feature that if a lattice has an even coordination number, then the Ising model boundary $`_u`$ is symmetric under $`uu`$; in the present case, the locus $`_u`$ for the $`L_y=2`$ strip does not have this property because the coordination number of the vertices on the strip is 3, while the locus $`_u`$ for the square lattice does have the property. We thus connect the curves in the left-hand half plane $`Re(u)<0`$ in such a way as to have this property. As discussed in , the intersection points at $`u=\pm i`$ are the same for both the strip and the square lattice, and, indeed, on wider strips. Evidently, each step of this procedure is based on fundamental principles, and, as long as one limits oneself to obtaining the qualitative features of the locus $`_u`$ for the 2D case from the $`L_y=2`$ strip, nothing is ad hoc. Of course, one cannot predict the actual values of the critical point $`u_{PMFM}=u_{PMAFM}^1`$, but this is was not the goal; there are powerful ways of determining the critical point via series analyses even in cases where a spin model cannot be solved exactly in 2D. (Note that for the purpose of obtaining information about $`_u`$ for a model in 2D, it can be advantageous to use strips with width $`L_y=2`$ rather than wider strips, since the locus $`_u`$ becomes progressively more complicated as $`L_y`$ increases, with more curves passing through $`u=0`$ .) Next, we demonstrate the corresponding comparison with the locus $`_u`$ for the $`q=2`$ Ising special case of our new results for the Potts model on the $`L_y=2`$ strip of the triangular lattice and $`_u`$ for the model on the full 2D triangular lattice. Our calculations are shown in Figs. 15 and 16. For the Ising model on the full 2D triangular lattice, $`_u`$ consists of the union of an oval curve that crosses the real (imaginary) $`u`$ axis at $`\pm i`$ with the line segment $`1/\sqrt{3}Im(u)<\mathrm{}`$ and its complex conjugate on the imaginary $`u`$ axis . Note that, just as was true in the case of the square strips and square lattice (and the honeycomb lattice ), the locus $`_u`$ for the Ising model on both the current strip and the full 2D triangular lattice has intersection points at $`u=\pm i`$. (In terms of the variable denoted $`u^2`$ in our current notation (and $`u`$ in the notation of ), the Ising locus for the triangular lattice is the union of the semi-infinite line segment $`\mathrm{}<u^2<1`$ and the circle $`\|u^2+(1/3)\|=2/3`$, and the above two intersection points map to the single intersection point at $`u^2=1`$.) Following the same procedure as for the square strip, the first step is to pull back each of the curves (of which there are now six rather than four) from the origin. On the imaginary axis, this retracting produces two complex-conjugate semi-infinite line segments, while for the other four curves, it allows one to connect them smoothly, via step 2, to form the oval. Step 3 is not necessary here since both the cyclic $`L_y=2`$ strip of the triangular lattice and the full triangular lattice have the property that each vertex has even coordination number ($`\mathrm{\Delta }=4`$ for the strip and $`\mathrm{\Delta }=6`$ for the 2D lattice) so that $`_u`$ is invariant under the replacement $`uu`$ . It is also interesting to observe that the similarities between the Ising model on the current strip and on the full 2D lattice are much stronger for the antiferromagnet than for the ferromagnet. In contrast to the case of the Ising ferromagnet, which has a finite-temperature critical point on 2D lattices but only a zero-temperature critical point on 1D and quasi-1D lattices, the Ising antiferromagnet has a zero-temperature critical point on both the full 2D triangular lattice and on the infinite-length, width $`L_y=2`$ strip of the triangular lattice, as a consequence of the frustration in this model. Thus in this case, there are not just similarities in the complex-temperature properties of the model, but in the actual physical thermodynamics also. Concerning the complex-temperature phase diagram, in the vicinity of this zero-temperature critical point, for the full 2D triangular lattice, $`_a`$ is just the inverse image of the locus $`_u`$ given above, namely the union of an oval intersecting the real (imaginary) $`a`$ axes at $`\pm \sqrt{3}`$ and $`\pm i`$, and a finite line segment on the imaginary axis stretching from $`a=\sqrt{3}i`$ to $`a=\sqrt{3}i`$. In the neighborhood of the origin, $`a=0`$, the locus $`_a`$ for the present strip is exactly the same, as can be seen in Fig. 16. Having shown the correspondences between the complex-temperature boundaries $`_u`$ and the associated complex-temperature phase diagrams for the exactly solved Ising case, we now use this result as a tool to suggest features of $``$ for the 2D Potts model for other values of $`q`$, where it has not been exactly solved. Since the complex-temperature zeros have usually been presented in the $`a`$ plane, we follow this convention here. The case $`q=3`$ on strips of the square lattice and on the full square lattice was discussed in . In Figs. 21 and 22 we show $`_u`$ and $`_a`$ for the $`L_y=2`$ strip of the square lattice. First, the fact that $`_a`$ crosses the real axis at $`a=2`$ for the strip is the same as is true for the square lattice ; in the latter case, this follows from the existence of the zero-temperature critical point for the $`q=3`$ antiferromagnet, the resultant fact that the singular locus $`_a`$ passes through $`a=0`$, and the duality of the model, according to which if a point $`a`$ is on $`_a`$, then the dual image $`a_d=(q+a1)/(a1)`$ is also on $`_a`$. Second, the existence of the intersection points at $`a=e^{\pm 2\pi i/3}`$ on $`_a`$ for the strip suggests that these also occur on $`_a`$ for the square lattice. We now extend these two comments to explore the global structure of $`_a`$. Since we know that the 2D $`q=3`$ Potts model on the square lattice is analytic in the neighborhood of the $`T=0`$ ferromagnetic point $`u=0`$, our first step is to pull back the four curves that pass through $`u=0`$ in Fig. 21 (equivalently, off to infinity in Fig. 22). Since we know that the 2D model has a finite-temperature critical point (the exact value, $`a=1\|sqrt3`$ is not crucial here), we connect the two ends of the curves in the $`Re(a)>0`$ half-plane together in Fig. 22 so that they cross the positive real $`a`$ plane, making the ferromagnetic critical point $`a_{PMFM}>1`$. The ends of the two curves pointing in the upper left and lower left directions in the $`Re(a)<0`$ half-plane are left as is, in accordance with the finding from series analyses that there are such singularities at $`a_e=1.71(1)\pm 1.46(q)i`$, consistent with lying on endpoints of curves on $`_a`$. Since we know that the $`q=3`$ Potts antiferromagnet has a $`T=0`$ critical point on the square lattice, and hence $`_a`$ passes through $`a=0`$, we move the four curves that pass through the point $`a=1/2`$ (forming an intersection point) over to the right, so that at least one branch passes through $`a=0`$; the exact solution for the strip suggests that $`_a`$ for the square lattice may also have an intersection (multiple) point at $`a=0`$. The complex-temperature zeros for the $`q=3`$ Potts model that have been calculated for patches of the square lattice are consistent with this suggestion (see, e.g., Fig. 1 of ). A third suggestion is that the intersection points on $`_a`$ at $`a1.85\pm 0.85i`$ for the strip have analogues on the locus $`_a`$ for the full square lattice. This suggestion is consistent with the patterns of complex-temperature zeros that have been calculated for patches of the square lattice, but the scatter is too great to provide a strong test. We next use our new results on the $`q=3`$ Potts model on the $`L_y=2`$ strip of the triangular lattice. ### 3.6 Thermodynamics of the Potts Model on the $`L_y=2`$ Strip of the Triangular Lattice #### 3.6.1 Ferromagnetic Case The Potts ferromagnet (with real $`q>0`$) on an arbitrary graph has $`v>0`$ so, as is clear from eq. (1.5), the partition function satisfies the constraint of positivity. In contrast, the specific heat $`C`$ is positive for the model on the (infinite-length limit of the) $`L_y=2`$ triangular strip is positive if and only if $`q>1`$ (for any choice of longitudinal boundary conditions). For $`q=1`$, $`f_{nq}=2\mathrm{ln}a=2K`$ and $`C`$ vanishes identically. Since a negative specific heat is unphysical, we therefore restrict to real $`q1`$. For general $`q`$ in this range, the reduced free energy is given for all temperatures by $`f=(1/2)\mathrm{ln}\lambda _{S,1}`$ as in (2.12) (independent of the different longitudinal boundary conditions, as must be true for the thermodynamic limit to exist). Recall that $`\lambda _{S,1}\lambda _{L,5}`$. It is straightforward to obtain the internal energy $`U`$ and specific heat from this free energy; since the expressions are somewhat complicated, we do not list them here. We show a plot of the specific heat (with $`k_B=1`$) in Fig. 23. One can observe that the value of the maximum is a monotonically increasing function of $`q`$. The high-temperature expansion of $`U`$ is $$U=\frac{2J}{q}\left[1+\frac{(q1)}{q}K+O(K^2)\right].$$ (3.6.1) For the specific heat we have $$C=\frac{2k_B(q1)K^2}{q^2}\left[1+\frac{(q+1)}{q}K+O(K^2)\right].$$ (3.6.2) The low-temperature expansions ($`K\mathrm{}`$) are $$U=J\left[2+(q1)e^{3K}\left[3+4e^K+5(q1)e^{2K}+O(e^{3K})\right]\right]\mathrm{as}K\mathrm{}$$ (3.6.3) and $$C=9k_BK^2(q1)e^{3K}\left[1+\frac{16}{9}e^K+\frac{25}{9}(q1)e^{2K}+O(e^{3K})\right]\mathrm{as}K\mathrm{}$$ (3.6.4) In general, the ratio $`\rho `$ of the largest subdominant to the dominant $`\lambda _j`$’s determines the asymptotic decay of the connected spin-spin correlation function and hence the correlation length $$\xi =\frac{1}{\mathrm{ln}\rho }$$ (3.6.5) Since $`\lambda _{L,5}`$ and $`\lambda _{L,2}`$ are the dominant and leading subdominant $`\lambda _j`$’s, respectively, we have $$\rho _{FM}=\frac{\lambda _{L,2}}{\lambda _{L,5}}$$ (3.6.6) and hence for the ferromagnetic zero-temperature critical point we find that the correlation length diverges, as $`T0`$, as $$\xi _{FM}(2q)^1e^{3K},\mathrm{as}K\mathrm{}$$ (3.6.7) #### 3.6.2 Antiferromagnetic Case In this section we first restrict to the real range $`q3`$ and the additional integer value $`q=2`$ (Ising case) where the Potts antiferromagnet exhibits physically acceptable behavior and then consider the remaining interval $`0<q<3`$ where (except for the trivial $`f_{nq}`$ for $`q=1`$) it exhibits unphysical properties. For $`q3`$, the free energy is given for all temperatures by (2.12), as in the ferromagnetic case but with $`J`$ negative, and is the same independent of the different longitudinal boundary conditions, as is necessary for there to exist a thermodynamic limit. We show plots of the specific heat, for several values of $`q`$, for the Potts antiferromagnet on the (infinite-length limit of the) $`L_y=2`$ strip of the triangular lattice in Fig 24. In contrast to the ferromagnetic case, the maxima of $`C`$ do not increase monotonically with $`q`$; they occur at 0.32 for $`q=2`$ and 0.79 for $`q=3`$, after which the values of the maxima decrease with increasing integral $`q`$ (e.g., they are 0.45 and 0.24 for $`q=4`$ and $`q=6`$). The fact that the curve for $`C`$ for the $`q=2`$ (Ising) antiferromagnet exhibits a very broad maximum can be inferred to be a consequence of the frustration that is present in this model. The high-temperature expansions of $`U`$ and $`C`$ are given by (3.6.1) and (3.6.2); more generally, these expansions also apply in the range $`0<q<3`$. As discussed above, the Ising case $`q=2`$ is one of the cases where one must take account of noncommutativity in the definition of the free energy and hence of thermodynamic quantities. If one sets $`q=2`$ first and then takes $`n\mathrm{}`$, then $`f=f_{nq}=(1/2)\mathrm{ln}\lambda _{L,5}\|_{q=2}`$ where $`\lambda _{L,5}\|q=2`$ was given in eq. (3.2.7), and the low-temperature expansions are $$U_{q=2}=\frac{J}{2}\left[1+\frac{1}{2}e^{K/2}+\frac{23}{16}e^{3K/2}2e^{2K}+O(e^{5K/2})\right]\mathrm{as}K\mathrm{}$$ (3.6.8) and $$C_{q=2}=\frac{k_BK^2e^{K/2}}{8}\left[1+\frac{69}{8}e^K16e^{3K/2}+\frac{1215}{128}e^{2K}+O(e^{3K})\right]\mathrm{as}K\mathrm{}.$$ (3.6.9) For the range $`q3`$, the low-temperature expansions are given by $$U=\frac{(J)e^K}{(q2)^2}\left[(2q3)\frac{(2q3)(4q5)}{(q2)^2}e^K+O(e^{2K})\right]\mathrm{as}K\mathrm{}$$ (3.6.10) and $$C=\frac{k_BK^2e^K}{(q2)^2}\left[(2q3)\frac{2(2q3)(4q5)}{(q2)^2}e^K+O(e^{2K})\right]\mathrm{as}K\mathrm{}.$$ (3.6.11) Note that for the antiferromagnetic case, $`U(T=0)=0`$ for $`q3`$, but $`U(T=0)=J/2=+\|J\|/2`$ for $`q=2`$. The vanishing value of $`U`$ at $`T=0`$ for $`q3`$ means that the Potts model can achieve its preferred ground state for this range of $`q`$, while the nonzero value of $`U(T=0)`$ for the Ising antiferromagnet is a consequence of the frustration that is present in this case. Similarly, the fact that the specific heat vanishes less rapidly for the Ising antiferromagnet than for the Potts model with $`q3`$ reflects the frustration that is present in the Ising case, which is not present for $`q3`$. Note that the apparent divergences that occur as $`q2`$ in eqs. (3.6.10) and (3.6.11) are not actually reached here since these expressions apply only in the region $`q3`$ (the discrete integral case $`q=2`$ was dealt with above). For the zero-temperature critical points in the $`q=2`$ and $`q=3`$ Potts antiferromagnet, $$\rho _{AFM,q=2,3}=\frac{\lambda _{L,2}}{\lambda _{L,5}}$$ (3.6.12) Using the respective expansions (3.2.19)-(3.2.20) and (3.2.21)-(3.2.24), we find that the correlation lengths defined as in (3.6.5) diverges, as $`T0`$, as $$\xi _{AFM,q=2}e^{K/2},\mathrm{as}K\mathrm{}$$ (3.6.13) and $$\xi _{AFM,q=3}e^K,\mathrm{as}K\mathrm{}$$ (3.6.14) Next, we consider the range of real nonintegral $`0<q<3`$. The first pathology is that the Potts antiferromagnet on the infinite-length limit of the $`L_y=2`$ triangular strip, defined with case of cyclic or Möbius longitudinal boundary conditions has a phase transition at the temperature $`T_{p,L}`$ given in eq. (3.3.7) for $`0<q<3`$ (except for $`f_{nq}`$ for the integral values$`q=1,2`$), while, in contrast, if one uses free boundary conditions, then there is no phase transition at this temperature and, although there is a phase transition for $`0<q<2`$, it occurs at the temperature $`T_{p,S}`$ given in eq. (2.2.6), which, in general, is not equal to $`T_{p,L}`$. It follows that there is no well-defined thermodynamic limit for the Potts model with $`0<q<3`$ and $`q1,2`$. The Ising case $`q=2`$ has been dealt with in the preceding subsection. Concerning the value $`q=1`$, as discussed earlier, one encounters noncommutativity in defining the free energy. If one takes $`q=1`$ to start with and then $`n\mathrm{}`$, the thermodynamic limit does exist, independent of boundary conditions, and $`f=f_{nq}=2K`$, $`U=2J=2\|J\|`$, and $`C=0`$. If one starts with $`q1`$, takes $`n\mathrm{}`$, calculates $`f_{qn}`$, and then takes $`q1`$, the thermodynamic limit does not exist since the result differs depending on whether one uses free longitudinal boundary conditions or cyclic (equivalently Möbius) longitudinal boundary conditions. Specifically, for this single value $`q=1`$, $`T_{p,S}`$ is equal to $`T_{p,L}`$, having the value given by $`k_BT_{p,S}=J/\mathrm{ln}[(1/2)(1+\sqrt{5})]2.078\|J\|`$; in the high-temperature phase, $`T>T_{p,S}`$, $`f_{qn}=(1/2)\mathrm{ln}\lambda _{L,5}`$, independent of longitudinal boundary conditions, but in the low-temperature phase, $`T<T_{p,S}`$, the expression for $`f_{qn}`$ is different for the open and cyclic (equivalently, Möbius) strips. There are also other unphysical properties, such as a negative specific heat and a negative partition function for certain ranges of temperature. These is similar to what was found for the analogous strip of the square lattice. ## 4 Summary In this paper we have presented exact calculations of the partition function $`Z`$ of the $`q`$-state Potts model and its generalization to real $`q`$, the Potts model, for arbitrary temperature on $`n`$-vertex strip graphs, of width $`L_y=2`$, of the triangular lattice with free, cyclic, and Möbius longitudinal boundary conditions. These partition functions are equivalent to Tutte/Whitney polynomials for these graphs. The free energy is calculated exactly for the infinite-length limit of these ladder graphs and the thermodynamics is discussed. Considering the full generalization to arbitrary complex $`q`$ and temperature, we determine the singular locus $``$ in the corresponding $`^2`$ space, arising as the accumulation set of partition function zeros as $`n\mathrm{}`$. In particular, we study the connection with the $`T=0`$ limit of the Potts antiferromagnet where $``$ reduces to the accumulation set of chromatic zeros. Comparisons are made with our previous exact calculation of Potts model partition functions for the corresponding strips of the square lattice. Our present calculations yield, as special cases, several quantities of graph-theoretic interest, such as the number of spanning trees, spanning forests, etc., which we record. Acknowledgment: The research of R. S. was supported in part at Stony Brook by the U. S. NSF grant PHY-97-22101 and at Brookhaven by the U.S. DOE contract DE-AC02-98CH10886.<sup>4</sup><sup>4</sup>4Accordingly, the U.S. government retains a non-exclusive royalty-free license to publish or reproduce the published form of this contribution or to allow others to do so for U.S. government purposes. ## 5 Appendix ### 5.1 General The formulas relating the Potts model partition function $`Z(G,q,v)`$ and the Tutte polynomial $`T(G,x,y)`$ were given in and hence we shall be brief here. The Tutte polynomial of $`G`$, $`T(G,x,y)`$, is given by - $$T(G,x,y)=\underset{G^{}G}{}(x1)^{k(G^{})k(G)}(y1)^{c(G^{})}$$ (5.1.1) where $`k(G^{})`$, $`e(G^{})`$, and $`n(G^{})=n(G)`$ denote the number of components, edges, and vertices of $`G^{}`$, and $$c(G^{})=e(G^{})+k(G^{})n(G^{})$$ (5.1.2) is the number of independent circuits in $`G^{}`$. As stated in the text, $`k(G)=1`$ for the graphs of interest here. Now let $$x=1+\frac{q}{v}$$ (5.1.3) and $$y=a=v+1$$ (5.1.4) so that $$q=(x1)(y1).$$ (5.1.5) Then $$Z(G,q,v)=(x1)^{k(G)}(y1)^{n(G)}T(G,x,y).$$ (5.1.6) For a planar graph $`G`$ the Tutte polynomial satisfies the duality relation $$T(G,x,y)=T(G^{},y,x)$$ (5.1.7) where $`G^{}`$ is the (planar) dual to $`G`$. As discussed in , the Tutte polynomial for recursively defined graphs comprised of $`m`$ repetitions of some subgraph has the form $$T(G_m,x,y)=\underset{j=1}{\overset{N_\lambda }{}}c_{T,G,j}(\lambda _{T,G,j})^m$$ (5.1.8) There are several special cases of the Tutte polynomial that are of interest. One that we have analyzed in the text and in previous papers is the chromatic polynomial $`P(G,q)`$. This is obtained by setting $`y=0`$, i.e., $`v=1`$, so that $`x=1q`$; the correspondence is $`P(G,q)=(q)^{k(G)}(1)^nT(G,1q,0)`$. A second special case is the flow polynomial $`F(G,q)`$, obtained by setting $`x=0`$ and $`y=1q`$: $`F(G,q)=(1)^{e(G)n(G)+k(G)}T(G,0,1q)`$ For planar $`G`$, given the relation (5.1.7), the flow polynomial is, up to a power of $`q`$, proportional to the chromatic polynomial: $`F(G,q)P(G^{},q)`$. A third special case is the reliability polynomial , which will not be considered here. ### 5.2 Triangular Strip with Free Longitudinal Boundary Conditions The generating function representation for the Tutte polynomial for the open strip of the triangular lattice comprised of $`m+1`$ squares with edges joining the lower-left to upper right vertices of each square, denoted $`S_m`$, is $$\mathrm{\Gamma }_T(S_m,x,y;z)=\underset{m=0}{\overset{\mathrm{}}{}}T(S_m,x,y)z^m.$$ (5.2.1) We have $$\mathrm{\Gamma }_T(S,x,y;z)=\frac{𝒩_T(S,x,y;z)}{𝒟_T(S,x,y;z)}$$ (5.2.2) where $$𝒩_T(S,x,y;z)=A_{T,S,0}+A_{T,S,1}z$$ (5.2.3) with $$A_{T,S,0}=x(x+1)^2+2xy+y(y+1)$$ (5.2.4) $$A_{T,S,1}=x^3y^2$$ (5.2.5) and $`𝒟_T(S,x,y,z)`$ $`=`$ $`1[(x+1)^2+y(y+2)]z+x^2y^2z^2`$ (5.2.6) $`=`$ $`{\displaystyle \underset{j=1}{\overset{2}{}}}(1\lambda _{T,S,j}z)`$ (5.2.8) with $$\lambda _{T,S,(1,2)}=\frac{1}{2}\left[(x+1)^2+y(y+2)\pm (1+x+y)\sqrt{R_T}\right]$$ (5.2.9) where $$R_T=(x+1)^2+(y+1)^212xy.$$ (5.2.10) The corresponding closed-form expression is given by the general formula $$T(S_m,x,y)=\left[\frac{A_{T,S,0}\lambda _{T,S,1}+A_{T,S,1}}{\lambda _{T,S,1}\lambda _{T,S,2}}\right](\lambda _{T,S,1})^m+\left[\frac{A_{T,S,0}\lambda _{T,S,2}+A_{T,S,1}}{\lambda _{T,S,2}\lambda _{T,S,1}}\right](\lambda _{T,S,2})^m.$$ (5.2.11) It is easily checked that this is a symmetric function of the $`\lambda _{S,j}`$, $`j=1,2`$. ### 5.3 Cyclic and Möbius Strips of the Triangular Lattice We have $$T(L_m,x,y)=\underset{j=1}{\overset{6}{}}c_{T,L,j}(\lambda _{T,L,j})^m$$ (5.3.1) and $$T(ML_m,x,y)=\underset{j=1}{\overset{6}{}}c_{T,ML,j}(\lambda _{T,ML,j})^m$$ (5.3.2) It is convenient to extract a common factor from the coefficients: $$c_{T,G,j}\frac{\overline{c}_{T,G,j}}{x1},G=L,ML.$$ (5.3.3) Of course, although the individual terms contributing to the Tutte polynomial are thus rational functions of $`x`$ rather than polynomials in $`x`$, the full Tutte polynomial is a polynomial in both $`x`$ and $`y`$. We have $$\lambda _{T,ML,j}=\lambda _{T,L,j},j=1,\mathrm{},6$$ (5.3.4) $$\lambda _{T,L,1}=1$$ (5.3.5) and $$\lambda _{T,5}=\lambda _{T,S,1},\lambda _{T,6}=\lambda _{T,S,2}$$ (5.3.6) where $`\lambda _{T,S,1}`$ and $`\lambda _{T,S,2}`$ were given in eq. (5.2.9) with (5.2.10). The $`\lambda _{T,L,j}`$, $`j=2,3,4`$, are solutions of the cubic equation $$\xi ^3(3+2x+2y+y^2)\xi ^2+[(x+1)^2+2xy(y+1)]\xi x^2y^2=0$$ (5.3.7) We label these $`\lambda _{T,L,j}`$’s in a manner corresponding to the $`\lambda _{L,j}`$’s occurring in the Potts model partition function, so that $`\lambda _{T,L,5}`$ is dominant in the region in $`(x,y)`$ equivalent to the PM phase, and so forth for the others. We note that $$\lambda _{T,L,2}\lambda _{T,L,3}\lambda _{T,L,4}=\lambda _{T,L,5}\lambda _{T,L,6}=x^2y^2.$$ (5.3.8) The coefficients for the cyclic strip are $$\overline{c}_{T,L,1}=[(x1)(y1)]^23(x1)(y1)+1$$ (5.3.9) $$\overline{c}_{T,L,j}=xyxy\mathrm{for}j=2,3,4$$ (5.3.10) and $$\overline{c}_{T,L,j}=1\mathrm{for}j=5,6.$$ (5.3.11) These are symmetric under interchange of $`xy`$, which is a consequence of the fact that the $`\overline{c}_{L,j}`$ are functions only of $`q`$, which, by eq. (5.1.5) is a symmetric function of $`x`$ and $`y`$. For the Möbius strip, we find $$\overline{c}_{T,ML,1}=1$$ (5.3.12) and $$\overline{c}_{T,ML,5}=\overline{c}_{T,ML,6}=1.$$ (5.3.13) The $`c_{T,ML,j}`$, $`j=2,3,4`$ can be calculated from the generating function using the formula given in . Their form is more complicated, as is evident from the results given in for the corresponding chromatic polynomial for this case. We write the generating function as $$\mathrm{\Gamma }_T(ML,x,y,z)=\underset{m=1}{\overset{\mathrm{}}{}}T(ML_m,x,y)z^{m1}.$$ (5.3.14) with $$\mathrm{\Gamma }_T(ML,x,y,z)=\frac{𝒩_T(ML,x,y,z)}{𝒟_T(ML,x,y,z)}$$ (5.3.15) where $$𝒟_T(ML,x,y,z)=𝒟_T(L,x,y,z)=\underset{j=1}{\overset{6}{}}(1\lambda _{L,j}z)$$ (5.3.16) and $$𝒩_T(ML,x,y,z)=\underset{j=0}{\overset{5}{}}A_{ML,j}z^j$$ (5.3.17) with $$A_{ML,0}=y(x+y+y^2)$$ (5.3.18) $$A_{ML,1}=yx4y^2x4y^33y^45y^3xy^52x^2yx^3yx^2y^2x^2y^3+2y+2x+3x^2+x^3$$ (5.3.19) $`A_{ML,2}`$ $`=`$ $`y(2x^2y+4yx+3x+3y+x^2y^32x^3y+6x^2y^2+2xy^4`$ (5.3.20) $`+`$ $`4y^3x+y^4x^2+3x^3y^2+5y^2+3y^3+y^4+6y^2x+2x^2)`$ (5.3.22) $`A_{ML,3}`$ $`=`$ $`y(2x^2y3yxxyx^4y+3x^2y^32x^3y+8x^2y^2+2xy^4`$ (5.3.23) $`+`$ $`2x^3y^4x^3y^3+5y^3xx^3+x^4y^2+2y^4x^2+4x^3y^22y^2`$ (5.3.25) $``$ $`y^3+3y^2x2x^2)`$ (5.3.27) $$A_{ML,4}=x^2y^2(1+x^2y^3x^2y^2x^2+y^22x+y+2y^3x+x^2y+y^3+yx)$$ (5.3.28) and $$A_{ML,5}=y^4x^4(y1).$$ (5.3.29) ### 5.4 Special Values of Tutte Polynomials for Strips of the Triangular Lattice For a given graph $`G=(V,E)`$, at certain special values of the arguments $`x`$ and $`y`$, the Tutte polynomial $`T(G,x,y)`$ yields quantities of basic graph-theoretic interest -. We recall some definitions: a spanning subgraph was defined at the beginning of the paper; a tree is a connected graph with no cycles; a forest is a graph containing one or more trees; and a spanning tree is a spanning subgraph that is a tree. We recall that the graphs $`G`$ that we consider are connected. Then the number of spanning trees of $`G`$, $`N_{ST}(G)`$, is $$N_{ST}(G)=T(G,1,1),$$ (5.4.1) the number of spanning forests of $`G`$, $`N_{SF}(G)`$, is $$N_{SF}(G)=T(G,2,1),$$ (5.4.2) the number of connected spanning subgraphs of $`G`$, $`N_{CSSG}(G)`$, is $$N_{CSSG}(G)=T(G,1,2),$$ (5.4.3) and the number of spanning subgraphs of $`G`$, $`N_{SSG}(G)`$, is $$N_{SSG}(G)=T(G,2,2).$$ (5.4.4) From the duality relation (5.1.7), one has, for planar graphs $`G`$ and their planar duals $`D^{}`$, $$N_{ST}(G)=N_{ST}(G^{}),N_{SSG}(G)=N_{SSG}(G^{})$$ (5.4.5) and $$N_{SF}(G)=N_{CSSG}(G^{}).$$ (5.4.6) From our calculations of Tutte polynomials, we find that $$N_{ST}(S_m)=N_{ST}(S_m^{})=\left(4+\frac{9\sqrt{5}}{5}\right)\left(\frac{7+3\sqrt{5}}{2}\right)^m+\left(4\frac{9\sqrt{5}}{5}\right)\left(\frac{73\sqrt{5}}{2}\right)^m$$ (5.4.7) $$N_{SF}(S_m)=N_{CSSG}(S_m^{})=\left(12+\frac{17\sqrt{2}}{2}\right)\left[2(3+2\sqrt{2})\right]^m+\left(12\frac{17\sqrt{2}}{2}\right)\left[2(32\sqrt{2})\right]^m$$ (5.4.8) $$N_{CSSG}(S_m)=N_{SF}(S_m^{})=(7+5\sqrt{2})\left[2(3+2\sqrt{2})\right]^m+(75\sqrt{2})\left[2(32\sqrt{2})\right]^m$$ (5.4.9) $$N_{SSG}(S_m)=N_{SSG}(S_m^{})=2^{4m+5}.$$ (5.4.10) For the cyclic $`L_y=2`$ strip of the triangular lattice, $`L_m`$, we first note that for $`m3`$, $`L_m`$ is a (proper) graph, but for $`m=1`$ and $`m=2`$, $`L_m`$ is not a proper graph, but instead, is a multigraph, with multiple edges (and, for $`m=1`$, loops). The following formulas apply for all $`m1`$: $$N_{ST}(L_m)=N_{ST}(L_m^{})=\frac{2m}{5}\left[\left(\frac{7+3\sqrt{5}}{2}\right)^m+\left(\frac{73\sqrt{5}}{2}\right)^m2\right]$$ (5.4.11) $$N_{SF}(L_m)=N_{CSSG}(L_m^{})=1\underset{j=2,3,4}{}(\lambda _{T,j})^m+\left[2(3+2\sqrt{2})\right]^m+\left[2(32\sqrt{2})\right]^m.$$ (5.4.12) where $`\lambda _{T,L,j}`$, $`j=2,3,4`$ are the roots of eq. (5.3.7) for $`x=2,y=1`$, viz., $`\xi ^310\xi ^2+17\xi 4=0`$: $$\lambda _{T,L,2}(x=2,y=1)=0.280176\mathrm{}$$ (5.4.13) $$\lambda _{T,L,3}(x=2,y=1)=1.803442\mathrm{}$$ (5.4.14) $$\lambda _{T,L,4}(x=2,y=1)=7.916382\mathrm{}$$ (5.4.15) We also calculate $`N_{CSSG}(L_m)=N_{SF}(L_m^{})`$ $`=`$ $`(2+{\displaystyle \frac{6m}{7}})+\left[1+m{\displaystyle \frac{(3\sqrt{2})}{7}}\right]\left[2(3+2\sqrt{2})\right]^m`$ (5.4.16) $`+`$ $`\left[1+m{\displaystyle \frac{(3+\sqrt{2})}{7}}\right]\left[2(32\sqrt{2})\right]^m`$ (5.4.18) $$N_{SSG}(L_m)=N_{SSG}(L_m^{})=2^{4m}.$$ (5.4.19) Since $`T(G_m,x,y)`$ grows exponentially as $`m\mathrm{}`$ for the families $`G_m=S_m`$ and $`L_m`$ for $`(x,y)=(1,1)`$, (2,1), (1,2), and (2,2), one defines the corresponding constants $$z_{set}(\{G\})=\underset{n(G)\mathrm{}}{lim}n(G)^1\mathrm{ln}N_{set}(G),set=ST,SF,CSSG,SSG$$ (5.4.20) where, as above, the symbol $`\{G\}`$ denotes the limit of the graph family $`G`$ as $`n(G)\mathrm{}`$ (and the $`z`$ here should not be confused with the auxiliary expansion variable in the generating function (5.2.1) or the Potts partition function $`Z(G,q,v)`$.) General inequalities for these were given in . Our results yield $$z_{ST}(\{G\})=\frac{1}{2}\mathrm{ln}\left(\frac{7+3\sqrt{5}}{2}\right)0.962424\mathrm{for}G=S,L,ML$$ (5.4.21) $$z_{SF}(\{G\})=z_{CSSG}(\{G\})=\frac{1}{2}\mathrm{ln}[2(3+2\sqrt{2})]1.22795\mathrm{for}G=S,L,ML$$ (5.4.22) and $$z_{SSG}(\{G\})=2\mathrm{ln}21.38629\mathrm{for}G=S,L,ML$$ (5.4.23) The cyclic $`L_y=2`$ strip of the triangular lattice is a $`\kappa `$-regular graph <sup>5</sup><sup>5</sup>5A $`\kappa `$-regular graph is a graph in which all of the vertices have the same degree, $`\kappa `$. with $`\kappa =4`$. It is therefore of interest to compare the value of $`z_{ST}`$ that we have obtained for the $`n\mathrm{}`$ limit of this family with an upper bound for $`\kappa `$-regular graphs with $`\kappa 3`$ , viz. $$z_{ST}(\{G\})z_{ST,u.b.}(\{G\}),\mathrm{where}z_{ST,u.b.}(\{G\})=\mathrm{ln}\left[\frac{(\kappa 1)^{\kappa 1}}{[\kappa (\kappa 2)]^{(\kappa /2)1}}\right]$$ (5.4.24) For this purpose we define $`r`$ as the ratio of the left- to the right-hand side of eq. (5.4.24). We have $$r_{ST}(\{L\})=\frac{\frac{1}{2}\mathrm{ln}\left(\frac{7+3\sqrt{5}}{2}\right)}{3\mathrm{ln}\left(\frac{3}{2}\right)}0.791$$ (5.4.25) Another comparison of interest is the ratio of $`z_{ST}`$ for these $`L_y=2`$ strips with $`z_{ST}`$ for the full 2D triangular lattice, which has the value $$z_{ST,tri}=2z_{hc}=\frac{3\sqrt{3}}{\pi }(15^2+7^211^2+13^2\mathrm{})=1.615329736\mathrm{}$$ (5.4.26) $$\frac{z_{ST}(\{L\})}{z_{ST}(tri)}0.596$$ (5.4.27) Finally, we may also compare the $`z`$ values above with those for the infinite-length limit of the $`L_y=2`$ strip of the square lattice (with free transverse boundary conditions). This is $$z_{ST}=\frac{1}{2}\mathrm{ln}(2+\sqrt{3})0.658479$$ (5.4.28) Further, the calculation of the Tutte polynomial for the open and/or cyclic $`L_y=2`$ strip of the square lattice yields $$z_{SF}=\frac{1}{2}\mathrm{ln}[2(2+\sqrt{3})]1.00505$$ (5.4.29) $$z_{CSSG}=\frac{1}{2}\mathrm{ln}\left(\frac{5+\sqrt{17}}{2}\right)0.758832$$ (5.4.30) $$z_{SSG}=\frac{3}{2}\mathrm{ln}2=1.03972$$ (5.4.31) Thus, one observes that each of these four quantities is greater for the (infinite-length limit of the) $`L_y=2`$ strip of the triangular lattice than for the (same limit of the) $`L_y=2`$ strip of the square lattice. For $`z_{ST}`$, this follows from the fact that one can obtain the strip of the triangular lattice from that for the square lattice by uniformly adding a diagonal bond to each of the squares . The inequalities for the other quantities $`z_{SF}`$, $`z_{CSSG}`$, and $`z_{SSG}`$ are presumably a consequence of this fact also. ### 5.5 Tutte Polynomials for Dual Graphs Since the Tutte polynomial satisfies the duality relation (5.1.7) for a planar graph $`G`$, our calculations of the Tutte polynomials $`T(G_m,x,y)`$ for the open and cyclic $`L_y=2`$ strips of the triangular lattice, $`S_m`$ and $`L_m`$, also yield the corresponding Tutte polynomials for the duals of these graphs. For $`m3`$, the dual of $`L_m`$ is a graph with $`n=2m+2`$ vertices comprised of the circuit graph $`C_{2m}`$ together with two additional vertices, which we denote $`v_e`$ and $`v_o`$ such that, if we label the vertices on $`C_{2m}`$ as $`v_j`$, $`j=1,..,2m`$, then each of the $`v_j`$ with $`j`$ even is connected by an edges to $`v_e`$ and each of the $`v_j`$ with $`j`$ odd is connected to $`v_o`$. We shall denote this graph as an “alternating wheel”, $`A_{2m+2}`$. Then $$T(A_{2m+2},x,y)=T(L_m,y,x),\mathrm{for}m3$$ (5.5.1) The cases $`m=1`$ and $`m=2`$ can be dealt with in a similar manner. The dual of the open strip graph $`S_m`$ is the multigraph formed by a line of $`2(m+1)`$ vertices, of which the interior $`2m`$ vertices are connected by single edges to a single additional vertex, and the two end vertices on the line are connected to this additional vertex by double edges. We have $`T((S_m)^{},x,y)=T(S_m,y,x)`$.
warning/0004/hep-ph0004105.html	ar5iv	text	# I Introduction ## I Introduction In the last few years there has been another wave of excitement regarding the question of neutrino masses. This is largely due to the many new experiments testing neutrino oscillations, most notably the positive indications obtained by Super Kamiokande on atmospheric neutrino oscillations . Similar indications come from other experiments . The solar neutrino experiments have for many years provided independent evidence for neutrino oscillations . Accelerator and reactor experiments have also played an important role. They have furnished strict bounds on neutrino oscillation parameters . In the case of the LSND experiment at Los Alamos evidence for $`\overline{\nu }_\mu \overline{\nu }_e`$ oscillation has been reported. See refs for recent reviews. It is hoped that new experimental results can be used to determine the neutrino squared mass differences and mixing angles. In turn, these may help to infer the neutrino mass matrix. This is presumably a possible gateway to a more fundamental theory beyond the standard model. Of course this is a highly speculative area, and even though there are many imaginative proposals , it seems fair to say that the the true answer is essentially unknown. In order to make progress in this direction, it seems useful to investigate various plausible Ansatze for the neutrino mass matrix. From this point of view we propose the Ansatz for the 3 generation neutrino mass matrix, $`M_\nu `$: $$\mathrm{Tr}(M_\nu )=0$$ (1) and investigate its consequences. We are considering the neutrinos to be represented by 2-component spinors so that, in the most general situation, $`M_\nu `$ is an arbitrary symmetric complex matrix. As we will see in section II, Eq. (1.1) can be motivated from an SO(10) grand unified model , in which it may be derived with some assumptions. Physically, Eq. (1.1) corresponds to the well known approximate signature of grand unification that $`\frac{m(b)}{m(\tau )}3`$. Furthermore we will see in sections IV and V that Eq. (1.1) can be straightforwardly combined with experimental information to get an idea of the neutrino masses themselves as well as the “texture” of $`M_\nu `$. Relevant matters of notation are discussed in section III while a summary is presented in section VI. ## II Plausibility of the Ansatz In the SO(10) grand unification model each generation contains one light massive two component neutrino and also a very heavy one which is “integrated out” according to the “seesaw mechanism” . The effective $`3\times 3`$ neutrino mass matrix takes the form: $$M_\nu =M_LM_D^TM_HM_D$$ (2) where $`M_L`$, $`M_H`$ and $`M_D`$ are respectively the mass matrices of the light neutrinos, heavy neutrinos and heavy-light mixing (or “Dirac matrix”). Generally the second, seesaw, term is considered to dominate. Here however we shall assume the first term to be the dominant one. This is necessary for the present derivation of Eq. (1.1) to hold. Also, a rough order of magnitude estimate for the second term would be $`\frac{m(\tau )^2}{10^{16}\mathrm{GeV}}`$ or about $`3\times 10^7`$ eV. Thus, the seesaw term could be negligible if neutrino masses turn out to be appreciably larger than this value. Now in SO(10), Higgs mesons belonging to the 10, 120 and 126 representations can contribute to the fermion masses at tree level. One has for the down quark, charged lepton and light neutrino mass matrices, $`M_d`$ $`=`$ $`aS(10)+bA(120){\displaystyle \frac{1}{3}}cS(126)`$ (3) $`rM_e`$ $`=`$ $`aS(10)+dA(120)+cS(126)`$ (4) $`sM_L`$ $`=`$ $`eS(126)`$ (5) where $`a`$, $`b`$, $`c`$, $`d`$, $`e`$ are numbers representing Higgs meson vacuum values. S(10), A(120) and S(126) are the matrices of the Yukawa type constants which couple the fermions to the 10, 120 and 126 Higgs mesons respectively; the matrices S(10) and S(126) must be symmetric while A(120) is antisymmetric. Finally, $`r3`$ is a renormalization factor for comparing the quark masses with the charged lepton masses at a low energy scale rather than at the grand unified scale; $`s`$ is a similar factor for the neutrino masses. With the stated assumption that the $`M_L`$ term dominates in Eq. (2.1) we get $$\mathrm{Tr}(M_\nu )\mathrm{Tr}(M_d)r\mathrm{Tr}(M_e)$$ (6) which clearly also holds when any number of 10’s or 120’s are present but only a single 126. The matrices appearing in Eq. (2.3) are so far essentially unrestricted complex ones. To proceed, we make the further assumption that the matrices are hermitian. Then $`M_d`$ and $`M_u`$ may each be brought to diagonal form by unitary transformations. Thus the right hand side of Eq. (2.3) may be evaluated to yield approximately, $$\mathrm{Tr}(M_\nu )m(b)rm(\tau )0$$ (7) according to a well known numerical success, based on the observation that $`r3`$, of grand unification . Note that we have not needed to assume that the mass matrix has any zero elements.<sup></sup><sup></sup>*In a similar mechanism was studied for $`M_H`$ where, in addition, a special combined Fritzsch-Stech Ansatz was used. Here we are not making any special Ansatz of this type for the mass matrices. Even if the cancellation on the right hand side of Eq. (2.4) is not perfect, it should still be a good approximation. In an SO(10) model where the mass matrices are hermitian, $`M_\nu `$ will be real symmetric. We will investigate this case and also the possibility that the more general case holds. ## III Some notation Our plan is to combine the Ansatz Eq. (1.1) with experimentally obtained results on neutrino oscillations in order to learn more about $`M_\nu `$ itself. For this purpose it may be helpful to set down our notation for the pieces of the effective $`\mathrm{SU}(2)_L\times \mathrm{U}(1)`$ theory involving neutrinos and to make some related remarks. The free Lagrangian containing three two component massive fields is: $$_{free}=i\rho ^{}\sigma _\mu _\mu \rho \frac{1}{2}(\rho ^T\sigma _2M_\nu \rho +h.c.),$$ (8) where $`M_\nu =M_\nu ^T`$ is the (not yet diagonalized) neutrino mass matrix of the underlying theory to be identified with the matrix in Eq. (1.1). Note that we are free to multiply the first mass term in Eq. (8) by an overall arbitrary phase which is a matter of convention. It is possible to find a unitary matrix $`U`$ which brings $`M_\nu `$ to real, positive, diagonal form in the following way: $$U^TM_\nu U=\widehat{M},\widehat{M}=\mathrm{diag}(m_1,m_2,m_3).$$ (9) The mass diagonal fields $`\nu `$ are then $$\rho =U\nu .$$ (10) Similarly, the column vector of left handed negatively charged leptons in the underlying theory, $`E_L`$ is related to the mass diagonal fields $`e_L`$ by $$E_L=\mathrm{\Omega }e_L,$$ (11) where $`\mathrm{\Omega }^{}=\mathrm{\Omega }^1`$. Combining factors from Eq. (3.3) and Eq. (3.4) we obtain the unitary mixing matrix, $`K`$ for the charged current weak interaction, $$K=\mathrm{\Omega }^{}U.$$ (12) This appears in the Lagrangian term, $$_{int}=\frac{ig}{\sqrt{2}}W_\mu ^{}\overline{e}_L\gamma _\mu K\nu +h.c.,$$ (13) where a conventional four component Dirac notation with $`\gamma _5`$ diagonal is being employed and $`\nu `$ has only the first two components non zero. Next we parameterize $`K`$. It is possible to restrict $`\mathrm{det}K=1`$ by adjusting an overall phase which can be absorbed in $`\overline{e}_L`$. Then we write $$K=\omega _0(\alpha )\omega _{23}(\eta _{23},\varphi _{23})\omega _{12}(\eta _{12},\varphi _{12})\omega _{13}(\eta _{13},\varphi _{13}),$$ (14) where $$\omega _0(\alpha )=\mathrm{diag}(e^{i\alpha _1},e^{i\alpha _2},e^{i\alpha _3}),$$ (15) with $`\alpha _3=(\alpha _1+\alpha _2)`$ and, for example $$\omega _{12}(\eta _{12},\varphi _{12})=\left[\begin{array}{ccc}\mathrm{cos}\eta _{12}& e^{i\varphi _{12}}\mathrm{sin}\eta _{12}& 0\\ e^{i\varphi _{12}}\mathrm{sin}\eta _{12}& \mathrm{cos}\eta _{12}& 0\\ 0& 0& 1\end{array}\right].$$ (16) Eq. (3.7) contains the eight parameters needed to characterize an arbitrary unitary unimodular matrix. From the standpoint of Eq. (3.6) it can be further simplified by using the freedom to rephase $`\overline{e}_L\overline{e}_L\omega _0^1(\alpha )`$ without changing the free part of the charged lepton Lagrangian. On the other hand, the form of the mass terms in Eq. (8) shows that the neutrino fields can not be rephased. Thus a suitable minimal parameterizationThis is written out in detail in Eq (2) of for $`K`$ in (3.6) is $$K=\omega _{23}(\eta _{23},\varphi _{23})\omega _{12}(\eta _{12},\varphi _{12})\omega _{13}(\eta _{13},\varphi _{13}),$$ (17) involving three “angles”, $`\eta _{ab}`$ and three “phases”, $`\varphi _{ab}`$. Note the identity $$\omega _0(\alpha )\omega _{ab}(\eta _{ab},\varphi _{ab})\omega _0^1(\alpha )=\omega _{ab}(\eta _{ab},\alpha _a+\varphi _{ab}\alpha _b).$$ (18) This identity may be used to transfer two of the phases $`\varphi _{ab}`$ in Eq. (3.7) to a diagonal matrix on the right of $`K`$ as, for example, $$K=\omega _{23}(\eta _{23},0)\omega _{12}(\eta _{12},0)\omega _{13}(\eta _{13},\varphi _{13})\omega _0^1(\tau ),$$ (19) where $`\tau _1+\tau _2+\tau _3=0`$, which may be used instead of Eq. (3.10). We also need the formula for the amplitude of neutrino oscillation. For the case when a neutrino, produced by a charged lepton of type $`a`$, “oscillates” to make at time $`t`$, a charged lepton of type $`b`$, we have $$\mathrm{amp}(ab)=\underset{\alpha }{}K_{a\alpha }^{}K_{b\alpha }e^{iE_\alpha t},$$ (20) where the sum goes over the neutrinos of definite mass, $`m_\alpha `$. Inserting the parameterization Eq. (3.12) into Eq. (3.13) shows that the effect of the factor $`\omega _0^1(\tau )`$ cancels out. Thus for ordinary oscillations, $`K`$ is parameterized by three angles and one CP violating phase as for the CKM quark mixing matrix. On the other hand, the two additional CP violating phases $`\tau _1`$ and $`\tau _2`$ show up if one considers neutrino-antineutrino oscillations or neutrinoless double beta decay . The formula for the probability, $`P_{ab}`$ is gotten by taking the squared magnitude of Eq. (3.13) and replacing the exponential factor $`E_\alpha t`$ by $`(E+m_\alpha ^2/(2E))L`$, where $`E`$ is the neutrino energy and $`L`$ is the oscillation distance. For practical reasons it is very important to take account of the experimental uncertainties in $`E`$ and $`L`$. The simplest approximation is to define $`b=L/(4E)`$ and assume that one can smear $`P_{ab}`$ with a Gaussian distribution in $`b`$. $`b_0`$ is defined as the mean value and $`\sigma _b`$ as the standard deviation appropriate to the particular physical setup. Then we find for the smeared probability $`P_{ab}`$ $`=`$ $`\delta _{ab}2{\displaystyle \underset{\alpha <\beta }{}}[\mathrm{Re}(f_{\alpha \beta ab})(1\mathrm{cos}(2b_0m_{\beta \alpha }^2)\mathrm{exp}(2\sigma _b^2(m_{\beta \alpha }^2)^2))`$ (21) $`+`$ $`\mathrm{Im}(f_{\alpha \beta ab})\mathrm{sin}(2b_0m_{\beta \alpha }^2)\mathrm{exp}(2\sigma _b^2(m_{\beta \alpha }^2)^2)],`$ (22) where $`f_{\alpha \beta ab}=K_{a\alpha }K_{a\beta }^{}K_{b\alpha }^{}K_{b\beta }`$ and $`m_{\beta \alpha }^2=m_\beta ^2m_\alpha ^2`$. Notice that when $`\mathrm{Im}(f_{\alpha \beta ab})=0`$, $`P_{ab}`$ is independent of the sign of $`m_{\beta \alpha }^2`$. ## IV Learning about $`M_\nu `$ from experiment Since $`\mathrm{Tr}\left(M_\nu \right)=0`$ provides only two real equations for 12 real parameters, it is clear that it has a relatively small amount of predictivity. In particular it cannot say much about the texture (e.g. possible zeroes) of $`M_\nu `$ which is suppposed to derive from a deeper theory than the standard model. On the other hand, we shall see that our Ansatz is complementary to the results which should emerge from analysis of neutrino oscillation experiments. Together, they should enable us to actually (with some conditions) reconstruct $`M_\nu `$. First, in this section, we shall consider $`M_\nu `$ to be hermitian so that the argument in favor of $`\mathrm{Tr}\left(M_\nu \right)=0`$ presented in section II holds without any further assumptions. Since $`M_\nu `$ is symmetric it must be real. It can be brought to real diagonal form via a real rotation $`R`$ as $`R^TM_\nu R`$. However there is no guarantee that all eigenvalues of $`M_\nu `$ will be positive. We can make them all positive by rephasing the diagonal fields with negative eigenvalues by a factor $`i`$ (see Eqs. (3.1) and (3.2)). This means that the general diagonalizing matrix $`U`$ in Eq. (3.2) now takes the form $$U=R\omega ,$$ (23) where $`\omega _{ab}=\delta _{ab}\eta _b^{\frac{1}{2}}`$ with $`\eta _b^{\frac{1}{2}}=1`$ for a positive eigenvalue and $`\eta _b^{\frac{1}{2}}=i`$ for a negative eigenvalue. We notice that Eq. (4.1) is of the form Eq. (3.12) for which we already noticed that the factor $`\omega `$ cancels out in the neutrino oscillation formula Eq. (3.14). Furthermore only the square of the mass is relevant in Eq. (3.14). Thus we choose to work in this section with some negative masses and no factor $`\omega `$ in Eq. (4.1). To avoid confusion, we remark that the factor $`\omega `$ in Eq. (4.1) does not introduce any CP violation in the theory since $`M_\nu `$ is real in any event. Now let us suppose that an experimental analysis of all neutrino oscillation experiments is made based on a formula like Eq. (3.14) (or one which treats the experimental uncertainties in a more sophisticated way). Furthermore assume that the CP violating phase $`\varphi _{13}`$ in Eq. (3.12) is negligible. Then we should know the magnitudes of the squared neutrino mass differences $$\left(m_2\right)^2\left(m_1\right)^2=A,\left(m_3\right)^2\left(m_2\right)^2=B,$$ (24) where $`A`$ and $`B`$ can be either positive or negative. Then, assuming the leptonic theory to be CP conserving, our Ansatz would imply $$0=\mathrm{Tr}\left(M_\nu \right)=\mathrm{Tr}\left(R\widehat{M}R^T\right)=m_1+m_2+m_3,$$ (25) where Eq. (3.2) was used. Eqs. (4.2) and (4.3) comprise three equations for the three neutrino masses $`m_1`$, $`m_2`$ and $`m_3`$. We can solve to get: $`\left(m_1\right)^2`$ $`=`$ $`{\displaystyle \frac{2}{3}}\left[\sqrt{A^2+B^2+AB}\left(A+{\displaystyle \frac{B}{2}}\right)\right],`$ (26) $`m_2`$ $`=`$ $`{\displaystyle \frac{B\left(m_1\right)^2}{2m_1}},`$ (27) $`m_3`$ $`=`$ $`m_1m_2.`$ (28) This leads to a limited number of solutions, depending on sign choices. If we make the further assumption that the charged lepton mixing matrix $`\mathrm{\Omega }`$ in Eq. (3.4) is approximately the unit matrix (this is expected to be a reasonable but not perfect approximation) we can identify $`RK`$ which would be obtained from experiment. Then, using the masses found in Eq. (4.4), we could reconstruct $`M_\nu `$ as $$M_\nu R\widehat{M}R^T.$$ (29) To proceed, we need only insert the experimental results for $`A`$, $`B`$ and $`K`$ in (4.4) and (4.5). Of course, it is presumably the task of the next decade to solidify the experimental determination of these quantities. We can, at the moment, only give a preliminary discussion. For this purpose we will use the results of a recent preliminary analysis of all neutrino experiments by Ohlsson and Snellman . These authors found, by a least square analysis, a best fit for (our notation) $`\left\|A\right\|`$, $`\left\|B\right\|`$ and the leptonic mixing matrix $`K`$. They used the formula Eq. (3.14) with a suitable choice of $`b_0`$ and $`\sigma _b^2`$ for each experiment. Furthermore, they made the simplifying assumption that $`K`$ is real. Finally they only searched for a fit in the range $`10^4\mathrm{eV}^2\left\|A\right\|10^3\mathrm{eV}^2`$, $`0.2\mathrm{eV}^2\left\|B\right\|2\mathrm{e}\mathrm{V}^2`$. This range corresponds to mass difference choices for which the MSW effect for solar and atmospheric neutrinos is not expected to be important and so greatly simplifies the analysis. Thus there is no guarantee that the solution of is unique. Altogether they fit sixteen different solar neutrino, atmospheric neutrino, accelerator and reactor experiments, including LSND. The best fit is: $$\left\|A\right\|=2.87\times 10^4\mathrm{eV}^2,\left\|B\right\|=1.11\mathrm{eV}^2,$$ (30) for the squared mass differences and $$K_{exp}=\left[\begin{array}{ccc}0.7052& 0.7052& 0.0732\\ 0.6441& 0.5940& 0.4820\\ 0.2964& 0.3871& 0.8731\end{array}\right]$$ (31) for the lepton mixing matrix $`K`$. As discussed above we will identify $`KR`$ here, keeping $`K`$ real but allowing for negative masses. The best fit matrix K was obtained to be similar but not identical to the “bimaximal mixing” matrix . With the best fit squared mass differences in Eq. (4.6), our model predicts, from the first of Eq. (4.4), eight different possibilities. These correspond to four different sign configurations for $`A`$ and $`B`$ times the two possible signs for $`m_1`$. However, only two of these eight are essentially different; these are $`\mathrm{type}`$ $`\mathrm{I}:m_1=0.6082\mathrm{eV},m_2=0.6084\mathrm{eV},m_3=1.2166\mathrm{eV}`$ (32) $`\mathrm{type}`$ $`\mathrm{II}:m_1=1.053701\mathrm{eV},m_2=1.053565\mathrm{eV},m_3=0.000136\mathrm{eV}`$ (33) The other solutions correspond to interchanging whichever of $`\left\|m_1\right\|`$ and $`\left\|m_2\right\|`$ is greater (which has only a negligible effect since they are almost degenerate) or reversing the signs of all masses. Physically it is clear what is happening: the smallness of $`\left\|A\right\|`$ compared to $`\left\|B\right\|`$ in Eq. (4.2) forces $`\left\|m_1\right\|\left\|m_2\right\|`$. Then we have either $`m_1m_2`$ with, using the constraint Eq. (4.3), $`m_32m_1`$ or $`m_1m_2`$ with $`m_3`$ very small. Since we have assumed the neutrinos to be of Majorana type for our plausibility argument in section II, their interactions will violate lepton number. Then they should mediate neutrinoless double beta decay $`\left(\beta \beta _{0\nu }\right)`$ . Such a process has not yet been observed and an upper bound has been set for the relevant quantity $$m_\nu \left\|\underset{\alpha =1}{\overset{3}{}}\left(K_{1\alpha }\right)^2m_\alpha \right\|.$$ (34) The best upper bound at present is $`m_\nu 0.20.6`$ eV, reflecting some uncertainty in the estimation of the needed nuclear matrix elements. Substituting the best fit for the matrix K from Eq. (4.7) together with our results in Eqs. (4.8) and (4.9) into Eq. (4.10) yields predictions for the two cases: $`\mathrm{typeI}:m_\nu `$ $`=`$ $`0.60\mathrm{eV},`$ (35) $`\mathrm{typeII}:m_\nu `$ $`=`$ $`6.7\times 10^5\mathrm{eV}.`$ (36) Both solutions seem to be acceptable, the type I case marginally but the type II case definitely. Note that the small value for $`m_\nu `$ in the type II case is due to the best fit prediction $`K_{11}=K_{12}`$ and also to the fact that $`m_2`$ is negative. The same value would clearly result if we made $`m_2`$ positive and set $`K_{12}=0.7052i`$ as discussed around Eq. (4.1) above. Finally, let us reconstruct the underlying neutrino mass matrices for each of the two cases. We use Eq. (4.5) based on the assumption that $`M_\nu `$ is real and also our ansatz to find (in units of eV): $`\mathrm{typeI}:M_\nu `$ $`=`$ $`M_\nu ^T\left[\begin{array}{ccc}0.5985& 0.0643& 0.1167\\ 0.0643& 0.1843& 0.7680\\ 0.1167& 0.7680& 0.7828\end{array}\right],`$ (40) $`\mathrm{typeII}:M_\nu `$ $`=`$ $`M_\nu ^T\left[\begin{array}{ccc}6.7\times 10^5& 0.9199& 0.5078\\ 0.9199& 0.0654& 0.0410\\ 0.5078& 0.0410& 0.0654\end{array}\right].`$ (44) The type I matrix does not have an excellent candidate for a “texture” zero. However the small value of $`\left(M_\nu \right)_{11}`$ in the type II case is certainly suggestive. These matrices lead to neutrino masses and a mixing matrix which give a best fit to all present data. It will be interesting to see if either of them hold up in the future. Incidentally, on comparing Eqs. (4.13) and (4.14) it is amusing to observe the large difference in two mass matrices “generated” in the same way except with respect to how $`\mathrm{Tr}\left(M_\nu \right)=0`$ is satisfied. ## V Case of complex $`M_\nu `$ It seems interesting to also investigate the Ansatz $`\mathrm{Tr}(M_\nu )=0`$ when $`M_\nu `$ is no longer restricted to be real. This also raises the problem of constructing the unitary diagonalizing matrix $`U`$ in Eq.(3.2), in terms of the experimentally measured lepton mixing matrix $`K_{exp}`$. For simplicity, as before, we will make the approximation that the charged lepton diagonalizing matrix $`\mathrm{\Omega }`$ is the unit matrix. Apart from an overall (conventional) phase we may write $$U=\omega _0(\sigma )K_{exp}\omega _0^1(\tau ),$$ (45) where the 2-parameter quantity $`\omega _0`$ was defined in Eq. (3.8). Since $`K_{exp}`$ has four parameters $`U`$ in Eq. (5.1) is described by eight parameters. As mentioned before, the two parameters in $`\omega _0^1(\tau )`$ are not measurable in neutrino oscillation experiments but show up when one considers $`(\beta \beta )_{o\nu }`$. The two parameters in $`\omega _0(\sigma )`$ may be eliminated, for experimental purposes, by rephasing the charged leptons. However, for the theoretical purpose of reconstructing the underlying neutrino mass matrix $`M_\nu `$, their existence cannot be ruled out. (They also do not contribute to $`(\beta \beta )_{0\nu }`$.) For the purpose of relating the Ansatz on $`M_\nu `$ to the physical neutrino masses in $`\widehat{M}`$, we note $$\mathrm{Tr}(M_\nu )=\mathrm{Tr}(K_{exp}^1\omega _0^1(2\sigma )K_{exp}^{}\omega _0(2\tau )\widehat{M}).$$ (46) For further simplicity of the analysis we adopt the special case $`\omega _0(2\sigma )=1`$ and also identify $`K_{exp}`$ with the real best fit in Eq. (4.7); our Ansatz now reads $$0=\mathrm{Tr}(M_\nu )=\mathrm{Tr}\left(\omega _0(2\tau )\widehat{M}\right).$$ (47) With the redefinitions $`\beta _1=4\tau _1+2\tau _2`$ and $`\beta _2=2\tau _1+4\tau _2`$, Eq. (5.3) becomes $$e^{i\beta _1}m_1+e^{i\beta _2}m_2+m_3=0$$ (48) This may be conveniently visualized as the vector triangle shown in Fig. 1. Combining Eqs. (5.4) and (4.2) gives four real equations for the five unknown quantities ($`m_1,m_2,m_3,\beta _1,\beta _2`$). Thus we have (for each set of ($`A,B`$) sign choices) a one parameter family of solutions. It is convenient to choose this parameter to be $`m_3`$. Then $`m_1`$ and $`m_2`$ may be found from the equations (4.2), provided that solutions exist. In this way all three sides of the triangle in Fig.1 are determined. The angles may finally be found as $`\mathrm{cos}\beta _2`$ $`=`$ $`{\displaystyle \frac{m_1^2m_2^2m_3^2}{2m_2m_3}}`$ (49) $`\mathrm{sin}\beta _1`$ $`=`$ $`{\displaystyle \frac{m_2}{m_1}}\mathrm{sin}\beta _2.`$ (50) We also need to investigate the constraint arising from the non-observation of $`(\beta \beta )_{0\nu }`$. Eq. (4.10) now becomes, with Eq. (5.1) as the mixing matrix $$m_\nu =\left\|m_1(K_{exp11})^2e^{i\beta _1}+m_2(K_{exp12})^2e^{i\beta _2}+m_3(K_{exp13})^2\right\|.$$ (51) Using the Ansatz constraint Eq. (5.4), Eq. (5.6) may be rewritten as $$m_\nu =\left\|\left[(K_{exp12})^2(K_{exp11})^2\right]m_2e^{i\beta _2}+\left[(K_{exp13})^2(K_{exp11})^2\right]m_3\right\|.$$ (52) This form is very convenient when identifying $`K_{exp}`$ with the best fit solution in Eq. (4.7). In the present context such an identification corresponds to CP violation for the $`(\beta \beta )_{0\nu }`$ process but not for usual neutrino oscillations. Since the (11) and (12) matrix elements are equal in Eq. (4.7) we find the simple result $$m_\nu =0.49m_3.$$ (53) Thus if the upper bound on $`m_\nu `$ is conservatively identified as in the $`0.20.6`$ eV range, we should have in this case $$m_30.411.22\mathrm{eV}.$$ (54) In the present complex case there is a continuum of possible solutions labelled by those values of $`m_3`$ satisfying Eq. (5.9), rather than just the two possibilities found in Eq. (4.8) and Eq. (4.9). Actually, the continuum separates roughly into two classes similar to either Eq. (4.8) or Eq. (4.9). In the generalized type I class, $`m_3`$ is of the order $`B=1.11`$ eV while $`m_1`$ and $`m_2`$ are related to nearly oppositely directed vectors in Fig. 1 and are also of the order $`B`$. In the generalized type II class, $`B`$ is negative; $`m_1`$ and $`m_2`$ correspond to vectors of order $`\|B\|`$ which are oppositely directed to each other, while $`m_3`$ ranges from very small to order $`\left\|B\right\|`$. Given the bound Eq. (5.9) from the non-observation of $`(\beta \beta )_{o\nu }`$, there are important limitations on the allowed $`m_3`$ values for type I solutions. In this case $`B`$ is positive so the equation $$(m_2)^2=(m_3)^2B=(m_3)^2(1.11\mathrm{eV})^2,$$ (55) will only allow solutions for $`m_3>1.11`$ eV. This range is barely compatible with Eq. (5.9). Thus the type II case where $`B<0`$ and $`m_1m_2>m_3`$ seems most probable. As an example of a solution for complex $`M_\nu `$, consider choosing $`m_3=0.2`$ eV. Then a solution is obtained with (compare with Fig. 1) $$m_1m_21.128\mathrm{eV},m_3=0.2\mathrm{eV},\beta _195.1^o,\beta _2264.9^o.$$ (56) The matrix $`U`$, which diagonalizes $`M_\nu `$ is obtained from Eq. (5.1), with the approximation $`\omega _0\left(\sigma \right)=1`$, and with now: $$\omega _0^1(\tau )=\mathrm{diag}(0.976+0.216i,0.3010.953i,0.500+0.866i)$$ (57) This factor introduces CP violation in the $`(\beta \beta )_{0\nu }`$ process but not in ordinary neutrino oscillations. Finally the underlying neutrino mass matrix, $`U\widehat{M}U^T`$ is “reconstructed” as (in units of eV): $$M_\nu =M_\nu ^T=\left[\begin{array}{ccc}0.0490.085i& 0.8550.481i& 0.459+0.287i\\ & 0.076+0.009i& 0.025+0.131i\\ & & 0.125+0.077i\end{array}\right].$$ (58) This is structurally similar to the real type II solution displayed in Eq. (4.14), although the suppression of the (11) element is not so pronounced. Notice that $`m_3`$ is considerably smaller than the almost degenerate pair $`m_1`$ and $`m_2`$. Furthermore $`m_1`$ and $`m_2`$ are large enough to possibly play some role in astrophysics. ## VI Summary and Discussion We investigated the Ansatz $`\mathrm{Tr}\left(M_\nu \right)=0`$ for the underlying (pre-diagonal) three generation neutrino mass matrix. It was motivated by noting that in an SO(10) grand unified model where $`M_\nu `$ was taken to be real (CP conserving), it corresponds to the well known unification of b quark and $`\tau `$ lepton masses. While not very predictive by itself it yields information complementary to what would be gotten from a complete three flavor analysis of all lepton number conserving neutrino oscillation experiments. Specifically from the specification of the magnitudes of two neutrino squared mass differences and also of the leptonic mixing matrix we can, with some assumptions, find the neutrino masses themselves and “reconstruct” $`M_\nu `$. This determination can be sharpened by consideration of the constraints imposed by non-observation of neutrinoless double beta decay. For the purpose of testing our Ansatz we employed the results of a reasonable best fit to all present neutrino experiments (including LSND) by Ohlsson and Snellmann . This fit will inevitably be improved in the next few years as new experiments are completed. They were able to fit the data without assuming any CP violation. This agrees with assuming $`M_\nu `$ to be real. We found two essentially different solutions in that case. The first features two neutrinos having approximately equal mass 0.608 eV and a third neutrino of mass 1.217 eV. This solution is on the borderline of being ruled out by non-observation of $`\left(\beta \beta _{0\nu }\right)`$. The second solution has two neutrinos with approximately degenerate mass 1.054 eV and a third neutrino with a mass $`1.36\times 10^4`$ eV. This solution is very safe from being ruled out by $`\left(\beta \beta _{0\nu }\right)`$ experiments. It also features a reconstructed $`M_\nu `$ which has an extremely small (11) element. Note that, for both solutions, even though $`M_\nu `$ is real there are some (CP conserving) factors of $`i`$ in the mixing matrix when all masses are taken to be positive. Alternatively one may have no $`i`$’s in the mixing matrix while allowing some masses to be negative. The latter form is useful for seeing intuitively how $`\mathrm{Tr}(M_\nu )`$=0 is possible. The case of matching the above best fit data to a complex $`M_\nu `$ was also considered. In this situation there are CP violating phases in the lepton mixing matrix which affect the $`(\beta \beta )_{0\nu }`$ process but do not affect ordinary total lepton number conserving neutrino oscillations. Such phases could also be measurable in principle with the observation of a decay like $`\tau ^{}\pi ^{}\pi ^{}\mu ^+`$. The case of complex $`M_\nu `$ allows a larger number of solutions. With a simplifying assumption there is a one parameter family of allowed neutrino mass sets. Roughly, these fall into one of the two types already encountered for real $`M_\nu `$. A question of some interest is whether the neutrinos are massive enough to play a role in cosmology. The relevant criterion for this to occur is usually stated as $`_am_a>6`$ eV. For the type II solutions with complex $`M_\nu `$ we have found the largest mass sum to be about 4.5 eV corresponding to $`m_1m_2=1.62`$ eV and $`m_31.22`$ eV. However this is on the very border of acceptability for non observation of $`(\beta \beta )_{0\nu }`$. Future best fits to the neutrino oscillation data can easily be accommodated in the present framework. Of course, the predictions for neutrino masses and mixings will depend on this input. ###### Acknowledgements. We would like to thank H. Benaoum for useful comments on the manuscript. This work has been supported in part by the US DOE under contract DE-FG-02-85ER 40231.
warning/0004/nucl-th0004012.html	ar5iv	text	# References Pion - Nucleon Bremsstrahlung beyond the Soft-Photon approximation A. Mariano<sup>a,b</sup> <sup>a</sup> Departamento de Física, Centro de Investigación y de Estudios Avanzados del IPN A.P. 14$``$740 México 07000 D.F. <sup>b</sup> Departamento de Física, Facultad de Ciencias Exáctas, Universidad Nacional de La Plata cc.67, 1900 La Plata, Argentina Abstract A dynamical model based on effective Lagrangians is proposed to describe the bremsstrahlung reaction $`\pi N\pi N\gamma `$ at low energies. The $`\mathrm{\Delta }(1232)`$ degrees of freedom are incorporated in a way consistent with both, electromagnetic gauge invariance and invariance under contact transformations. The model also includes the initial and final state rescattering of hadrons via a T-matrix with off the momentum-shell effects. The double differential distribution of photons is computed for three different T-matrix models and the results are compared with the soft photon approximation, and with experimental data. The aim of this analysis is to test the off-shell behaviour of the different T-matrices under consideration. Finally an alternative simpler dynamical model that incorporates the unstable character of the isobar-$`\mathrm{\Delta }(1232)`$ through a complex mass, is presented. As we will see it is suitable for the study of the magnetic moment of the resonance. PACS numbers: 25.80.Ek, 13.60.-n, 13.75.-r 1. INTRODUCTION In order to extract resonant parameters of the nucleon resonances(N) from the $`\gamma N\pi N`$ reaction, it is important to evaluate the background contribution to isolate the resonant peak. An important contribution to the background of this photoproduction reaction is provided by the final state rescattering (FSI) of the $`\pi N`$ system . Consequently, we require the knowledge of the T-matrix that describes this rescattering process, in the off the momentum-shell regime(off-shell),i.e, we need $`T(\stackrel{}{q}{}_{}{}^{},\stackrel{}{q};z(\stackrel{}{q}))`$ with $`\|\stackrel{}{q}\|\|\stackrel{}{q}{}_{}{}^{}\|`$ where $`z`$ is the total energy of the $`\pi N`$ system as a function of the relative momentum $`\|\stackrel{}{q}\|`$ of the initial state. This particular rescattering amplitude is more properly called half-off-shell T-matrix. In all cases the so called ‘realistic’ interactions are fitted to reproduce the phase shifts in elastic $`\pi N`$ scattering, which only depends on the on-shell ($`\|\stackrel{}{q}\|=\|\stackrel{}{q}{}_{}{}^{}\|`$) values of the relative momenta. Thus, elastic scattering is not useful to constrain the off-shell behavior of the T-matrix because interactions working similarly in elastic scattering, may have different behaviour in the off-shell regime. Another reaction where the $`\pi N`$ off-shell T-matrix is required is $`\pi N\pi N\gamma `$ bremsstrahlung. This process has been studied within the Soft Photon Approximation (SPA). Within this approximation the full amplitude, expressed as a power expansion in the photon energy, depends only on the electromagnetic static multipoles of $`\pi `$ and $`N`$ and on the T-matrix (and its derivatives) of the corresponding non-radiative process . The SPA reproduces very well the experimental data on radiative $`\pi N`$ scattering, in spite of some objections that have been raised recently regarding the departures of the formulation of ref. from the original Low’s prescription. Because of the power expansion in the photon energy, the total SPA amplitude depends on the derivatives of the T-matrix, and thus on off-shell effects. Since the non-leading terms (of the power expansion) are fixed by imposing electromagnetic gauge invariance, the off-shell effects of the T-matrix cancel up mutually. Therefore, the information on the off-shell behaviour of the T-matrix can be tested by adding the contributions to the radiative $`\pi N`$ scattering within the framework of an specific dynamical model. The purpose of the present paper is to check the off-shell behaviour of three different T-matrices for $`\pi N`$ rescattering in the reaction $`\pi N\pi N\gamma `$. We use a dynamical model to describe the $`\pi N\pi N\gamma `$ reaction. The gauge invariant electromagnetic current is constructed explicitly, with vertices and propagators derived from the relevant hadronic and electromagnetic Lagrangians. We include also two-body meson exchange currents, and the full energy-momentum dependence of the T-matrix for the elastic $`\pi N`$ scattering which exhibits its off-shell behaviour. Finally we implement this model with different T-matrices in order to compare their different off-shell dependence. On the other hand since different T-matrices depend on many parameters (bare masses and coupling constants of the hadrons, cut-off form factor parameters, etc.), results difficult to use a dynamical model based on a T-matrix input for analyzing resonance unknown properties. For this reason we also present a simpler formalism that assumes the isobar-$`\mathrm{\Delta }`$ as the main degree of freedom, and gives it an unstable character through a complex mass. This approach will be used to study the anomalous magnetic moment of the resonance. This paper is organized as follows. In section 2 we will construct the gauge-invariant amplitude of radiative $`\pi N`$ scattering.The Lagrangians used to construct the gauge invariant current of our process, are provided also in this section. The second simpler dynamical model will be described in section 3 . Finally the results and conclusions are given in section 4. 2. GAUGE INVARIANT BREMSSTRAHLUNG AMPLITUDE AND DYNAMICAL MODEL In the pion-nucleon bremsstrahlung process we deal with a problem of scattering by two potentials : the strong pion-nucleon and the electromagnetic interactions. The cross section for $`\pi N\pi N\gamma `$ process reads $`d\sigma `$ $`=`$ $`{\displaystyle \frac{d\stackrel{}{k}}{\omega _\gamma }\frac{d\stackrel{}{q}_f}{\omega (\stackrel{}{q}_f)}\frac{d\stackrel{}{p}_f}{E(\stackrel{}{p}_f)}(2\pi )^4\delta ^4(p_i+q_ip_fq_fk)}`$ (1) $`\times {\displaystyle \frac{1}{2}}{\displaystyle \underset{ϵ_\lambda ,ms_f,ms_i}{}}\left\|{\displaystyle \frac{m_N^2}{2\sqrt{2}}}M_{\pi N\gamma ,\pi N}(ϵ_\lambda ,k;q_f,p_f,ms_f;q_i,p_i,ms_i)\right\|^2,`$ where $`q=(\omega ,\stackrel{}{q})`$, $`p=(E,\stackrel{}{p})`$ and $`k=(\omega _\gamma ,\stackrel{}{k})`$ denote pion, nucleon and photon four-momenta, respectively; $`ms`$ is the nucleon’s spin projection and $`ϵ_\lambda `$ indicates the polarization four-vector of the photon. The subindexes $`i,f`$ refer to initial and final state quantities. The Lorentz invariant amplitude<sup>1</sup><sup>1</sup>1Throughout this paper, $`M`$ will denote the amplitude generated by the operator $`\widehat{M}`$,ie., $`M=\overline{u}\|\widehat{M}\|u`$. $`M_{\pi N\gamma ,\pi N}`$ explicitly reads $`M_{\pi N\gamma ,\pi N}=\overline{u}(\stackrel{}{p}_f,ms_f)\|\widehat{M}_{\pi N\gamma ,\pi N}(ϵ_\lambda ,k;q_f,p_f;q_i,p_i)\|u(\stackrel{}{p}_i,ms_i),`$ (2) where $`u(\stackrel{}{p},ms)`$ denote nucleon Dirac spinors, and the amplitude operator $`\widehat{M}_{\pi N\gamma ,\pi N}`$ is obtained from the coupled channel Bethe-Salpeter equation for the $`\pi N\gamma `$ system as follows (we consider electromagnetic interactions at the lowest order) $`\widehat{M}_{\pi N\gamma ,\pi N}`$ $`=`$ $`\widehat{V}_{\pi N\gamma ,\pi N}`$ (3) $`+i{\displaystyle \frac{dq^4}{(2\pi )^4}\left[\widehat{V}_{\pi N\gamma ,\pi N}(q)\widehat{G}_{\pi N}(q)\widehat{M}_{\pi N,\pi N}(q)+\widehat{M}_{\pi N,\pi N}(q)\widehat{G}_{\pi N}(q)\widehat{V}_{\pi N\gamma ,\pi N}(q)\right]}`$ $`+i^2{\displaystyle \frac{dq^4}{(2\pi )^4}\frac{dq^4}{(2\pi )^4}\left[\widehat{M}_{\pi N,\pi N}(q^{})\widehat{G}_{\pi N}(q^{})\widehat{V}_{\pi N\gamma ,\pi N}(q^{},q)\widehat{G}_{\pi N}(q)\widehat{M}_{\pi N,\pi N}(q)\right]}.`$ In terms of the above operator amplitude the T-matrix, defined as $$\widehat{T}(q_f,p_f;q_i,p_i)=\frac{1}{(2\pi )^3}\widehat{M}_{\pi N,\pi N}(q_f,p_f;q_i,p_i),$$ (4) satisfies the integral equation $`\widehat{T}`$ $`=`$ $`\widehat{U}+i{\displaystyle \frac{dq^4}{(2\pi )^4}\widehat{U}(q)\widehat{G}(q)\widehat{T}(q)},`$ (5) $`\widehat{U}`$ $`=`$ $`{\displaystyle \frac{1}{(2\pi )^3}}\widehat{V}_{\pi N,\pi N},`$ $`\widehat{G}`$ $`=`$ $`(2\pi )^3\widehat{G}_{\pi N}.`$ In the previous equations $`\widehat{V}_{ij}`$ denote $`\widehat{M}`$-matrix elements corresponding to the irreducible Feynman diagrams for each process, while $`\widehat{G}_i`$ is the product of Feynman propagators of intermediate particles. Following the Thompson’s prescription, we can set the above integrals in a three-dimensional form as follows (we set in the center of mass frame of the $`\pi N`$ system) $$\widehat{T}(\stackrel{}{q}{}_{}{}^{},\stackrel{}{q},z)=\widehat{U}(\stackrel{}{q}{}_{}{}^{},\stackrel{}{q})+d^3\stackrel{}{q}\mathrm{"}\widehat{U}(\stackrel{}{q}{}_{}{}^{},\stackrel{}{q}\mathrm{"})\widehat{G}_{TH}(z,\stackrel{}{q}\mathrm{"})\widehat{T}(\stackrel{}{q}\mathrm{"},\stackrel{}{q},z),$$ (6) with, $$\widehat{G}_{TH}(z,\stackrel{}{q}\mathrm{"})=\frac{m_N}{2\omega (\stackrel{}{q}\mathrm{"})E(\stackrel{}{q}\mathrm{"})}\frac{\underset{ms\mathrm{"}}{}\|u(\stackrel{}{q}\mathrm{"},ms\mathrm{"})\overline{u}(\stackrel{}{q}\mathrm{"},ms\mathrm{"})\|}{zz\mathrm{"}+i\eta },$$ (7) where $`z\mathrm{"}=E(\stackrel{}{q}\mathrm{"})+\omega (\stackrel{}{q}\mathrm{"})`$. In the above expressions $`\widehat{G}_{TH}`$ denotes the Thompson propagator replacing the full $`\widehat{G}_{\pi N}`$ Feynman propagator which, as a consequence of the three-dimensional reduction, eliminates the propagation of antiparticles and put intermediate particles on their mass-shell. The kernel function $`\widehat{U}(\stackrel{}{q}{}_{}{}^{},\stackrel{}{q})`$ contains all the $`\pi N`$-interaction irreducible diagrams to be iterated in the T-matrix calculation, but usually only second-order contributions are kept. The electromagnetic current $`\widehat{V}_{\pi N\gamma ,\pi N}`$ can be broken into two pieces $$\widehat{V}_{\pi N\gamma ,\pi N}\widehat{V}_{\pi N\gamma ,\pi N}^{(1)}+\widehat{V}_{\pi N\gamma ,\pi N}^{(2)},$$ (8) where the upper indices denote one- and two-body contributions, respectively which are obtained by coupling the photon to all the internal lines of $`\widehat{U}`$. As is known, the operator $`\widehat{V}_{\pi N\gamma ,\pi N}^{(2)}`$ must be added to the electromagnetic current in order to satisfy the electromagnetic gauge invariance of the total amplitude, while the one-body amplitude $`V_{\pi N\gamma ,\pi N}^{(1)}`$ vanishes for free hadrons. Both contributions to the total amplitude are illustrated in Fig. 1. Fig.1 One- and two-body contributions to the bremsstrahlung current amplitude. We follow a procedure that put the bremsstrahlung amplitude manifestly gauge invariant. By replacing eqs. (4) and (5) into eq.(3), only in the one-body component of the amplitude $`M_{\pi N\gamma ,\pi N}^{(1)}M_{\pi N\gamma ,\pi N}(\widehat{V}_{\pi N\gamma ,\pi N}^{(1)})`$, we can isolate the lowest order nonzero contribution of the one-body current. After the three dimensional reduction, the total amplitude can be rewritten as $$M_{\pi N\gamma ,\pi N}\left[\stackrel{~}{V}_{\pi N\gamma ,\pi N}+\stackrel{~}{M}_{\pi N\gamma ,\pi N}^{pre}+\stackrel{~}{M}_{\pi N\gamma ,\pi N}^{post}+\stackrel{~}{M}_{\pi N\gamma ,\pi N}^{double}\right],$$ (9) with $`\stackrel{~}{V}_{\pi N\gamma ,\pi N}`$ $`=`$ $`\overline{u}(\stackrel{}{q}{}_{}{}^{}\stackrel{}{k}/2,ms_f)\|\widehat{\stackrel{~}{V}}_{\pi N\gamma ,\pi N}(ϵ_\lambda ,\stackrel{}{k},\stackrel{}{q}{}_{}{}^{},\stackrel{}{q})\|u(\stackrel{}{q},ms_i),`$ $`\stackrel{~}{M}_{\pi N\gamma ,\pi N}^{pre}`$ $`=`$ $`{\displaystyle }dq^{\prime \prime 3}\overline{u}(\stackrel{}{q}{}_{}{}^{},ms_f)\|\widehat{T}^{()}(\stackrel{}{q}{}_{}{}^{},\stackrel{}{q}^{\prime \prime },z^{})\widehat{G}_{TH}(z^{},\stackrel{}{q}^{\prime \prime })\widehat{\stackrel{~}{V}}_{\pi N\gamma ,\pi N}(ϵ_\lambda ,\stackrel{}{k},\stackrel{}{q}^{\prime \prime },\stackrel{}{q})\|u(\stackrel{}{q}+\stackrel{}{k}/2),ms_i,`$ $`\stackrel{~}{M}_{\pi N\gamma ,\pi N}^{post}`$ $`=`$ $`{\displaystyle }dq^{\prime \prime 3}\overline{u}(\stackrel{}{q}{}_{}{}^{}\stackrel{}{k}/2,ms_f)\|\widehat{\stackrel{~}{V}}_{\pi N\gamma ,\pi N}(ϵ_\lambda ,\stackrel{}{k},\stackrel{}{q}{}_{}{}^{},\stackrel{}{q}^{\prime \prime })\widehat{G}_{TH}(z,\stackrel{}{q}^{\prime \prime })\widehat{T}(\stackrel{}{q}^{\prime \prime },\stackrel{}{q},z)\|u(\stackrel{}{q},ms_i),`$ $`\stackrel{~}{M}_{\pi N\gamma ,\pi N}^{double}`$ $`=`$ $`{\displaystyle }dq^{\prime \prime 3}{\displaystyle }dq^{\prime \prime \prime 3}\overline{u}(\stackrel{}{q}{}_{}{}^{}\stackrel{}{k}/2,ms_f)\|\widehat{T}^{()}(\stackrel{}{q}{}_{}{}^{}\stackrel{}{k}/2,\stackrel{}{q}^{\prime \prime },z^{})`$ (10) $`\widehat{G}_{TH}(z^{},\stackrel{}{q}^{\prime \prime })\widehat{\stackrel{~}{V}}_{\pi N\gamma ,\pi N}(ϵ_\lambda ,\stackrel{}{k},\stackrel{}{q}^{\prime \prime },\stackrel{}{q}^{\prime \prime \prime })\widehat{G}_{TH}(z,\stackrel{}{q}^{\prime \prime \prime })\widehat{T}(\stackrel{}{q}^{\prime \prime \prime },\stackrel{}{q},z)\|u(\stackrel{}{q},ms_i)`$ where the current $`\widehat{\stackrel{~}{V}}_{\pi N\gamma ,\pi N}=i\widehat{V}_{\pi N\gamma ,\pi N}^{(1)}\widehat{G}\widehat{U}+i\widehat{U}^{}\widehat{G}\widehat{V}_{\pi N\gamma ,\pi N}^{(1)}+\widehat{V}_{\pi N\gamma ,\pi N}^{(2)},`$ (11) generates a gauge invariant Born amplitude $`\stackrel{~}{V}_{\pi N\gamma ,\pi N}`$, that involve all the possible ways of attaching a photon to the $`\pi N`$ scattering amplitude $`U`$, and contains the full propagator $`\widehat{G}\widehat{G}_{\pi N}`$. The operator $`\widehat{T}^{()}(z^{})`$, where $`z^{}=z+\omega _\gamma `$, obeys eq.(6) if we change $`\eta \eta `$ in eq. (7). The superscript pre (post) in Eq. (9) indicate that the photon is emitted before (after) the action of the T-matrix, while the superscript double refers to a double-scattering term where the photon is emitted from internal lines between two T-matrices. In the above equations the Born, pre and double amplitudes were evaluated in the initial center of mass frame ($`\stackrel{}{q}_f=\stackrel{}{q}{}_{}{}^{}\stackrel{}{k}/2,\stackrel{}{p}_f=\stackrel{}{q}{}_{}{}^{}\stackrel{}{k}/2`$, $`\stackrel{}{q}_i=\stackrel{}{p}_i=\stackrel{}{q}`$), while the other (post) amplitude was evaluated in the corresponding final frame($`\stackrel{}{q}_f=\stackrel{}{p}_f=\stackrel{}{q}^{}`$, $`\stackrel{}{q}_i=\stackrel{}{q}+\stackrel{}{k}/2,\stackrel{}{p}_i=\stackrel{}{q}+\stackrel{}{k}/2`$). The different terms in Eq. (9) are illustrated in Fig. 2a ($`\stackrel{~}{V}_{\pi N\gamma ,\pi N}`$) and in Fig. 2b (remaining terms). Fig.2 (a) Gauge-invariant bremsstrahlung current amplitude. (b) Post-,pre- and double-scattering amplitude contributions (se eq. (9)). Within the SPA, the total bremsstrahlung amplitude can be split into external ($`E`$) and internal ($`I`$) contributions: $`M_{\pi N\gamma ,\pi N}M_{\pi N\gamma ,\pi N}^E+M_{\pi N\gamma ,\pi N}^I,`$ (12) where we can identify $`M_{\pi N\gamma ,\pi N}^E`$ $``$ $`\stackrel{~}{M}_{\pi N\gamma ,\pi N}^{pre}(\widehat{\stackrel{~}{V}}_{\pi N\gamma ,\pi N}\widehat{V}_{\pi N\gamma ,\pi N}^{(1)})+\stackrel{~}{M}_{\pi N\gamma ,\pi N}^{post}(\widehat{\stackrel{~}{V}}_{\pi N\gamma ,\pi N}\widehat{V}_{\pi N\gamma ,\pi N}^{(1)}),`$ and the internal contribution $`M_{\pi N\gamma ,\pi N}^I`$ can be obtained “by enforcing” the gauge invariance condition $$M_{\pi N\gamma ,\pi N}(ϵ_\lambda ^\mu =k^\mu )=0.$$ (13) The SPA the $`M_{\pi N\gamma ,\pi N}`$ amplitude depends only on the elastic T-matrix, because derivative terms of $`\widehat{T}`$ cancels in the addition of internal and external contributions. Let us emphasize that any dependence on the structure of internal contributions (in particular, the dependence of off-shell effects) are of higher order in $`\omega _\gamma `$ and must be included explicitly in the amplitude in a gauge invariant way. The bremsstrahlung amplitude will be computed along the lines developed in eqs. (9-11), using a potential $`\widehat{U}`$ obtained from effective Lagrangians, and three specific models for the T-matrix to describe the $`\pi N`$ rescattering. The three models advocated for the T-matrix will be called OBQA, SEP and NEW, respectively. The OBQA version for the T-matrix interaction, is based on a model that includes $`\pi `$ and $`\rho `$ mesons exchange through a correlated $`2\pi `$ exchange potential. The SEP model for $`\pi N`$ rescattering is generated by a phenomenological separable potential. Finally, the NEW model is obtained from exchange of $`\pi `$ and (sharp) $`\rho `$ mesons. The operator $`\widehat{U}`$ is constructed from a Lagrangian density that includes the nucleon($`N`$), the $`\mathrm{\Delta }`$-isobar, and the $`\pi `$, $`\rho `$ and $`\sigma `$ mesons $$\widehat{}_{hadr}=\widehat{}_{\pi NN}+\widehat{}_{N\mathrm{\Delta }\pi }+\widehat{}_{\rho NN}+\widehat{}_{\rho \pi \pi }+\widehat{}_{\sigma NN}+\widehat{}_{\sigma \pi \pi },$$ (14) while the electromagnetic currents are obtained from the hadronic Lagrangian density through minimal coupling of the photon. The Scattering amplitude $`U`$ is depicted in Fig.3, while the amplitude $`\stackrel{~}{V}_{\pi N\gamma ,\pi N}`$ can be obtained by coupling the photon to all diagrams in $`U`$ as shown in fig.4. Fig.3 Born amplitude corresponding to the $`\pi N`$ potential operator $`\widehat{U}`$. The first diagram denote the nucleon-pole, and the $`\mathrm{\Delta }`$-pole corresponds to the third diagram. The fifth and sixth diagrams correspond to $`\rho `$ and $`\sigma `$ mesons exchange. Fig.4 The gauge-invariant amplitude obtained by coupling a photon to the Born terms in fig.3, together with the two-meson exchange currents (fifth, sixth and seventh graphs in each line). (a) Contributions obtained from the N and $`\mathrm{\Delta }`$ direct-pole diagrams in fig.3; (b) Terms generated by the cross-pole diagrams in fig.3, and (c) diagrams obtained from $`\rho `$ and $`\sigma `$ exchange contributions in fig.3. The good convergence properties of the scattering equations given in Eqs. (10) can be obtained by introducing hadronic form factors, which are supposed to describe the composite nature of hadrons. It is a common practice to use different parametrizations of the form factors for different T-matrices. For example, the OBQA model uses monopole and dipole forms with cutoff parameters ranging from $`\mathrm{\Lambda }=12001600MeV`$. In the case of the SEP interaction different form factors are introduced for each partial wave component , while in the NEW form factors usually advocated are of monopolar form with $`\mathrm{\Lambda }=13002300MeV`$ . However, the introduction of form factors replacing point vertices in $`\stackrel{~}{V}_{\pi N\gamma ,\pi N}`$ spoils the gauge invariance of the total amplitude. Fortunately, the gauge invariance of the amplitude can be recovered by using the method of Gross and Riska which, however, does not yield unique electromagnetic couplings to hadrons. Therefore, we follow the more simple prescription of using a common form factor of monopole type $$f(\stackrel{}{q}^{\prime \prime })=\frac{\mathrm{\Lambda }^2}{\mathrm{\Lambda }^2+\stackrel{}{q}^{\prime \prime 2}},$$ (15) where the scale $`\mathrm{\Lambda }`$ can be adjusted at a given incident energy for each T-matrix model. 3. DYNAMICAL MODEL II The dynamical model described in the previous section has the adventage of decribing $`\pi N`$ -FSI with great detail by the inclussion of many components in the T-matrix potential. Nevertheless, these various effective contributions depend on time on certain parameters that are fitted in order to reproduce scattering phase shifts. This makes difficult to analyze (by fitting) unknown parameters of the resonances using a scattering matrix, on which we have additional parameters to be controlled. In this section we will adopt an alternative dynamical model that is simpler, and gives us a tool to study the, until this moment, unknown magnetic moment of the isobar-$`\mathrm{\Delta }(1232)`$ resonance. Taking into account that in the range of energies we will consider ($`T_{lab}300MeV`$) the main contribution comes from the $`\mathrm{\Delta }`$-intermediate terms, since we are in the resonance region, we only consider the $`\mathrm{\Delta }`$-pole contribution in the hadronic potential (third contribution in fig. 3). This reads $`\widehat{U}=\widehat{\mathrm{\Gamma }}(\pi N\mathrm{\Delta })_\nu \widehat{G}^{\nu ,\mu }\widehat{\mathrm{\Gamma }}(\mathrm{\Delta }\pi N)_\mu ,`$ (16) where $`\widehat{G}^{\mu ,\nu }(q)`$ $`=`$ $`{\displaystyle \frac{1}{q^2m_\mathrm{\Delta }^2}}\widehat{O}^{\mu ,\nu }(q)`$ (17) $`\widehat{O}^{\mu ,\nu }(q)`$ $`=`$ $`(\overline{)}q+m_\mathrm{\Delta })\left[g^{\mu \nu }+{\displaystyle \frac{1}{3}}\gamma ^\mu \gamma ^\nu +{\displaystyle \frac{2}{3}}{\displaystyle \frac{q^\mu q^\nu }{m_\mathrm{\Delta }^2}}{\displaystyle \frac{1}{3}}{\displaystyle \frac{q^\mu \gamma ^\nu q^\nu \gamma ^\mu }{m_\mathrm{\Delta }}}\right]`$ (18) $`{\displaystyle \frac{2}{3m_\mathrm{\Delta }^2}}(q^2m_\mathrm{\Delta }^2)\left[{\displaystyle \frac{}{}}(\gamma ^\mu q^\nu \gamma ^\nu q^\mu )+(\overline{)}q+m_\mathrm{\Delta })\gamma ^\mu \gamma ^\nu \right],`$ is the $`\mathrm{\Delta }`$ propagator . Iteration of this contribution to all orders through the T-matrix eq. (6)leads to $`\widehat{\stackrel{~}{U}}=\widehat{\mathrm{\Gamma }}(\pi N\mathrm{\Delta })_\nu \widehat{𝒢}^{\nu ,\mu }\widehat{\mathrm{\Gamma }}(\mathrm{\Delta }\pi N)_\mu `$ (19) where $`\widehat{𝒢}^{\mu ,\nu }(q)={\displaystyle \frac{1}{(q^2m_\mathrm{\Delta }^2)g_\alpha ^\mu \mathrm{\Sigma }(q)_\alpha ^\mu }}\widehat{O}^{\alpha ,\nu }(q)`$ (20) is a modified $`\mathrm{\Delta }`$ propagator with a self-energy $`\mathrm{\Sigma }(q)_\alpha ^\mu =\widehat{O}^{\mu \beta }(q){\displaystyle d^3p\widehat{\mathrm{\Gamma }}(\mathrm{\Delta }\pi N)(p)_\beta \widehat{G}(p,Z)\widehat{\mathrm{\Gamma }}(\pi N\mathrm{\Delta })(p)_\alpha }.`$ In order to simplify we adopt the renormalization recipe $`m_\mathrm{\Delta }^2g_\alpha ^\mu +\mathrm{\Sigma }(q)_\alpha ^\mu (M_\mathrm{\Delta }^2iM_\mathrm{\Delta }\mathrm{\Gamma })g_\alpha ^\mu `$ with $`M_\mathrm{\Delta }`$ and $`\mathrm{\Gamma }`$ constants. This should be equivalent to make the replacement $`m_\mathrm{\Delta }^2M_\mathrm{\Delta }^2iM_\mathrm{\Delta }\mathrm{\Gamma }`$ in the denominator of the $`\widehat{G}^{\mu ,\nu }(q)`$, with which we are considering the unstable character of the $`\mathrm{\Delta }`$. In order to get gauge invariance at the moment of coupling the photon to each line of the potential, we must make this replacement in the full propagator. Finally to calculate the bremsstrahlung amplitude within this model we must replace $`\widehat{\stackrel{~}{U}}`$ by $`\widehat{U}`$ in eq.(11) and keep only the born contribution in eq.(9). The anomalous magnetic moment of the $`\mathrm{\Delta }`$, which is not wellknown will be moved in orther to fit the region of high-energy photons since in the $`\mathrm{\Delta }`$ electromagnetic vertex $`\widehat{\mathrm{\Gamma }}_{\nu \mu \alpha }`$ $`=`$ $`\widehat{e}_\mathrm{\Delta }\left[(\gamma _\alpha g_{\nu \mu }{\displaystyle \frac{1}{3}}\gamma _\alpha \gamma _\nu \gamma _\mu {\displaystyle \frac{1}{3}}\gamma _\nu g_{\mu \alpha }+{\displaystyle \frac{1}{3}}\gamma _\mu g_{\nu \alpha }){\displaystyle \frac{\widehat{\kappa }_\mathrm{\Delta }}{2m_\mathrm{\Delta }}}\sigma _{\alpha \beta }k^\beta g_{\nu \mu }\right],`$ where $`\widehat{\kappa }_\mathrm{\Delta }`$ and $`\widehat{e}_\mathrm{\Delta }`$ are the anomalous magnetic moment<sup>2</sup><sup>2</sup>2We restrict ourselves to the $`\mathrm{\Delta }^{++}`$ contribution, the only one for which we have experimental information on $`\kappa _\mathrm{\Delta }`$ . and charge operators whose action upon the Rarita-Schwinger field give as eigenvalues the corresponding values of these properties, the $`\widehat{\kappa }_\mathrm{\Delta }`$ term goes as $`\omega _\gamma `$. 4. NUMERICAL RESULTS AND CONCLUSIONS The differential cross section $`d\sigma /d\mathrm{\Omega }_\pi d\mathrm{\Omega }_\gamma d\omega _\gamma `$ to be compared with experimental data can be obtained from eq.(1), where the amplitude $`M_{\pi N\gamma ,\pi N}`$ is calculated from eqs.(9-11). The dynamical model approximation (DMA I) advocated in the present paper contains the following steps. The current operator $`\widehat{\stackrel{~}{V}}_{\pi N\gamma ,\pi N}`$ is computed from the effective Lagrangian described in section 2, and the monopole form factor given in eq.(15) to have good convergence of the intermediate momentum integrals. As is known, double scattering terms have a significant contribution in the case of $`protonproton`$ bremsstrahlung, mainly in the end point region of $`\omega _\gamma `$ . In the present work we will neglect $`\stackrel{~}{M}_{\pi N\gamma ,\pi N}^{double}`$ in eq.(10)), because the numerical calculation of the two three-dimensional integrals requires an enormous computation effort. Nevertheless, we keep double scattering-like contributions in post and pre amplitudes coming from current components $`i\widehat{U}^{}\widehat{G}\widehat{V}_{\pi N\gamma ,\pi N}^{(1)}`$, and $`i\widehat{V}_{\pi N\gamma ,\pi N}^{(1)}\widehat{G}\widehat{U}`$, respectively. In order to compare the approaches provided by the DMA I and SPA, we fix $`M_{\pi N\gamma ,\pi N}`$ to coincide quantitatively at low photon energies. In addition we are going to study the sensibility of the bremsstrahlung cross section with the value of the $`\mathrm{\Delta }`$ magnetic moment. For this purpose we also evaluate the cross section through the alternative dynamical model described in section 3 (DMA II). For illustration purposes, we will implemented the DMA I approach with the OBQA, SEP and NEW T-matrices for the specific example of $`\pi ^+p\pi ^+p\gamma `$. The different coupling constants and masses needed to evaluate $`\widehat{V}_{\pi N\gamma ,\pi N}`$ were taken from the model II of ref., from , and from . For direct pole diagrams we use bare masses and coupling constants since they get dynamically dressed in the T-matrix scattering eq.(5). In the OBQA case we replace the $`2\pi `$ correlated exchange potential, by $`\pi `$ and $`\rho `$ sharp mass exchange terms, since they do not lead to sizable differences as shown in ref.. In the SEP case we use the same coupling constants and masses, because the scattering potential is not generated from a dynamical model. We will compare the theoretical predictions to the experimental cross sections measured by Nefkens(EXP), which have been reported for different kinematical configurations. In EXP the pions were detected at fixed angles $`\theta _\pi =50.5^0,\varphi _\pi =180^0`$, for three different energies of incident pions (269, 298 and 324 MeV), and the photons were detected at various $`\theta _\gamma ,\varphi _\gamma `$ angles in the range of energies $`\omega _\gamma =0150`$ MeV. Our results for the cross sections in the DMA I approach, for the different ansatz of the T-matrix, are shown in Fig. 5. The SPA and experimental values are also plotted for comparison. The predictions of the DMA I approach shown in fig. 5 using the three models of the T-matrix, are compared to the results of EXP for photon angles given by $`G_{14}\theta _\gamma =103^0,\varphi _\gamma =180^0`$. Since the parameters entering T-matrices are usually quoted to reproduce the elastic phase shifts, the cutoff parameters $`\mathrm{\Lambda }`$ were fixed in order to have good coincidence between the predictions of the different interactions and the SPA at low $`\omega _\gamma `$ values. We get (in MeV units) $`\mathrm{\Lambda }_{OBQA}=750,700,600`$, $`\mathrm{\Lambda }_{SEP}=700,600,500`$, and $`\mathrm{\Lambda }_{NEW}=550,500,450`$, for incident pion energies $`T_{lab}=269,298,324`$ respectively. Observe that the value of $`\mathrm{\Lambda }`$ found in the case of the SEP interaction for $`T_{lab}=298MeV`$ is consistent with the one previously found for pion photo-production at $`T_{lab}=300MeV`$, while the OBQA values are roughly consistent with the form factors used in ref., which corresponds to a monopole form-factor with $`\mathrm{\Lambda }800MeV`$. Fig.5 $`\pi N\gamma `$ cross section for $`T_{lab}=269`$, $`298`$ and $`324MeV`$ and $`G_{14}\theta _\gamma =103^0,\varphi _\gamma =180^0`$ in EXP, calculated in the DMA I for the different T-matrices. We also include the SPA cross section and the measured values. Results for the cross section within the DMA II are shown in Fig. 6 for different photon detectors, $`G_{14}\theta _\gamma =103^0,\varphi _\gamma =180^0`$;$`G_{11}\theta _\gamma =160^0,\varphi _\gamma =180^0`$; and $`G_7\theta _\gamma =120^0,\varphi _\gamma =0^0`$. We will fix $`M_\mathrm{\Delta }=1232MeV`$ and $`\mathrm{\Gamma }=110MeV`$ as reported in previous experimental works. The coupling constant $`f_{N\mathrm{\Delta }\pi }`$, was determined fitting the low-energy photon emission region (elastic scattering). We obtained $`f_{N\mathrm{\Delta }\pi }=0.24`$, which is in the range of previous works. Independently, the $`\mathrm{\Delta }`$ magnetic moment $`\kappa _\mathrm{\Delta }`$ was shifted in order to show how the high energy photon region can be fitted. Fig.6 $`\pi N\gamma `$ cross section for $`T_{lab}=269MeV`$ and $`G_7\theta _\gamma =120^0,\varphi _\gamma =0^0`$, $`G_{11}\theta _\gamma =160^0,\varphi _\gamma =0^0`$, and $`G_{14}\theta _\gamma =103^0,\varphi _\gamma =180^0`$ in EXP, calculated in the DMA II for two different values of $`\kappa _\mathrm{\Delta }`$. In almost all the cases, the predictions of the DMA I (fig.5) lies above the experimental cross section and the SPA for energies $`\omega _\gamma >20MeV`$. One of the reasons for this may be the use of an overall form factor to cure the gauge invariance problems. The total bremsstrahlung amplitude is built up, as can be seen from eqs. (9) and (10), by adding different components. It is not expected that the common form factor works satisfactorily by adding up these components as it does the one used to generate the individual T-matrices, which change their values from vertex to vertex. The comparison between the results of the SPA and DMA I schemes shows that the additional off-shell effects, added coherently to the lowest order contributions, may have important contributions since they do not cancel exactly with the derivate terms of the T-matrix appearing from the soft-photon expansion. From fig.5 we can check that the SEP interaction provides the closest results to the experimental cross section with a departure starting for $`\omega _\gamma >40MeV`$. This indicates that the dynamical model involved in the SEP interaction gives the smallest off-shell effects. On the other hand, the strongest off-shell effects appear in the OBQA model. This conclusion agrees with a previous study on observables in pion photo-production experiments. As it was discussed in section 2, the off-shell contributions to the external and internal amplitudes within the SPA cancel each other, thus we cannot study these effects within this approximation. In addition since we get gauge invariance in the SPA by adjusting the internal amplitude, the gauge-invariant electromagnetic currents remain hidden. In the DMA I approach, these cancellation must occur explicitly between the different components of the amplitude (the so-called born,pre,post contributions). The departure of the different T-matrices from the SPA can be used to estimate the size of unbalanced off-shell terms and provide a test of their off-shell behaviors. Also, since the electromagnetic gauge-invariant current is constructed explicitly from effective Lagrangians, we can use the radiative $`\pi N`$ reaction to study the relevance of the degrees of freedom and the parameters involved in this dynamical model. On the other hand from fig.6 we can see as this reaction within the DMA II (tree-label + complex $`M_\mathrm{\Delta }`$), is sensible to the changes of $`\kappa _\mathrm{\Delta }`$. This model seems appropiated to analyze the anomalous magnetic moment of the $`\mathrm{\Delta }`$ resonance, since within the SPA the the ’enforcing’ gauge invariant procedure does not fix this term, which is self-gauge invariant. In order to give a realistic value , a more carefull fitting procedure of $`M_\mathrm{\Delta }`$, $`\mathrm{\Gamma }`$, and $`\kappa _\mathrm{\Delta }`$, must be done. ACKNOWLEDGMENTS The work of A. Mariano was supported in part by Conacyt (México) through the Fondo de Cátedras Patrimoniales de Excelencia Nivel II, and Conicet (Argentina). He is also grateful to G. López Castro for important discussions.
warning/0004/math0004181.html	ar5iv	text	# The Connes-Higson construction is an isomorphism ## 1. Introduction The fundamental homotopy functors on the category of separable $`C^{}`$-algebras are all based on extensions — either a priori or a posteriori . So also the $`E`$-theory of Connes and Higson; in the words of the founders: ’La $`E`$-theorie est ainsi le quotient par homotopie de la théorie des extensions’, cf. \[CH\]. The connection between the asymptotic homomorphisms which feature explicitly in the definition of $`E`$-theory, and $`C^{}`$-extensions, appears as a fundamental construction which associates an asymptotic homomorphism $`SAB`$ to a given extension of $`A`$ by $`B`$. While it is easy to see that the homotopy class of the asymptotic homomorphism only depends on the homotopy class of the extension it is not so easy to decide if the converse is also true; if the extensions must be homotopic when the asymptotic homomorphisms which they give rise to via the Connes-Higson construction are. A part of the main result in the present paper asserts that this is the case when $`A`$ is a suspension and $`B`$ is stable. Rather unexpectedly it turned out that the methods we developed for this were also able to characterize $`E`$-theory as the quotient of all extensions of $`SA`$ by $`B`$ by an algebraic relation which is very similar to the algebraic relation which has been considered on the set of extensions since the way-breaking work of Brown, Douglas and Fillmore, \[BDF\]. Recall that in the BDF-approach two $`C^{}`$-extensions are identified when they become unitarily equivalent after addition by extensions which are split, meaning that the quotient map admits a $``$-homomorphism as a right-inverse. In the algebraic relation, on the set of all $`C^{}`$-extensions of $`SA`$ by $`B`$, which we will show gives rise to $`E`$-theory, two extensions are identified when they become unitarily equivalent after addition by extensions which are asymptotically split, where we call an extension asymptotically split when there is an asymptotic homomorphism $`\pi =(\pi _t)_{t[1,\mathrm{})}:AE`$ such that $`p\pi _t=\mathrm{id}_A`$ for all $`t`$. We emphasize that with this relation all extensions of $`SA`$ by $`B`$ admit an inverse. In contrast, Kirchberg has shown, \[Ki\], that the unitary equivalence classes of extensions of $`SA`$ by $`𝒦`$, modulo the split extensions, do not form a group when $`A`$ is the reduced group $`C^{}`$-algebra of a discrete non-amenable subgroup of a connected Lie-group. Since our results show that the algebraic relation we have just described is the same as homotopy, our main result can also be considered as a result on homotopy invariance and it is therefore noteworthy that the proof is self-contained, and in particular does not depend on the homotopy invariance results of Kasparov. Since there is also an equivariant version of $`E`$-theory, \[GHT\], which is being used in connection with the Baum-Connes conjecture, we formulate and prove our results in the equivariant case. With the present technology this does not require much additional work, but since some of the material which we shall build on does not explicitly consider the equivariant setting, notably \[DL\] and \[H-LT\], there are a few places where we leave the reader to check that the results from these sources can be adapted to the equivariant case. ## 2. An alternative to the BDF extension group Let $`G`$ be a locally compact, $`\sigma `$-compact group, and let $`A`$ and $`B`$ be separable $`G`$-algebras, i.e. separable $`C^{}`$-algebras with a pointwise norm-continuous action of $`G`$ by automorphisms. Assume also that $`B`$ is weakly stable, i.e. that $`B`$ is equivariantly isomorphic to $`B𝒦`$ where $`𝒦`$ denotes the compact operators of $`l_2`$ with the trivial $`G`$-action. Let $`M(B)`$ denote the multiplier algebra of $`B`$, $`Q(B)=M(B)/B`$ the corresponding corona algebra and $`q_B:M(B)Q(B)`$ the quotient map. Then $`G`$ acts by automorphisms on both $`M(B)`$ and $`Q(B)`$<sup>1</sup><sup>1</sup>1These actions are not pointwise normcontinuous in general.. It follows from \[Th1\] that we can identify the set of equivariant $``$-homomorphisms, $`\mathrm{Hom}_G(A,Q(B))`$, from $`A`$ to $`Q(B)`$ with the set of $`G`$-extensions of $`A`$ by $`B`$. Two $`G`$-extensions $`\phi ,\psi :AQ(B)`$ are unitarily equivalent when there is a unitary $`wM(B)`$ such that $`q_B(w)Q(B)`$ is $`G`$-invariant and $`\mathrm{Ad}q_B(w)\phi =\psi `$. Since $`B`$ is weakly stable the set of unitary equivalence classes of extensions of $`A`$ by $`B`$ form a semi-group; the addition is obtained by choosing two $`G`$-invariant isometries $`V_1,V_2M(B)`$ such that $`V_1V_1^{}+V_2V_2^{}=1`$ and setting $`\phi \psi =q_B(V_1)\phi ()q_B(V_1)^{}+q_B(V_2)\psi ()q_B(V_2)^{}`$. A $`G`$-extension $`\phi :AQ(B)`$ will be called asymptotically split when there is an asymptotic homomorphism $`\pi =\{\pi _t\}_{t[1,\mathrm{})}:AM(B)`$ such that $`q_B\pi _t=\phi `$ for all $`t`$. All asymptotic homomorphisms we consider in this paper will be assumed to be equivariant in the sense that $`lim_t\mathrm{}g\pi _t(a)\pi _t(ga)=0`$ for all $`aA`$ and $`gG`$. As in \[MT2\] we say that a $`G`$-extension $`\phi :AQ(B)`$ is semi-invertible when there is a $`G`$-extension $`\psi \mathrm{Hom}_G(A,Q(B))`$ such that $`\phi \psi :AQ(B)`$ is asymptotically split. Two semi-invertible extensions, $`\phi ,\psi `$, are called stably unitary equivalent when they become unitarily equivalent after addition by asymptotically split extensions, i.e. when there is an asymptotically split extension $`\lambda `$ such that $`\phi \lambda `$ is unitarily equivalent to $`\psi \lambda `$. This is an equivalence relation on the subset of semi-invertible extensions in $`\mathrm{Hom}_G(A,Q(B))`$ and the corresponding equivalence classes form an abelian group which we denote by $`\mathrm{Ext}^{1/2}(A,B)`$. For any locally compact space $`X`$ we consider $`C_0(X)A`$ as a $`G`$-algebra with the trivial $`G`$-action on the tensor factor $`C_0(X)`$. When $`X=(0,1]`$ we denote $`C_0(0,1]A`$ by $`\mathrm{cone}(A)`$. Similarly, we set $`SA=C_0(0,1)A`$. ###### Lemma 2.1. Let $`\lambda :\mathrm{cone}(A)Q(B)`$ be a $`G`$-extension. It follows that there is an asymptotic homomorphism $`\pi =(\pi _t)_{t[1,\mathrm{})}:\mathrm{cone}(A)M_2(M(B))`$ such that $$q_{M_2(B)}\pi _t=\left(\begin{array}{cc}\lambda & \\ & 0\end{array}\right)$$ for all $`t[1,\mathrm{})`$. ###### Proof. The proof is based on an idea of Voiculescu, cf. \[V\]. Let $`\mu :\mathrm{cone}(A)M(B)`$ be a continuous, self-adjoint and homogeneous lift of $`\lambda `$ such that $`\mu (x)2x`$ for all $`x\mathrm{cone}(A)`$. Such $`\mu `$ exists by the Bartle-Graves selection theorem, cf. \[L\]. Define $`\phi _s:\mathrm{cone}(A)\mathrm{cone}(A)`$ such that $`\phi _s(f)(t)=f((1s)t),s[0,1]`$. Choose continuous functions $`f_i:[1,\mathrm{})[0,1],i=0,1,2,\mathrm{}`$, such that 1. $`f_0(t)=0`$ for all $`t[1,\mathrm{})`$, 2. $`f_nf_{n+1}`$ for all $`n`$, 3. for each $`n`$, there is an $`m_n`$ such that $`f_i(t)=1`$ for all $`im_n`$, and all $`t[1,n+1]`$, 4. $`lim_t\mathrm{}\mathrm{max}_i\|f_i(t)f_{i+1}(t)\|=0`$. Let $`F_1F_2F_3\mathrm{}`$ be an increasing sequence of finite subsets with dense union in $`\mathrm{cone}(A)`$. Write $`G=_nK_n`$ where $`K_1K_2K_3\mathrm{}`$ are compact subsets of $`G`$. For each $`n`$, choose $`m_n`$ as in 3). We may assume that $`m_{n+1}>m_n`$. By Lemma 1.4 of \[K\] we can choose elements $$X_0^nX_1^nX_2^n\mathrm{}$$ in $`B`$ such that $`0X_i^n1`$ for all $`i`$ and $`X_i^n=0`$ for $`im_n`$, and 1. $`X_i^nX_{i+1}^n=X_{i+1}^n`$ for all $`i`$, 2. $`X_i^nbb\frac{1}{n}`$ for all $`i=0,1,2,\mathrm{},m_n1`$, and all $`bS_n`$, 3. $`X_i^nyyX_i^n\frac{1}{n}`$ for all $`i`$ and all $`yL_n`$, 4. $`gX_i^nX_i^n\frac{1}{n},gK_n`$, for all $`i`$, 5. $`X_i^n(g\mu (a)\mu (ga))(g\mu (a)\mu (ga))\frac{1}{n},gK_n,aF_n`$, for all $`i=0,1,2,\mathrm{},m_n1`$, where $`L_n`$ and $`S_n`$ are the compact sets $`L_n=\{\mu (\phi _s(a)):s[0,1],aF_n\}`$ and $$\begin{array}{cc}& S_n=\{\mu (\phi _s(a))+\mu (\phi _s(b))\mu (\phi _s(a+b)):a,bF_n,s[0,1]\}\hfill \\ & \{\mu (\phi _s(ab))\mu (\phi _s(a))\mu (\phi _s(b)):a,bF_n,s[0,1]\}.\hfill \end{array}$$ Since we choose the $`X`$’s recursively we can arrange that $`X_i^{n+1}X_k^n=X_k^n`$ for all $`k`$ and all $`im_{n+1}`$. By connecting first $`X_0^n`$ to $`X_0^{n+1}`$ via the straight line between them, then $`X_1^n`$ to $`X_1^{n+1}`$ via a straight line, then $`X_2^n`$ to $`X_2^{n+1}`$ etc., we obtain norm-continuous pathes, $`X(t,i),t[n,n+1],i=0,1,2,3,\mathrm{}`$, in $`B`$ such $`X(n,i)=X_i^n,X(n+1,i)=X_i^{n+1}`$ for all $`i`$ and 1. $`X(t,i)X(t,i+1)=X(t,i+1),t[n,n+1]`$, for all $`i`$, 2. $`X(t,i)bb\frac{1}{n}`$ for all $`i=0,1,2,\mathrm{},m_n1,t[n,n+1]`$ and all $`bS_n`$, 3. $`X(t,i)yyX(t,i)\frac{1}{n}`$ for all $`i`$, all $`t[n,n+1]`$ and all $`yL_n`$, 4. $`gX(t,i)X(t,i)\frac{1}{n},gK_n,t[n,n+1]`$, for all $`i`$, 5. $`X(t,i)(g\mu (a)\mu (ga))(g\mu (a)\mu (ga))\frac{1}{n},gK_n,aF_n,t[n,n+1]`$, for all $`i=0,1,\mathrm{},m_n1`$. In addition, $`X(t,i)=0,im_{n+1},t[n,n+1]`$. Let $`l_2(B)`$ denote the Hilbert $`B`$-module of sequences $`(b_1,b_2,b_3,\mathrm{})`$ in $`B`$ such that $`_{i=1}^{\mathrm{}}b_i^{}b_i`$ converges in norm. Writing an element $`(b_1,b_2,b_3,\mathrm{})l_2(B)`$ as the sum $`_{i=0}^{\mathrm{}}b_ie_i`$ we define a representation $`V`$ of $`G`$ on $`l_2(B)`$ such that $`V_g(_{i=0}^{\mathrm{}}b_ie_i)=_{i=0}^{\mathrm{}}(gb_i)e_i`$. Then $`G`$ acts by automorphisms on $`𝕃(l_2(B))`$ ( = the adjoinable operators on $`l_2(B)`$) such that $`gm=V_gmV_{g^1}`$. Set $$T_t=\left(\begin{array}{cccc}\sqrt{1X(t,0)}& \sqrt{X(t,0)X(t,1)}& \sqrt{X(t,1)X(t,2)}& \mathrm{}\\ 0& 0& 0& \mathrm{}\\ 0& 0& 0& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\end{array}\right)𝕃(l_2(B)).$$ Then $`P_t=T_t^{}T_t`$ is a projection in $`𝕃(l_2(B))`$ since $`T_tT_t^{}`$ clearly is. Note that $`P_t`$ is tri-diagonal because of condition a) above, and that the entries of $`P_t`$ are all in $`B`$, with the notable exception of the $`1\times 1`$-entry which is equal to $`1`$ modulo $`B`$. We define $`\delta _t:\mathrm{cone}(A)𝕃(l_2(B))`$ by $$\delta _t(a)(\underset{i=0}{\overset{\mathrm{}}{}}b_ie_i)=\underset{i=0}{\overset{\mathrm{}}{}}\mu (\phi _{f_i(t)}(a))b_ie_i.$$ Set $`\pi _t(a)=P_t\delta _t(a)P_t`$ for $`a\mathrm{cone}(A)`$ and $`t[1,\mathrm{})`$. We assert that $`\pi =(\pi _t)_{t[1,\mathrm{})}`$ is an asymptotic homomorphism. By using the continuity of $`\mu `$ and that $`\{\phi _s(a):s[0,1]\}`$ is a compact set for fixed $`a`$, it follows readily that the family of maps $`a\pi _t(a),t[1,\mathrm{})`$, is an equicontinuous family. Since each $`\pi _t`$ is self-adjoint and homogeneous, it suffices therefore to take an $`n`$ and elements $`a,bF_n`$, $`gK_n`$, and check that $$\underset{t\mathrm{}}{lim}P_t\delta _t(a)P_t\delta _t(b)P_tP_t\delta _t(ab)P_t=0,$$ $$\underset{t\mathrm{}}{lim}P_t\delta _t(a+b)P_tP_t\delta _t(a)P_tP_t\delta _t(b)P_t=0,$$ and $$\underset{t\mathrm{}}{lim}P_t\delta _t(ga)P_tg(P_t\delta _t(a)P_t)=0.$$ The first two limits are zero by 4), b) and c), the third by d) and e). For each $`a,t`$, $`P_t\delta _t(a)P_t=\mathrm{diag}(\mu (a),0,0,\mathrm{})`$ modulo $`𝕂(l_2(B))`$ ( = the ideal of ’compact’ operators on $`l_2(B)`$). Since $`B`$ is weakly stable there is an equivariant isomorphism $`l_2(B)BB`$ of Hilbert $`B`$-modules which leaves the first coordinate invariant. We can therefore transfer $`\pi `$ to an asymptotic homomorphism $`\pi =(\pi _t)_{t[1,\mathrm{})}:\mathrm{cone}(A)𝕃(BB)=M_2(M(B))`$ with the stated property. Two $`G`$-extensions $`\phi ,\psi \mathrm{Hom}_G(A,Q(B))`$ are strongly homotopic when there is a path $`\mathbb{\Phi }_t\mathrm{Hom}_G(A,Q(B)),t[0,1]`$, such that $`\mathbb{\Phi }_0=\phi ,\mathbb{\Phi }_1=\psi `$ and $`t\mathbb{\Phi }_t(a)`$ is continuous for all $`aA`$. ###### Theorem 2.2. Let $`\phi :AQ(B)`$ be a $`G`$-extension which is strongly homotopic to $`0`$ in $`\mathrm{Hom}_G(A,Q(B))`$. It follows that there is an asymptotic homomorphism $`\pi =(\pi _t)_{t[1,\mathrm{})}:AM_2(M(B))`$ such that $$q_{M_2(B)}\pi _t=\left(\begin{array}{cc}\phi & \\ & 0\end{array}\right)$$ for all $`t[1,\mathrm{})`$. ###### Proof. Since $`\phi `$ is strongly homotopic to $`0`$ there is an equivariant $``$-homomorphism $`\mu :A\mathrm{cone}(D)`$, where $`DQ(B)`$ is a separable $`G`$-algebra containing $`\phi (A)`$, and an equivariant $``$-homomorphism $`\lambda :\mathrm{cone}(D)Q(B)`$ such that $`\phi =\lambda \mu `$. Apply Lemma 2.1 to $`\lambda `$. ∎ ###### Corollary 2.3. Every $`G`$-extension $`\phi :SAQ(B)`$ is semi-invertible. ###### Proof. Let $`\alpha \mathrm{Aut}SA`$ be the automorphism of $`SA`$ given by $`\alpha (f)(t)=f(1t)`$. It is wellknown that $`\phi (\phi \alpha )`$ is strongly homotopic to $`0`$. Hence $`\phi (\phi \alpha )0`$ is asymptotically split by Theorem 2.2. ∎ Because of Corollary 2.3 we drop the superscript $`1/2`$ and write $`\mathrm{Ext}(SA,B)`$ instead of $`\mathrm{Ext}^{1/2}(SA,B)`$. ###### Lemma 2.4. Let $`\phi ,\psi :SAQ(B)`$ be two $`G`$-extensions which are strongly homotopic. It follows that $`\phi `$ and $`\psi `$ are stably unitarily equivalent. ###### Proof. It follows from Theorem 2.2 that $`\lambda _1=(\phi \alpha )\phi 0`$ and $`\lambda _2=(\phi \alpha )\psi 0`$ are both asymptotically split. Since $`\psi \lambda _1`$ and $`\phi \lambda _2`$ are unitarily equivalent, the conclusion follows because infinite direct sums are well-defined for asymptotically split extensions. ∎ Set $`IB=C[0,1]B`$ and let $`e_t:IBB`$ denote evaluation at $`t[0,1]`$ and note that $`e_t`$ defines a equivariant $``$-homomorphisms $`M(IB)M(B)`$ and $`Q(IB)Q(B)`$ which we again denote by $`e_t`$. Two $`G`$-extensions $`\phi ,\psi \mathrm{Hom}_G(A,Q(B))`$ are homotopic when there is a $`G`$-extension $`\mathbb{\Phi }\mathrm{Hom}_G(A,Q(IB))`$ such that $`e_0\mathbb{\Phi }=\phi `$ and $`e_1\mathbb{\Phi }=\psi `$. As in \[MT2\] we denote the set of homotopy classes of $`G`$-extensions by $`\mathrm{Ext}(A,B)_h`$. In general this is merely an abelian semigroup, but $`\mathrm{Ext}(SA,B)_h`$ is a group. The Connes-Higson construction associates to any $`G`$-extension $`\phi \mathrm{Hom}_G(A,Q(B))`$ an asymptotic homomorphism $`CH(\phi ):SAB`$ in the following way, cf. \[CH\], \[GHT\]: By Lemma 1.4 of \[K\] or Lemma 5.3 of \[GHT\] there is a norm-continuous path $`\{u_t\}_{t[1,\mathrm{})}`$ of elements in $`B`$ such that $`0u_t1`$ for all $`t`$, $`lim_t\mathrm{}u_tbb=0`$ for all $`bB`$, $`lim_t\mathrm{}u_tmmu_t=0`$ for all $`mq_B^1(\phi (A))`$ and $`lim_t\mathrm{}gu_tu_t=0`$ for all $`gG`$. From these data $`CH(\phi )`$ is determined up to asymptotic equality as the equicontinuous<sup>2</sup><sup>2</sup>2Equicontinuity of an asymptotic homomorphism $`\pi =(\pi _t)_{t[1,\mathrm{})}:AB`$ means that $`A\times G(a,g)g\pi _t(a),t[1,\mathrm{})`$, is an equicontinuous family of maps. asymptotic homomorphism $`CH(\phi ):SAB`$ which satisfies that $$\underset{t\mathrm{}}{lim}CH(\phi )_t(fa)f(u_t)x=0,xq_B^1(\phi (a)),$$ for all $`fC_0(0,1)`$ and all $`aA`$. Let $`[[SA,B]]`$ denote the abelian group of homotopy classes of asymptotic homomorphisms, $`SAB`$, cf. \[CH\], \[GHT\]. The Connes-Higson construction defines in the obvious way a semi-group homomorphism $`CH:\mathrm{Ext}(A,B)_h[[SA,B]]`$. Since there is a canonical (semi-group) homomorphism $`\mathrm{Ext}^{1/2}(A,B)\mathrm{Ext}(A,B)_h`$ we may also consider the Connes-Higson construction as a homomorphism $`CH:\mathrm{Ext}^{1/2}(A,B)[[SA,B]]`$. Notice that $`\mathrm{Ext}(SA,B)`$ and $`\mathrm{Ext}(SA,B)_h`$ are both abelian groups and the canonical map $`\mathrm{Ext}(SA,B)\mathrm{Ext}(SA,B)_h`$ is a surjective group homomorphism by Corollary 2.3. In Corollary 5.4 below we show that it is an isomorphism. ## 3. On equivalence of asymptotic homomorphisms ###### Lemma 3.1. Let $`A`$ and $`B`$ be separable $`G`$-algebras, $`B`$ weakly stable. Let $`\phi =(\phi _t)_{t[1,\mathrm{})}:AB`$ be an asymptotic homomorphism which is homotopic to $`0`$. It follows that there is an asymptotic homomorphism $`\psi =(\psi _t)_{t[1,\mathrm{})}:AB`$ and a norm-continuous path $`\{W_t\}_{t[1,\mathrm{})}`$ of $`G`$-invariant unitaries in $`M(M_2(B))`$ such that $$\underset{t\mathrm{}}{lim}\left(\begin{array}{cc}\phi _t\left(a\right)& \\ & \psi _t\left(a\right)\end{array}\right)W_t\left(\begin{array}{cc}0& \\ & \psi _t\left(a\right)\end{array}\right)W_t^{}=0$$ for all $`aA`$. ###### Proof. Let $`\mathbb{\Phi }=(\mathbb{\Phi }_t)_{t[1,\mathrm{})}:AIB`$ be an asymptotic homomorphism such that $`e_0\mathbb{\Phi }_t(a)=0,e_1\mathbb{\Phi }_t(a)=\phi _t(a)`$ for all $`t[1,\mathrm{}),aA`$. We may assume that both $`\phi `$ and $`\mathbb{\Phi }`$ are equicontinuous, cf. Proposition 2.4 of \[Th2\]. Let $`F_1F_2F_3\mathrm{}`$ be a sequence of finite subsets with dense union in $`A`$. For each $`n`$ there is $`\delta _n>0`$ with the property that $$e_x\mathbb{\Phi }_t(a)e_y\mathbb{\Phi }_t(a)<\frac{1}{n}$$ when $`\|xy\|<\delta _n,t[1,n],aF_n`$. Choose then a sequence of functions $`f_k:[1,\mathrm{})[0,1]`$ such that $`f_1(t)=1,f_kf_{k+1},\|f_k(t)f_{k+1}(t)\|<\delta _n,t[1,n]`$ for all $`k,n`$ and such that $`f_k\|_{[1,n]}=0`$ for all but finitely many $`k`$’s for all $`n`$. Set $`\lambda _t^n(a)=e_{f_n(t)}\mathbb{\Phi }_t(a)`$ for all $`aA,n,t[1,\mathrm{})`$. Note that $`\lambda _t^i(a)\lambda _t^{i+1}(a)<\frac{1}{n},aF_n,t[1,n]`$, for all $`i`$ and $`n`$. Then $$\mu _t(a)=\mathrm{diag}(\phi _t(a),\lambda _t^1(a),\lambda _t^2(a),\lambda _t^3(a),\mathrm{})𝕂(l_2(B))$$ and $$\delta _t(a)=\mathrm{diag}(0,\lambda _t^1(a),\lambda _t^2(a),\lambda _t^3(a),\mathrm{})𝕂(l_2(B))$$ define asymptotic homomorphisms $`\mu ,\delta :A𝕂(l_2(B))`$. By connecting appropriate permutation unitaries, acting on $`l_2(B)`$ by permutations of $`B`$-coordinates, we get a norm-continuous path of $`G`$-invariant unitaries $`\{S_t\}_{t[1,\mathrm{})}𝕃(l_2(B))`$ such that $$S_t\delta _t(a)S_t^{}=\mathrm{diag}(\lambda _t^1(a),\lambda _t^2(a),\lambda _t^3(a),\mathrm{})$$ for all $`a,t`$. Then $`lim_t\mathrm{}\mu _t(a)S_t\delta _t(a)S_t^{}=0`$ for all $`aA`$. Since $`B`$ is weakly stable there is an isomorphism $`l_2(B)BB`$ of Hilbert $`B,G`$-algebras which fixes the first coordinate. Applying this isomorphism in the obvious way and remembering the identifications $`𝕂(BB)=M_2(B)`$ and $`𝕃(BB)=M(M_2(B))`$ gives the result. ∎ ###### Theorem 3.2. Let $`A`$ and $`B`$ be separable $`G`$-algebras, $`B`$ weakly stable. Assume that $`[[A,B]]`$ is a group. Two asymptotic homomorphisms, $`\phi =(\phi _t)_{t[1,\mathrm{})},\psi =(\psi _t)_{t[1,\mathrm{})}:AB`$, are homotopic if and only if there is an asymptotic homomorphism $`\lambda =(\lambda _t)_{t[1,\mathrm{})}:AB`$ and a norm-continuous path $`\{W_t\}_{t[1,\mathrm{})}`$ of $`G`$-invariant unitaries in $`M(M_2(B))`$ such that $$\underset{t\mathrm{}}{lim}\left(\begin{array}{cc}\phi _t\left(a\right)& \\ & \lambda _t\left(a\right)\end{array}\right)W_t\left(\begin{array}{cc}\psi _t\left(a\right)& \\ & \lambda _t\left(a\right)\end{array}\right)W_t^{}=0$$ for all $`aA`$. ###### Proof. The ’if’ part is easy and the ’only if’ part follows from Lemma 3.1 in the same way as Lemma 2.4 follows from Theorem 2.2. ∎ ###### Lemma 3.3. Let $`B`$ be a weakly stable $`G`$-algebra and $`D_0`$ a separable $`G`$-subalgebra of $`C_b([1,\mathrm{}),B)`$. Let $`V_1,V_2,\mathrm{},V_NM(B)`$ be $`G`$-invariant isometries. There is then a weakly stable separable $`G`$-subalgebra $`D`$ of $`C_b([1,\mathrm{}),B)`$ such that $`V_iDV_i^{}DD_0D`$ for all $`i=1,2,\mathrm{},N`$. ###### Proof. Since $`B`$ is weakly stable we can write $`B=B𝒦`$ with $`G`$ acting trivially on the tensor-factor $`𝒦`$. We embed $`𝒦`$ into $`M(B𝒦)`$ via $`x1_Bx`$. Let $`\{f_n\}C_b([1,\mathrm{}),B𝒦)`$ be a dense sequence in $`D_0`$. For each $`n`$ there is a function $`g_nC_b([1,\mathrm{}),𝒦)`$ such that $`g_nf_nf_n<\frac{1}{n}`$. Let $`E_{00}`$ be the $`C^{}`$-algebra generated by $`\{g_n\}_{n=1}^{\mathrm{}}`$. Then $`E_{00}C_b([1,\mathrm{}),𝒦)C_b([1,\mathrm{}),B^+𝒦)`$. Consider a positive element $`fE_{00}`$ and an $`ϵ>0`$. Set $`U_j=]j,j+2[[1,\mathrm{}[,j=0,1,2,\mathrm{}`$. We can then find a sequence $`p_0p_1p_2\mathrm{}`$ of projections in $`𝒦`$ such that $$\underset{x\overline{U_j}}{sup}p_jf(x)p_jf(x)<ϵ.$$ Let $`\{h_j\}`$ be a partition of unity in $`C_b[1,\mathrm{})`$ subordinate to the cover $`\{U_j\}`$ and set $`g(t)=_{j=0}^{\mathrm{}}h_j(t)p_jf(t)p_j`$. Then $`gC_b([1,\mathrm{}),𝒦)`$, $`g0,gf<ϵ`$. For each $`j`$ we choose a partial isometry $`v_j𝒦`$ such that $`v_jv_j^{}=p_{j+2}`$, $`v_j^{}v_jp_{j+2}=0`$ and $`v_j^{}v_jv_k^{}v_k=0,k<j`$. Set $`h(t)=_{j=0}^{\mathrm{}}\sqrt{h_j(t)}v_j`$. Then $`hh^{}g=g`$ and $`h^{}hg=0`$. It follows that we can find a sequence $`E_{00}=X_1X_2X_3\mathrm{}`$ of separable $`C^{}`$-subalgebras of $`C_b([1,\mathrm{}),𝒦)`$ and for each $`n`$ have a dense sequence $`\{f_1,f_2,\mathrm{}\}`$ in the positive part of $`X_n`$ and elements $`\{v_1,v_2,\mathrm{}\}`$ in $`X_{n+1}`$ such that $`f_kv_k^{}v_k<\frac{1}{k}`$ and $`v_k^{}v_kv_kv_k^{}=0`$ for all $`k`$. It follows then from Proposition 2.2 and Theorem 2.1 of \[HR\] that $`E_0=\overline{_nX_n}`$ is a separable stable $`C^{}`$-subalgebra of $`C_b([1,\mathrm{}),𝒦)`$ such that $`E_{00}E_0`$. Note that $`E_0`$ contains a sequence $`\{r_n\}`$ with the property that $`lim_n\mathrm{}r_nx=x`$ for all $`xD_0`$ since $`E_{00}`$ does. Set $`W=\{V_1,V_2,\mathrm{},V_N\}\{V_1^{},V_2^{},\mathrm{},V_N^{}\}`$. By repeating the above argument with $`D_0`$ substituted by the $`G`$-algebra $`D_1`$ generated by $`D_0WD_0E_0D_0`$, we get a stable $`C^{}`$-subalgebra $`E_1C_b([1,\mathrm{}),𝒦)`$ which contains a sequence $`\{r_n\}`$ such that $`lim_n\mathrm{}r_ny=y`$ for all $`yD_1`$. It is clear from the construction that we can arrange that $`E_0E_1`$. We can therefore continue this procedure to obtain sequences of separable $`G`$-algebras, $`D_0D_1D_2D_3\mathrm{}`$ in $`C_b([1,\mathrm{}),B𝒦)`$, and $`E_0E_1E_2E_3\mathrm{}`$ in $`C_b([1,\mathrm{}),𝒦)C_b([1,\mathrm{}),B^+𝒦)`$ such that each $`E_n`$ is stable and contains a sequence $`\{r_k\}`$ such that $`lim_k\mathrm{}r_kx=x,xD_n`$, and $`D_nWD_nE_nD_nD_{n+1}`$ for all $`n`$. Set $`E_{\mathrm{}}=\overline{_nE_n}`$ and $`D=\overline{_nD_n}`$. It follows from Corollary 4.1 of \[HR\] that $`E_{\mathrm{}}`$ is stable. By construction $`V_iDV_i^{}DD`$ for all $`i`$ and $`E_{\mathrm{}}DD`$. The last property ensures that $`D`$ is an ideal in the $`G`$-algebra $`E`$ generated by $`E_{\mathrm{}}`$ and $`D`$. There is therefore a $``$-homomorphism $`\lambda :E_{\mathrm{}}M(D)`$. By construction an approximate unit for $`E_{\mathrm{}}`$ is also an approximate unit for $`D`$ so $`\lambda `$ extends to a $``$-homomorphism $`\lambda :M(E_{\mathrm{}})M(D)`$ which is strictly continuous on the unit ball of $`M(E_{\mathrm{}})`$. Since $`E_{\mathrm{}}`$ is stable there is a sequence $`P_i,i=1,2,\mathrm{}`$, of orthogonal and Murray-von Neumann equivalent projections in $`M(E_{\mathrm{}})`$ which sum to $`1`$ in the strict topology. Then $`Q_i=\lambda (P_i),i=1,2,\mathrm{}`$, is a sequence of orthogonal and Murray-von Neumann equivalent projections in $`M(D)`$ which sum to $`1`$ in the strict topology. Since $`E_{\mathrm{}}`$ consists entirely of $`G`$-invariant elements it follows that all the $`Q_i`$’s are $`G`$-invariant. Consequently $`DQ_1DQ_1𝒦`$ as $`G`$-algebras, proving that $`D`$ is weakly stable. ∎ Two asymptotic homomorphisms $`\phi =(\phi _t)_{t[1,\mathrm{})},\psi =(\psi _t)_{t[1,\mathrm{})}:AB`$ will be called equi-homotopic when there is a family $`\mathbb{\Phi }^\lambda =(\mathbb{\Phi }_t^\lambda )_{t[1,\mathrm{})}:AB,\lambda [0,1]`$, of asymptotic homomorphisms such that the family of maps, $`[0,1]\lambda \mathbb{\Phi }_t^\lambda (a),t[1,\mathrm{})`$, is equicontinuous for each $`aA`$. ###### Theorem 3.4. Let $`A`$ and $`B`$ be separable $`G`$-algebras, $`B`$ weakly stable. Let $`\phi =(\phi _t)_{t[1,\mathrm{})},\psi =(\psi _t)_{t[1,\mathrm{})}:SAB`$ be asymptotic homomorphisms. Then the following are equivalent: 1. $`\phi `$ and $`\psi `$ are homotopic $`(`$i.e. $`[\phi ]=[\psi ]`$ in $`[[SA,B]]`$$`)`$. 2. $`\phi `$ and $`\psi `$ are equi-homotopic. 3. There is an asymptotic homomorphism $`\lambda =(\lambda _t)_{t[1,\mathrm{})}:SAB`$ and a norm-continuous path $`\{W_t\}_{t[1,\mathrm{})}`$ of $`G`$-invariant unitaries in $`M(M_2(B))`$ such that $$\underset{t\mathrm{}}{lim}\left(\begin{array}{cc}\phi _t\left(a\right)& \\ & \lambda _t\left(a\right)\end{array}\right)W_t\left(\begin{array}{cc}\psi _t\left(a\right)& \\ & \lambda _t\left(a\right)\end{array}\right)W_t^{}=0$$ for all $`aA`$. ###### Proof. The equivalence 1) $``$ 3) follows from Theorem 3.2 and the implication 2) $``$ 1) is trivial, so we need only prove that 1) $``$ 2). To this end, let $`[[SA,B]]^e`$ denote the set of equi-homotopy classes of asymptotic homomorphisms $`SAB`$. Choose $`G`$-invariant isometries $`V_1,V_2M(B)`$ such that $`V_1V_1^{}+V_2V_2^{}=1`$ and define a composition in $`[[SA,B]]^e`$ by $$[\phi ]+[\psi ]=[(V_1\phi _tV_1^{}+V_2\psi _tV_2^{})_{t[1,\mathrm{})}].$$ It follows from Lemma 3.3 that $`[[SA,B]]^e`$ is a group. It suffices therefore to show that the natural map $`[[SA,B]]^e[[SA,B]]`$ has trivial kernel. If $`\phi `$ is an asymptotic homomorphism representing an element in the kernel we conclude from Lemma 3.1 that there is a norm-continuous path $`W_t,t[1,\mathrm{})`$, of $`G`$-invariant unitaries in $`M_2(M(B)))`$ and an asymptotic homomorphism $`\psi `$ such that $$\underset{t\mathrm{}}{lim}\left(\begin{array}{cccc}\phi _t\left(a\right)& & & \\ & \psi _t\left(a\right)& & \\ & & 0& \\ & & & 0\end{array}\right)\left(\begin{array}{cc}W_t& \\ & W_t^{}\end{array}\right)\left(\begin{array}{cccc}0& & & \\ & \psi _t\left(a\right)& & \\ & & 0& \\ & & & 0\end{array}\right)\left(\begin{array}{cc}W_t^{}& \\ & W_t\end{array}\right)=0$$ for all $`aSA`$. By a standard rotation argument we can remove the unitaries $`\left(\begin{array}{cc}W_t& \\ & W_t^{}\end{array}\right)`$ via an equi-homotopy and we see in this way that $`[\phi ]+[\psi ]=[\psi ]`$ in $`[[SA,B]]^e`$. Hence $`[\phi ]=0`$ in $`[[SA,B]]^e`$. Simple examples show that the implications 1) $``$ 2) and 1) $``$ 3) of Theorem 3.4 generally fail in $`[[A,B]]`$. ## 4. Making genuine homomorphisms out of asymptotic ones Let $`A`$ and $`B`$ be separable $`C^{}`$-algebras. Set $$M(B)_G=\{xM(B):Gggx\text{is norm-continuous}\}$$ and $$Q(B)_G=\{xQ(B):Gggx\text{is norm-continuous}\}.$$ Then (1) is a short exact sequence of $`G`$-algebras. (This is not trivial - the surjectivity of the quotient map follows from Theorem 2.1 of \[Th1\].) We are going to construct a map $`\alpha :[[SA,Q(B)_G𝒦]]\mathrm{Ext}(SA,B𝒦)_h`$. The key to this is another variant of the Voiculescu’s tri-diagonal projection trick from \[V\]. Let $`b`$ be a strictly positive element of $`B𝒦`$, $`0b1`$. A unit sequence in $`B𝒦`$ is a sequence $`\{u_n\}_{n=0}^{\mathrm{}}B𝒦`$ such that 1. there is a continuous function $`f_n:[0,1][0,1]`$ which is zero in a neighbourhood of $`0`$ and $`u_n=f_n(b)`$, 2. $`0u_n1`$ for all $`n=0,1,2,3,\mathrm{}`$, 3. $`u_{n+1}u_n=u_n`$ for all $`n`$, 4. $`lim_n\mathrm{}u_nx=x,xB𝒦`$, 5. $`lim_n\mathrm{}gu_nu_n=0,gG`$. Let $`\{e_{ij}\}_{i,j=0}^{\mathrm{}}`$ be the matrix units acting on $`l_2(B𝒦)`$ in the standard way. ###### Lemma 4.1. Let $`𝒰=\{u_n\}`$ be a unit sequence in $`B𝒦`$. Then $$\sqrt{u_0}e_{00}+\underset{j=1}{\overset{\mathrm{}}{}}\sqrt{u_ju_{j1}}e_{0j}$$ converges in the strict topology to a partial isometry $`V`$ in $`𝕃(l_2(B𝒦))`$ such that $`VV^{}=e_{00}`$. ###### Proof. Let $`b=(b_0,b_1,b_2,\mathrm{})=_{i=0}^{\mathrm{}}b_ie_il_2(B𝒦)`$. Then $$\begin{array}{cc}& \underset{j=n}{\overset{m}{}}\sqrt{u_ju_{j1}}e_{0j}(b)^2=\underset{k,j=n}{\overset{m}{}}b_k^{}\sqrt{u_ku_{k1}}\sqrt{u_ju_{j1}}b_j\hfill \\ & \\ & =\underset{k=n}{\overset{m}{}}b_k^{}(u_ku_{k1})b_k+\underset{k=n}{\overset{m1}{}}b_k^{}\sqrt{u_ku_{k1}}\sqrt{u_{k+1}u_k}b_{k+1}+\hfill \\ & \underset{k=n}{\overset{m1}{}}b_{k+1}^{}\sqrt{u_{k+1}u_k}\sqrt{u_ku_{k1}}b_k\hfill \\ & \\ & \underset{k=n}{\overset{m}{}}b_k^{}b_k+2\sqrt{\underset{k=n}{\overset{m1}{}}b_k^{}b_k}\sqrt{\underset{k=n+1}{\overset{m}{}}b_k^{}b_k},\hfill \end{array}$$ proving that $`_{j=1}^{\mathrm{}}\sqrt{u_ju_{j1}}e_{0j}(b)`$ converges in $`l_2(B𝒦)`$. And $$(\underset{j=n}{\overset{m}{}}\sqrt{u_ju_{j1}}e_{0j})^{}(b)^2=\underset{j=n}{\overset{m}{}}b_0^{}(u_ju_{j1})b_0,$$ proving that also $`(_{j=1}^{\mathrm{}}\sqrt{u_ju_{j1}}e_{0j})^{}(b)`$ converges in $`l_2(B𝒦)`$. It follows that $$V=\sqrt{u_0}e_{00}+\underset{j=1}{\overset{\mathrm{}}{}}\sqrt{u_ju_{j1}}e_{0j}$$ exists as a strict limit in $`𝕃(l_2(B𝒦))`$. It it then straightforward to check that $`VV^{}=e_{00}`$. ∎ Let $`P_𝒰=V^{}V`$ and note that $`P_𝒰`$ is tri-diagonal with respect to the matrix units $`\{e_{ij}\}`$. Fix now a continuous and homogeneous section $`\chi `$ for the map $`q_B\mathrm{id}_𝒦:M(B)_G𝒦Q(B)_G𝒦`$. Consider an equicontinuous asymptotic homomorphism $`\phi =(\phi _t)_{t[1,\mathrm{})}:AQ(B)_G𝒦`$. Let $`F_1F_2F_3\mathrm{}`$ be a sequence of finite sets with dense union in $`A`$ and $`K_1K_2K_3\mathrm{}`$ a sequence of compact subsets in $`G`$ such that $`_nK_n=G`$. It is easy to see that there is a unit sequence $`\{u_n\}`$ in $`B𝒦`$ with the following properties : 1. $`u_n\chi (\phi _t(a))\chi (\phi _t(a))u_n\frac{1}{n},aF_n,t[1,n+1]`$, 2. $`(1u_n)(\chi (\phi _t(ab))\chi (\phi _t(a))\chi (\phi _t(b)))\phi _t(ab)\phi _t(a)\phi _t(b)+\frac{1}{n},t[1,n+1],a,bF_n`$, 3. $`(1u_n)(\chi (\phi _t(a+b))\chi (\phi _t(a))\chi (\phi _t(b)))\phi _t(a+b)\phi _t(a)\phi _t(b)+\frac{1}{n},t[1,n+1],a,bF_n`$, 4. $`(1u_n)(g\chi (\phi _t(a))\chi (\phi _t(ga)))g\phi _t(a)\phi _t(ga)+\frac{1}{n},t[1,n],aF_n,gK_n`$. Let $`\{\phi _{t_n}\}_n`$ be a discretization of $`\phi `$, cf. Lemma 5.1 of \[MT1\], such that 1. $`t_nn`$ for all $`n`$. Set $$\stackrel{~}{\phi }(a)=P_𝒰(\underset{j=0}{\overset{\mathrm{}}{}}\chi (\phi _{t_{j+1}}(a))e_{jj})P_𝒰.$$ Then $`\stackrel{~}{\phi }:A𝕃(l_2(B𝒦))`$ is an equivariant $``$-homomorphism modulo $`𝕂(l_2(B𝒦))`$. By identifying $`𝕃(l_2(B𝒦))`$ with $`M(B𝒦)`$, $`𝕂(l_2(B𝒦))`$ with $`B𝒦`$ and the quotient $`𝕃(l_2(B𝒦))/𝕂(l_2(B𝒦))`$ with $`Q(B𝒦)`$, we can consider $`\stackrel{~}{\phi }`$ as a map $`\stackrel{~}{\phi }:AM(B𝒦)`$ with the property that $`q_{B𝒦}\stackrel{~}{\phi }\mathrm{Hom}_G(A,Q(B𝒦))`$. ###### Lemma 4.2. The class of $`q_{B𝒦}\stackrel{~}{\phi }`$ in $`\mathrm{Ext}(A,B𝒦)_h`$ is independent of the choice of unit sequence, subject to the conditions 0)-8), and of the chosen discretization, subject to condition 9), and depends only on the class $`[\phi ]`$ of $`\phi `$ in $`[[A,Q(B)_G𝒦]]`$. ###### Proof. Let $`\{v_n\}`$ be another unit sequence satisfying 0)-8). There is then a unit sequence $`\{w_n\}`$ in $`B𝒦`$ such that $`w_nv_n=v_n,w_nu_n=u_n`$ for all $`n`$. Connect $`u_0`$ to $`w_0`$ by a straight line, then $`u_1`$ to $`w_1`$ by a straight line, etc. This gives a path $`\{w_n^t\}_{t[0,1[}`$ of unit sequences. For each $`t[0,1[`$ we get then a map $`\mu _t:AM(B𝒦)`$ such that $`q_{B𝒦}\mu _t\mathrm{Hom}_G(A,Q(B𝒦))`$ and $`[q_{B𝒦}\mu _0]=[q_{B𝒦}\stackrel{~}{\phi }]`$ in $`\mathrm{Ext}(A,B𝒦)`$. Let $`\delta :AM(B𝒦)`$ be the map obtained from $`\phi `$ as $`\stackrel{~}{\phi }`$ was, but by using $`\{w_n\}`$ instead of $`\{u_n\}`$. Then $`lim_{t1}\mu _t(a)=\delta (a)`$ in the strict topology for all $`aA`$, and $$\underset{t1}{lim}\mu _t(a)\mu _t(b)\mu _t(ab)=\delta (a)\delta (b)\delta (ab),$$ $$\underset{t1}{lim}\mu _t(a+\lambda b)\mu _t(a)\lambda \mu _t(b)=\delta (a+b)\delta (a)\lambda \delta (b),$$ $$\underset{t1}{lim}\mu _t(a^{})\mu _t(a)^{}=\delta (a^{})\delta (a)^{},$$ $$\underset{t1}{lim}\mu _t(ga)g\mu _t(a)=\delta (ga)g\delta (a),$$ in norm for all $`a,bA,\lambda ,gG`$. Hence $`[q_{B𝒦}\delta ]=[q_{B𝒦}\stackrel{~}{\phi }]`$ in $`\mathrm{Ext}(A,B𝒦)_h`$. The same argument with the unit sequence $`\{u_n\}`$ replaced by $`\{v_n\}`$ shows that the class of $`[q_{B𝒦}\stackrel{~}{\phi }]`$ in $`\mathrm{Ext}(A,B𝒦)_h`$ is independent of the choice of unit sequence. Once this is established it is clear that a homotopy of asymptotic homomorphisms $`AQ(B)_G𝒦`$ gives rise, by an appropriate choice of unit sequence, to a homotopy which shows that $`[q_{B𝒦}\stackrel{~}{\phi }]\mathrm{Ext}(A,B𝒦)_h`$ only depends on the homotopy class of $`\phi `$. That $`[q_{B𝒦}\stackrel{~}{\phi }]`$ is also independent of the discretization and only depends on the homotopy class of $`\phi `$ follows in the same way as in Lemma 5.3 and Lemma 5.4 of \[MT1\]. It follows that we have the desired map $`\alpha :[[A,Q(B)_G𝒦]]\mathrm{Ext}(A,B𝒦)_h`$ which is easily seen to be a semi-group homomorphism. ###### Lemma 4.3. Let $`\phi :SAQ(B)𝒦`$ be an equivariant $``$-homomorphism which we consider as a (constant) asymptotic homomorphism. Let $`X`$ be a compact subset with dense span in $`SA`$ and choose a unit sequence $`𝒰=\{u_n\}`$ in $`B𝒦`$ such that $$\sqrt{u_nu_{n1}}\chi (\phi (a))\chi (\phi (a))\sqrt{u_nu_{n1}}<2^n$$ (2) for all $`aX`$ and $$\underset{j=1}{\overset{\mathrm{}}{}}g\sqrt{u_ju_{j1}}\sqrt{u_ju_{j1}}^2<\mathrm{}$$ (3) for all $`gG`$. Then $`[q_{B𝒦}\stackrel{~}{\phi }]=[\iota \phi ]`$ in $`\mathrm{Ext}(SA,B𝒦)`$, where $`\iota :Q(B)_G𝒦Q(B𝒦)_G`$ is the natural embedding. ###### Proof. $`\stackrel{~}{\phi }`$ has the form $`\stackrel{~}{\phi }(a)=P_𝒰(_{j=0}^{\mathrm{}}\chi (\phi (a))e_{jj})P_𝒰`$. Let $`V𝕃(l_2(B𝒦))`$ be the partial isometry defining $`P_𝒰`$ and note that $`gVV𝕂(l_2(B𝒦))`$ for all $`gG`$ because of (3). Thus $$\left(\begin{array}{cc}V& 1VV^{}\\ 1V^{}V& V^{}\end{array}\right)$$ is a unitary in $`M_2(𝕃(l_2(B𝒦)))`$ which is $`G`$-invariant modulo $`M_2(𝕂(l_2(B𝒦)))`$ and satisfies that $$\left(\begin{array}{cc}V& 1VV^{}\\ 1V^{}V& V^{}\end{array}\right)\left(\begin{array}{cc}\stackrel{~}{\phi }& 0\\ 0& 0\end{array}\right)\left(\begin{array}{cc}V^{}& 1V^{}V\\ 1VV^{}& V\end{array}\right)=\left(\begin{array}{cc}\phi _0& 0\\ 0& 0\end{array}\right),$$ where $`\phi _0(a)=(\sqrt{u_0}\chi (\phi (a))\sqrt{u_0}+_{j=1}^{\mathrm{}}\sqrt{u_ju_{j1}}\chi (\phi (a))\sqrt{u_ju_{j1}})e_{00}`$. Thanks to (2) we have that $$\underset{j=1}{\overset{\mathrm{}}{}}\sqrt{u_ju_{j1}}\chi (\phi (a))\sqrt{u_ju_{j1}}(u_ju_{j1})\chi (\phi (a))<\mathrm{}$$ for all $`aX`$. Since $`_{j=1}^{\mathrm{}}(u_ju_{j1})\chi (\phi (a))+u_0\chi (\phi (a))=\chi (\phi (a))`$ (with convergence in the strict topology) we find that $`\phi _0(a)=\chi (\phi (a))e_{00}`$ modulo $`𝕂(l_2(B𝒦))`$ for all $`aX`$, and hence in fact for all $`aSA`$. This proves the lemma. ## 5. The main results Since $`A`$ is separable, $`[[SA,X𝒦]]=\underset{}{\mathrm{lim}}_D[[SA,D𝒦]]`$ for any $`G`$-algebra $`X`$, when we take the limit over all separable $`G`$-subalgebras $`D`$ of $`X`$. It follows from \[DL\] that the suspension map $`S:[[SA,X𝒦]][[S^2A,SX𝒦]]`$ is an isomorphism.<sup>3</sup><sup>3</sup>3Dadarlat and Loring did not consider the equivariant theory in \[DL\], but it is easy to check that their arguments carry over unchanged. Hence $`[[SA,𝒦]]`$ is a homotopy invariant and half-exact functor on the category of $`G`$-algebras (and not only separable $`G`$-algebras). There is therefore a map $$:[[SA,SQ(B)_G𝒦]][[SA,B𝒦]]$$ arising as the boundary map coming from the extension (1), cf. e.g. \[GHT\]. Well-known arguments from the K-theory of $`C^{}`$-algebras, cf. \[Bl\], show that $`[[SA,SM(B)_G𝒦]]=[[SA,M(B)_G𝒦]]=0`$, so the six-terms exact sequence obtained by applying $`[[SA,𝒦]]`$ to (1) shows that $``$ is an isomorphism. For any $`G`$-algebra $`D`$ we let $`s:DD𝒦`$ be the stabilising $``$-homomorphism given by $`s(d)=de`$ for some minimal projection $`e𝒦`$. Since $`B`$ is weakly stable there is an equivariant $``$-isomorphism $`\gamma _0:B𝒦B`$ such that $`s\gamma _0:B𝒦B𝒦`$ is equivariantly homotopic to $`\mathrm{id}_{B𝒦}`$. Let $`\gamma :Q(B𝒦)_GQ(B)_G`$ the $``$-isomorphism induced by $`\gamma _0`$. ###### Lemma 5.1. The composition of the maps is the identity. ###### Proof. We are going to use Theorem 2.3 of \[H-LT\].<sup>4</sup><sup>4</sup>4The equivariant theory was not explicitly considered in \[H-LT\], but all arguments carry over unchanged. Let $`x=s_{}([\mathrm{id}_{SB}])[[SB,SB𝒦]]`$, where $`[\mathrm{id}_{SB}][[SB,SB]]`$ is the element represented by the identity map of $`SB`$ and $`s:SBSB𝒦`$ is the stabilising $``$-homomorphism. By Theorem 2.3 of \[H-LT\] it suffices to identify the image of $`x`$ under the Bott-periodicity isomorphism $`[[SB,SB𝒦]][[S^2B,B𝒦]]`$ and show that the image of that element is not changed under the map we are trying to prove is always the identity. This is what we do. Under the isomorphism $`[[SB,SB𝒦]][[S^2B,B𝒦]]`$, coming from Bott-periodicity, the image of $`x`$ is represented by the asymptotic homomorphism $`S^2BB𝒦`$ arising by applying the Connes-Higson construction to the Toeplitz extension tensored with $`B`$ : (4) In other words, if $`\phi :SBQ(B𝒦)`$ is the Busby invariant of (4) the image of $`x`$ in $`[[S^2B,B𝒦]]`$ is $`[CH(\phi )]`$. For each separable $`G`$-subalgebra $`DQ(B)_G`$ we let $`\iota _D:DQ(B)_G`$ denote the inclusion. Then the boundary map $`:[[S^2B,SQ(B)_G𝒦]][[S^2B,B𝒦]]`$ is given by $$(z)=\underset{D}{lim}[CH(\iota _D)\mathrm{id}_𝒦]z,$$ where $``$ denote the composition product in $`E`$-theory. Hence $`^1[CH(\phi )]`$ is the element $`z[[S^2B,SQ(B)_G𝒦]]`$ with the property that $$\underset{D}{lim}[CH(\iota _D)\mathrm{id}_𝒦]z=[CH(\phi )]$$ for all large enough $`D`$. Let $`\iota :Q(B)_G𝒦Q(B𝒦)_G`$ be the natural embedding. By the naturality of the Connes-Higson construction, $$[CH(\iota _D)\mathrm{id}_𝒦]S([s\gamma \phi ])=[CH(\iota s\gamma \phi )]$$ for all separable $`G`$-subalgebras $`DQ(B)_G`$ which contains $`\gamma \phi (SB)`$. Since $`s\gamma _0`$ is equivariantly homotopic to the identity map, we have that $$[CH(\iota s\gamma \phi )]=(s\gamma _0)_{}[CH(\phi )]=[CH(\phi )],$$ so we conclude that $`^1[CH(\phi )]=S([s\gamma \phi ])`$. Hence $`\alpha S^1^1[CH(\phi )]=[\iota s\gamma \phi ]`$ by Lemma 4.3. Thus the image of $`[CH(\phi )]`$ in $`[[S^2B,B𝒦]]`$ under the composite map is $`CH[\iota s\gamma \phi ]=(s\gamma _0)_{}[CH(\phi )]=[CH(\phi )]`$. The proof is complete. ∎ ###### Lemma 5.2. Let $`\lambda \mathrm{Ext}(SA,B𝒦)`$. Then $`\phi =s\gamma \lambda `$ is an equivariant $``$-homomorphism $`\phi :SAQ(B)_G𝒦`$ such that $`\alpha [\phi ]=s_{}\gamma _{}[\lambda ]`$ in $`\mathrm{Ext}(SA,B𝒦)_h`$ and such that $`[\phi ]=0`$ in $`[[SA,Q(B)_G𝒦]]`$ implies that $`[\lambda ]=0`$ in $`\mathrm{Ext}(SA,B𝒦)`$. ###### Proof. If $`[\phi ]=0`$ in $`[[SA,Q(B)_G𝒦]]`$, there is a path $`\mu ^t,t[0,1]`$, of asymptotic homomorphisms $`SAQ(B)_G𝒦`$ such that $`\mu ^0=\phi `$ and $`\mu ^1=0`$ and a unit sequence $`𝒰=\{u_n\}`$ in $`B𝒦`$ such that $$q_{B𝒦}\stackrel{~}{\mu ^t},t[0,1],$$ (5) connects $`q_{B𝒦}\stackrel{~}{\phi }`$ to $`0`$. By Theorem 3.4 we may assume that $`\mu `$ is an equi-homotopy and it is then easy to see that (5) is a strong homotopy. By Lemma 2.4 we conclude from this that $`[q_{B𝒦}\stackrel{~}{\phi }]=0`$ in $`\mathrm{Ext}(SA,B𝒦)`$. But $`[q_{B𝒦}\stackrel{~}{\phi }]=[\phi ]`$ in $`\mathrm{Ext}(SA,B𝒦)`$ by Lemma 4.3. Hence $`\alpha [\phi ]=s_{}\gamma _{}[\lambda ]`$ in $`\mathrm{Ext}(SA,B𝒦)_h`$ and $`[\phi ]=0s_{}\gamma _{}[\lambda ]=0`$ in $`\mathrm{Ext}(SA,B𝒦)`$. To complete the proof it suffices to show that $`s_{}\gamma _{}:\mathrm{Ext}(SA,B𝒦)\mathrm{Ext}(SA,B𝒦)`$ is injective. However, $`\gamma `$ is an equivariant $``$-isomorphism and therefore $`\gamma _{}`$ is an isomorphism. The injectivity of $`s_{}:\mathrm{Ext}(SA,B)\mathrm{Ext}(SA,B𝒦)`$ follows from the weak stability of $`B`$ : There is a $`G`$-invariant isometry $`VM(B𝒦)`$ such that $`xV^{}s(x)V`$ is an equivariant $``$-automorphism $`B𝒦B𝒦`$ and $`s(x)=\mathrm{Ad}V(V^{}s(x)V)`$. Since $`\mathrm{Ad}V`$ induces the identity map on $`\mathrm{Ext}(SA,B𝒦)`$ we see that $`s_{}:\mathrm{Ext}(SA,B)\mathrm{Ext}(SA,B𝒦)`$ is an isomorphism. ###### Lemma 5.3. The map $`CH:\mathrm{Ext}(SA,B)[[S^2A,B]]`$ is injective. ###### Proof. Consider an extension $`\lambda \mathrm{Ext}(SA,B𝒦)`$ and assume that $`[CH(\lambda )]=0`$ in $`[[S^2A,B𝒦]]`$. With the notation from Lemma 5.2 we find that $`CH\alpha [\phi ]=CH[s\gamma \lambda ]=s_{}\gamma _{}[CH(\lambda )]=0`$. But then Lemma 5.1 implies that $`[\phi ]=0`$ in $`[[SA,Q(B)_G𝒦]]`$. By Lemma 5.2 this yields the conclusion that $`[\lambda ]=0`$ in $`\mathrm{Ext}(SA,B𝒦)`$. Thus $`CH:\mathrm{Ext}(SA,B𝒦)[[S^2A,B𝒦]]`$ is injective. But $`B`$ is weakly stable so the result follows. ∎ ###### Corollary 5.4. $`\mathrm{Ext}(SA,B)=\mathrm{Ext}(SA,B)_h`$. The surjectivity of $`CH:\mathrm{Ext}(SA,B)[[S^2A,B]]`$ follows from Lemma 5.1. Furthermore, it follows from Lemma 5.3 that $`\alpha `$ is well-defined as a map $`\alpha :[[SA,Q(B)_G𝒦]]\mathrm{Ext}(SA,B𝒦)`$ and then Lemma 5.1 tells us that $$CH^1=\alpha S^1^1.$$ Another description of $`CH^1`$ can be obtained from \[MT2\]. The crucial construction for this is the map $`E`$ which was considered in \[MT1\] and \[MT2\], inspired by \[MM\] and \[MN\]. However, in \[MT1\] and \[MT2\] we only defined $`E`$ as a map into homotopy classes of extensions, so to see that the $`E`$-construction can also invert the CH-map of Lemma 5.3 we must show that it is well-defined as a map from homotopy classes of asymptotic homomorphisms to stable unitary equivalence classes of extensions. Let us therefore review the construction. Given an equicontinuous asymptotic homomorphism $`\phi =\{\phi _t\}_{t[1,\mathrm{})}:AB`$ we choose a discretization $`\{\phi _{t_i}\}_i`$ such that $`lim_i\mathrm{}t_i=\mathrm{}`$ and $`lim_i\mathrm{}sup_{t[t_i,t_{i+1}]}\phi _t(a)\phi _{t_i}(a)=0`$ for all $`aA`$. Since $`G`$ is $`\sigma `$-compact (and $`\phi `$ equicontinuous) we can also arrange that $$\underset{i\mathrm{}}{lim}\underset{t[t_i,t_{i+1}]}{sup}\underset{gK}{sup}g\phi _t(a)\phi _t(ga)=0$$ for all $`aA`$ and all compact subsets $`KG`$. To define from such a discretization a map $`\mathbb{\Phi }:A𝕃(l_2()B)`$ we introduce the standard matrix units $`e_{ij},i,j`$, which act on the Hilbert $`B`$-module $`l_2()B`$ in the obvious way. Then $$\mathbb{\Phi }(a)=\underset{i1}{}\phi _{t_i}(a)e_{ii}$$ defines a map $`\mathbb{\Phi }:A𝕃(l_2()B)`$. As in the proof of Lemma 2.1 we can define a representation of $`G`$ on $`l_2()B`$ and in this way obtain a representation of $`G`$ as automorphisms of $`𝕃(l_2()B)`$. Since $`B`$ is weakly stable we can identify $`B`$ with $`𝕂(l_2()B))`$, the $`B`$-compact operators in $`𝕃(l_2()B)`$. Observe that $`\mathbb{\Phi }`$ is then an equivariant $``$-homomorphism modulo $`B`$. Furthermore, $`\mathbb{\Phi }(a)`$ commutes modulo $`B`$ with the two-sided shift $`T=_je_{j,j+1}`$ which is $`G`$-invariant. So we get in this way a $`G`$-extension $`E(\phi ):AQ(B)=𝕃(l_2()B)/𝕂(l_2()B)`$ such that $$E(\phi )(fa)=f(\underset{¯}{T})\underset{¯}{\mathbb{\Phi }(a)}$$ for all $`fC(𝕋),aA`$. Here and in the following we denote by $`\underset{¯}{S}`$ the image in $`Q(B)=𝕃(l_2()B)/𝕂(l_2()B)`$ of an element $`S𝕃(l_2()B)`$. ###### Lemma 5.5. $`E(\phi )`$ is a semi-invertible $`G`$-extension, and up to stable unitary equivalence it does not depend on the chosen discretization of $`\phi `$. ###### Proof. Consider another discretization $`(\phi _{s_i})_i`$ of $`\phi `$ and define $`\mathbb{\Psi }:A𝕃(l_2()B)`$ by $$\mathbb{\Psi }(a)=\underset{i0}{}\phi _{s_{i+1}}(a)e_{ii}.$$ There is then a $`G`$-extension $`E(\phi ):C(𝕋)A𝕃(l_2()B)/𝕂(l_2()B)`$ such that $`E(\phi )(fa)=f(\underset{¯}{T})\underset{¯}{\mathbb{\Psi }(a)}`$. It suffices to show that $`E(\phi )E(\phi )`$ is unitarily equivalent to an asymptotically split $`G`$-extension. Define $`\mathrm{\Lambda }:A𝕃(l_2()B)`$ such that $$\mathrm{\Lambda }(a)=\underset{i1}{}\phi _{t_i}(a)e_{ii}+\underset{i0}{}\phi _{s_{i+1}}(a)e_{ii}.$$ There is then a $`G`$-extension $`\pi _0:C(𝕋)A𝕃(l_2()B)/𝕂(l_2()B)`$ such that $`\pi _0(fa)=f(\underset{¯}{T})\underset{¯}{\mathrm{\Lambda }(a)}`$. $`E(\phi )E(\phi )`$ is clearly unitarily equivalent (via a $`G`$-invariant unitary) to $`\pi _00`$, so it suffices to show that $`\pi _0`$ is asymptotically split. For each $`n`$ we define $`\mathrm{\Lambda }_n:A𝕃(l_2()B)`$ by $$\begin{array}{cc}& \mathrm{\Lambda }_n(a)=\hfill \\ & \underset{i>n}{}\phi _{t_i}(a)e_{ii}+\underset{1in}{}\phi _{t_n}(a)e_{ii}+\underset{\{i0:s_{i+1}t_n\}}{}\phi _{t_n}(a)e_{ii}+\underset{\{i0:s_{i+1}>t_n\}}{}\phi _{s_i}(a)e_{ii}.\hfill \end{array}$$ Then $`\{\mathrm{\Lambda }_n\}_n`$ is a discrete asymptotic homomorphism such that $`lim_n\mathrm{}\mathrm{\Lambda }_n(a)\mathrm{\Lambda }_{n+1}(a)=0`$, $`lim_n\mathrm{}g\mathrm{\Lambda }_n(a)\mathrm{\Lambda }_n(ga)=0,gG`$, $`lim_n\mathrm{}T\mathrm{\Lambda }_n(a)\mathrm{\Lambda }_n(a)T=0`$ and $`\mathrm{\Lambda }_n(a)=\mathrm{\Lambda }(a)`$ modulo $`𝕂(l_2()B)`$. By convex interpolation and an obvious application of the $`C^{}`$-algebra $$\{fC_b([1,\mathrm{}),M(B)):q_B(f(t))=q_B(f(1)),t[1,\mathrm{})\}/C_0([1,\mathrm{}),B)$$ we get an asymptotic homomorphism $`(\pi _t)_{t[1,\mathrm{})}:C(𝕋)AM(B)=𝕃(l_2()B)`$ such that $`\pi _0=q_B\pi _t`$ for all $`t`$. ∎ Theorem 3.4 and Lemma 5.5 in combination show that there is group homomorphism $`E:[[SA,B]]\mathrm{Ext}(C(𝕋)SA,B)`$ such that $`E[\phi ]=[E(\phi )]`$ for any equicontinuous asymptotic homomorphism $`\phi :SAB`$. By pulling extensions back along the inclusion $`S^2AC(𝕋)SA`$ we can also consider $`E`$ as a map $`E:[[SA,B]]\mathrm{Ext}(S^2A,B)`$. Let $`\chi :SAS^3M_2(A)`$ be a $``$-homomorphism which is invertible in KK-theory. By weak stability of $`B`$ there is also an isomorphism $`\beta :[[S^2A,B]][[S^2M_2(A),B]]`$. Let $`\xi :S^2𝒦`$ be the asymptotic homomorphism which arises from the Connes-Higson construction applied to the Toeplitz extension. By changing $`\chi `$ ’by a sign’ we may assume that the composite map is the identity. Consider the diagram The square commutes by the naturality of the Connes-Higson construction, and it follows from Lemma 2.3 of \[MT2\] (or Lemma 5.5 of \[MT1\]) that $`(S\chi )^{}CHE\beta =\mathrm{id}`$. We conclude therefore that $`CH\chi ^{}E\beta =\mathrm{id}`$. We have now obtained our main results : ###### Theorem 5.6. Let $`A`$ and $`B`$ be separable $`G`$-algebras, $`B`$ weakly stable. $`CH:\mathrm{Ext}(SA,B)[[S^2A,B]]`$ is an isomorphism with inverse $`\chi ^{}E\beta `$. It follows, of course, that the bifunctor $`\mathrm{Ext}(SA,B)`$ has the same properties as $`E`$-theory, such as excision and Bott periodicity in both variables, for example. ###### Theorem 5.7. Let $`A`$ and $`B`$ be separable $`G`$-algebras, $`B`$ weakly stable, and let $`\phi ,\psi :SAQ(B)`$ be two $`G`$-extensions. The following conditions are equivalent : 1. $`[\phi ]=[\psi ]`$ in $`\mathrm{Ext}(SA,B)`$ (i.e. $`\phi `$ and $`\psi `$ are stably unitarily equivalent). 2. $`\phi 0`$ and $`\psi 0`$ are strongly homotopic. 3. $`\phi `$ and $`\psi `$ are homotopic. ###### Proof. 1) $``$ 2): Assuming 1) there is an asymptotically split extension $`\lambda `$ such that $`\phi \lambda `$ and $`\psi \lambda `$ are unitarily equivalent. By Lemma 6.1 of \[Th1\] this implies that $`\phi \lambda 0`$ and $`\psi \lambda 0`$ are strongly homotopic. Then $`\phi \lambda (\lambda \alpha )0`$ and $`\psi \lambda (\lambda \alpha )0`$ are also strongly homotopic, where $`\alpha \mathrm{Aut}SA`$ inverts the orientation of the suspension. 2) follows by observing that $`\lambda (\lambda \alpha )`$ is strongly homotopic to $`0`$. 2) $``$ 3) follows because an invariant isometry in $`M(B)`$ can be connected to $`1`$ via a strictly continuous path of $`G`$-invariant isometries, cf. e.g. Lemma 3.3 2) of \[Th1\]. 3) $``$ 1) follows from Lemma 5.3. ∎ ###### Remark 5.8. It is easy to extend Theorem 5.6 and Theorem 5.7 to the case where $`B`$ is only $`\sigma `$-unital (i.e. contains a strictly positive element). In fact, it suffices to observe that $$\mathrm{Ext}(SA,B)\underset{}{\mathrm{lim}}_D\mathrm{Ext}(SA,D),$$ where we take the limit over all weakly stable separable $`G`$-subalgebras $`D`$ of $`B`$ with the property that $`D`$ contains a positive element which is strictly positive in $`B`$. ## 6. $`K`$-homology It follows from Theorem 5.6 and Theorem 5.7 that $`\mathrm{Ext}(SA,B)=[[S^2A,B]]`$ can also be identified with the homotopy classes of equivariant $``$-homomorphisms $`\psi :SAQ(B)`$ with the property that $`\psi 0`$ is strongly homotopic to $`\psi `$. As a consequence we conclude that $$[[S^2A,B]]\underset{}{\mathrm{lim}}_n[SA,Q(B)M_n()],$$ where $`[,]`$ denotes homotopy classes of equivariant $``$-homomorphisms. In the important special case where $`B=𝒦`$, and the group $`G`$ is trivial, we can even do better. Let $`Q`$ denote the Calkin algebra, $`Q=𝕃(l_2)/𝕂(l_2)`$. ###### Lemma 6.1. Let $`\phi :SAQ`$ be a $``$-homomorphism. There is then an isometry $`V𝕃(l_2)`$ with infinite dimensional co-kernel and a $``$-homomorphism $`\phi _0:SAQ`$ such that $`\phi `$ is homotopic to $`\mathrm{Ad}q_𝒦(V)\phi _0`$. ###### Proof. We may assume that $`\phi `$ is not homotopic to $`0`$. Let $`\iota _i:SASA,i=1,2`$, be $``$-homomorphisms with orthogonal ranges, both homotopic to the identity map. Then $`\phi \iota _1`$ and $`\phi \iota _2`$ are homotopic to $`\phi `$, and in particular non-zero. Let $`a`$ be a non-zero positive element in the range of $`\phi \iota _2`$ and let $`b𝕃(l_2)`$ be a positive lift of $`a`$. By spectral theory $`b𝕃(l_2)b`$ contains a projection $`E`$ with non-zero image in $`Q`$. Since $`(1q_𝒦(E))x=x`$ for all $`x\phi \iota _1(SA)`$, we conclude that $`1q_𝒦(E)`$ is non-zero in $`Q`$. It follows that there is an isometry $`V`$ with infinite dimensional co-kernel such that $`VV^{}=1E`$. Set $`\phi _0=\mathrm{Ad}q_𝒦(V^{})\phi \iota _1`$. ###### Theorem 6.2. Let $`A`$ be a separable $`C^{}`$-algebra. Then $`E(A,)`$ is naturally isomorphic to the group $`[SA,Q]`$ of homotopy classes of $``$-homomorphisms from $`SA`$ to $`Q`$. ###### Proof. It follows from Lemma 6.1 that $`\phi `$ is strongly homotopic to $`\phi 0`$ for any $``$-homomorphism $`\phi :SAQ`$. By using that the unitary group of $`𝕃(l_2)`$ is norm-connected, it follows from this and Theorem 5.7 that $`\mathrm{Ext}(SA,𝒦)`$ is naturally isomorphic to $`[SA,Q]`$. Since $`\mathrm{Ext}(SA,𝒦)`$ is naturally isomorphic to $`E(A,)`$ by Theorem 5.6, this completes the proof. ∎ A weak version of Theorem 6.2 was conjectured by Rosenberg in \[R\]. V. M. Manuilov Dept. of Mech. and Math., Moscow State University, Moscow, 119899, Russia e-mail: manuilov@mech.math.msu.su K. Thomsen Institut for matematiske fag, Ny Munkegade, 8000 Aarhus C, Denmark e-mail: matkt@imf.au.dk
warning/0004/cond-mat0004028.html	ar5iv	text	# A new approach to the study of the ground-state properties of 2D Ising spin glass ## 1 Introduction The equilibrium properties of spin glass have remained a great challenge in numerical simulations. Investigating the equilibrium ground-state structure of spin glass is also important and interesting. In the last 20 years, there has been a great deal of work on spin glass . It is generally agreed that the simplest spin glass system for most theoretical work is the Edwards-Anderson (EA) model, whose Hamiltonian is $$H=\underset{<i,j>}{}J_{ij}\sigma _i\sigma _j,$$ (1) where $`\sigma _i`$ takes on the values $`\pm 1`$ and the sum goes over the nearest neighbors. The $`J_{ij}`$ are dimensionless variables which describe the random interactions between the spins and are taken as $`J_{ij}=\pm 1`$. In two dimensions, a phase transition occurs only at zero temperature for this kind of $`\pm J`$ Ising spin glass with nearest neighbor interactions. This model has been studied previously by the transfer matrix method , replica Monte Carlo method , multicanonical ensemble method and many other methods (see ref for a review). The traditional Monte Carlo methods mostly concentrate on generating standard statistical ensembles, e.g., the canonical ensemble or microcanonical ensemble. Using the canonical ensemble simulations, we need to simulate at different temperatures to get full information about the system. It is tedious to calculate certain thermodynamic quantities like the free energy and the entropy since the density of states cannot be obtained directly from the simulation data. The correlation between subsequent configurations generated by canonical ensemble simulations also causes the ergodicity problem for some systems. In 1991, Berg proposed the multicanonical ensemble method to overcome the above shortcomings of simulations on canonical ensemble. The multicanonical ensemble is an ensemble where the probability $`P(E)`$ of having energy $`E`$ at equilibrium is a constant. The multicanonical method has been very successful in solving the systems that involve energy barriers. Recently, Wang proposed a dynamics which can generate a flat histogram in the energy space as the multicanonical method. This dynamics has some connections with the broad histogram method , which does not give the correct microcanonical average . Similar to the broad histogram method, the new dynamics is also based on $`N(\sigma ,\mathrm{\Delta }E)`$, the (microcanonical) average number of potential moves which increase the energy by $`\mathrm{\Delta }E`$ in a single spin flip. A cumulative average (over Monte Carlo steps) can be used as a first approximation to the exact microcanonical average in the flip rate. Thermodynamic quantities can be then calculated from the simulation data with ease. In this paper, we use the new method to study the thermodynamics as well as ground-state properties for the two-dimensional Ising spin glass system. In Section 2, the flat histogram transition matrix Monte Carlo dynamics is described. Using the flat histogram sampling, we get the average number of potential moves $`N(\sigma ,\mathrm{\Delta }E)_E`$, which can be used to construct a transition matrix Monte Carlo dynamics in the energy space . We apply the new method to two-dimensional Ising spin glass and present some numerical results in Section 3. In the last section, we give a conclusion to the new method. ## 2 The transition matrix Monte Carlo dynamics with the flat histogram sampling To connect our dynamics with single-spin-flip Glauber dynamics , we restrict the protocol of each move to be single-spin flip in the following discussion. For a given state $`\sigma `$ with energy $`E`$, consider all possible single-spin flips. The single-spin flips change the current state into $`N`$ possible new states, with new energy $`E^{}=E+\mathrm{\Delta }E`$. For two-dimensional Ising spin glass, $`\mathrm{\Delta }E`$ = $`0`$, $`\pm 4`$, and $`\pm 8`$. We classify the $`N`$ new states according to $`\mathrm{\Delta }E`$ and count the number of $`N(\sigma ,\mathrm{\Delta }E)`$. Since each move from the state $`\sigma `$ of energy $`E`$ to the state $`\sigma ^{}`$ of energy $`E^{}`$ and the reverse move are both allowed, the total number of moves from all the states with energy $`E`$ to $`E^{}`$ is the same as from $`E^{}`$ to $`E`$. Thus, we have $$\underset{E(\sigma )=E}{}N(\sigma ,\mathrm{\Delta }E)=\underset{E(\sigma ^{})=E+\mathrm{\Delta }E}{}N(\sigma ^{},\mathrm{\Delta }E).$$ (2) The microcanonical average of a quantity $`A(\sigma )`$ is defined as $$A_E=\frac{1}{n(E)}\underset{E(\sigma )=E}{}A(\sigma ),$$ (3) where the summation is over all the configurations having energy $`E`$ and $`n(E)`$ is the density of states. In terms of the microcanonical averages, we can rewrite Eq. (2) as $$n(E)N(\sigma ,\mathrm{\Delta }E)_E=n(E+\mathrm{\Delta }E)N(\sigma ^{},\mathrm{\Delta }E)_{E+\mathrm{\Delta }E}.$$ (4) Eq. (4) is the basic result of the broad histogram method . While the broad histogram random walk algorithm is not correct, Eq. (4) is not problematic and taken as the starting point of the flat histogram sampling. We select a site to flip at random. The flip rate for a single-spin flip from state $`\sigma `$ with energy $`E`$ to $`\sigma ^{}`$ with energy $`E^{}=E+\mathrm{\Delta }E`$ is chosen as $$r(E^{}\|E)=\mathrm{min}(1,\frac{N(\sigma ^{},\mathrm{\Delta }E)_E^{}}{N(\sigma ,\mathrm{\Delta }E)_E}).$$ (5) Then the detailed balance condition for this rate $$r(E^{}\|E)P(\sigma )=r(E\|E^{})P(\sigma ^{})$$ (6) is satisfied for $`P(\sigma )1/n(E(\sigma ))`$. Thus the energy histogram is flat , $$P(E)=\underset{E(\sigma )=E}{}P(\sigma )n(E)\frac{1}{n(E)}=\text{const}.$$ (7) Since $`N(\sigma ,\mathrm{\Delta }E)_E`$ is not known in general, an approximation scheme should be used to start the simulation. For those $`E`$ which we have not visited yet, we simply set $`r(E^{}\|E)=1`$. Then a cumulative average (over Monte Carlo steps) can be used as an approximation to the exact microcanonical average in the flip rate. We have numerical evidence that this procedure converges to the exact result. We can then construct a transition matrix Monte Carlo dynamics in the energy space with $`N(\sigma ,\mathrm{\Delta }E)_E`$. For a single-spin-flip Glauber dynamics with energy change $`\mathrm{\Delta }E`$, the flip rate is given as $$w(\mathrm{\Delta }E)=\frac{1}{2}\left[1\mathrm{tanh}\left(\frac{\mathrm{\Delta }E}{2k_BT}\right)\right].$$ (8) Since there are (on average) $`N(\sigma ,\mathrm{\Delta }E)_E`$ different ways of going from $`E`$ to $`E^{}=E+\mathrm{\Delta }E`$, the total probability for transition from $`E`$ to $`E^{}`$ is $$W(E+\mathrm{\Delta }E\|E)=w(\mathrm{\Delta }E)N(\sigma ,\mathrm{\Delta }E)_E,for\mathrm{\Delta }E0.$$ (9) The diagonal elements can be determined by $`_{\mathrm{\Delta }E}W(E+\mathrm{\Delta }E\|E)=1`$, since the total probability from $`E`$ to $`E^{}`$ is 1. This new dynamics in the space of energy $`E`$ is related to single-spin-flip dynamics by $$W(E^{}\|E)=\frac{1}{n(E)}\underset{E(\sigma )=E}{}\underset{E(\sigma ^{})=E^{}}{}\mathrm{\Gamma }(\sigma ^{}\|\sigma ).$$ (10) where $`\mathrm{\Gamma }(\sigma ^{}\|\sigma )`$ is the transition matrix of the single-spin-flip dynamics. The equilibrium state of the transition matrix gives the canonical probability distribution of energy $`P_T(E)n(E)\mathrm{exp}(E/k_BT)`$. An important aspect of this dynamics is that we can calculate the thermodynamic quantities easily by just performing one simulation for each coupling state $`J_{ij}`$. The density of states $`n(E)`$ can be obtained through Eq. (4). Once we have the density of states $`n(E)`$, we can obtain $`P_T(E)`$ and then calculate any thermodynamic quantities of interest. In actual implementation, we usually determine $`P_T(E)`$ directly from the detailed balance equation $$W(E+\mathrm{\Delta }E\|E)P_T(E)=W(E\|E+\mathrm{\Delta }E)P_T(E+\mathrm{\Delta }E)$$ (11) instead of solving Eq. (4). From Eq. (9), we know, the transition matrix $`W(E^{}\|E)`$ can be formed at any temperature once the quantity $`N(\sigma ,\mathrm{\Delta }E)_E`$ is computed accurately. In other words, the Monte Carlo computation is uncorrelated to thermodynamics. The temperature dependence enters only after simulation in the weighting formula. Like Berg’s multicanonical ensemble simulations, our dynamics also generate a multicanonical ensemble in the energy space. From this point, both of the two dynamics have the same goal of flattening the space of energy. But they are quite different in implementation. In the multicanonical ensemble method, the flip rate is chosen as the inverse of the density of states $`n(E)`$, parametrized in some way. To start the simulation, we need give an estimate of $`n(E)`$, since $`n(E)`$ is not initially known. Thus, the efficiency of this method is determined by the goodness of the estimated $`n(E)`$. If $`n(E)`$ is not given properly, say far off the true density, the simulations may get stuck in some region. With our method, we sample the energy space with a flip rate which is related to the density of states through Eq. (4). The central quantity is $`N(\sigma ,\mathrm{\Delta }E)_E`$ which can be quite accurate in a short simulation time. And the accuracy of this quantity is improved by further simulations. We then provide an alternate for the estimate of $`n(E)`$, which leads to a more efficient way for simulating the multicanonical ensemble. The flat histogram also generalizes easily to multi-variate models . An example is the Ising Spin Glass model with overlap parameter $`q`$ which has the Hamiltonian $$H_2=\underset{<i,j>}{}J_{ij}\sigma _i^1\sigma _j^1\underset{<i,j>}{}J_{ij}\sigma _i^2\sigma _j^2h\underset{q}{\underset{}{\underset{i}{}\sigma _i^1\sigma _i^2}}.$$ (12) $`E`$ refers to the first two interaction terms involving the coupling constants $`J_{ij}`$. In this bivariate case, the quantity $`N(\sigma ,\mathrm{\Delta }E)`$ generalize to $`N(\sigma ,\mathrm{\Delta }E,\mathrm{\Delta }q)`$. It can be easily shown that the detailed balance condition is now $$n(E,q)N(\sigma ,\mathrm{\Delta }E,\mathrm{\Delta }q)_{E,q}=n(E+\mathrm{\Delta }E,q+\mathrm{\Delta }q)N(\sigma ^{},\mathrm{\Delta }E,\mathrm{\Delta }q)_{E+\mathrm{\Delta }E,q+\mathrm{\Delta }q}$$ (13) with $`n(E,q)`$ as the new “density of states”. The algorithm gives a flat histogram in both $`E`$ and $`q`$. ## 3 Numerical results We have performed simulations on lattices of size $`L=4,10,16,24`$ and $`32`$. Each simulation starts with independent random numbers. To illustrate the performance of our algorithm, we define the time $`\tau _L`$ as the average number (over coupling constant $`J_{ij}`$) of Monte Carlo steps needed to reach the ground-states. A Monte Carlo step is defined as flipping each spin on the lattice once (on the average). Table 1 gives an overview of typical time in Monte Carlo steps to reach the ground-states, starting from an arbitrary energy level. The time to reach the ground-states depends on the size of the system and also the random interactions. We consider a large number of random coupling states to make the statistical error small enough in Table 1. The simulations are long enough to ensure that the ground states are really reached. In Fig.1 we plot the time $`\tau _L`$ versus lattice size $`L`$ on a double log scale. The data are consistent with a straight-line fit, which gives the finite-size behavior $$\tau _LL^{4.71},\text{MC steps}.$$ (14) The corresponding CPU time for a Digital Alpha 600M workstation is also shown in Table 1. For accuracy, 5 independent runs are performed for each lattice size to obtain the average CPU time. Up to $`L=32`$ the CPU time can be approximated by a polynomial function of $`L^{6.08}`$. We also consider the tunneling time which is defined as the average Monte Carlo steps needed to move from $`E_{max}`$ to $`E_{min}`$, or from $`E_{min}`$ to $`E_{max}`$. Note that $`E_{min}`$ is the same as ground state energy and $`E_{max}`$ is $`NE_{min}`$. We note that Berg’s definition about tunneling time is slightly different from ours. During the simulation, Berg imposed a constraint $`_{ij}J_{ij}=0`$. But for our method, both the time $`\tau _L`$ and the tunneling time will not be affected significantly by the imposition of the constraint. We start the simulations from an arbitrary energy level. Table 2 gives an overview of the tunneling time obtained using the two methods. The power law fits are $$\tau _{M.C.}L^{4.43},\text{and}\tau _{F.H.}L^{5.03},$$ (15) for Berg’s method and current flat histogram method, respectively. It shows that they basically give the same tunneling time. We also compared with Hatano’s result that autocorrelation time scales approximately as volume $`N`$ of the system. We look at the tunneling time which is a better measure of the algorithm’s efficiency in our case. From our results given in Table 2, we found no support for Hatano’s result. Instead, the power law fit (see Fig. 2) $$\tau _qL^{4.45}$$ (16) is almost the same as the monovariate case. This is not surprising as Berg mentioned that the optimal performance for multicanonical algorithm is $`N`$(=$`L^2`$) based on random walk picture. The ground-state energy and entropy of the infinite system are also estimated using our method. It is straightforward to obtain the ground-state energy in the simulation stage. We calculate the ground-state entropy from $$S(E)=\frac{k_B}{N}\mathrm{ln}n(E).$$ (17) Since $`n(E)`$ can be calculated from the simulation data directly, we then obtain $`S(E)`$ with ease. To compare with the results obtained in the literature, we fit our data using the form $`f_L=f_{\mathrm{}}+c/L^2`$ and get $`e^0=1.4007\pm 0.0085`$, $`s^0=0.0709\pm 0.006`$. The energy fit is plotted in Fig. 3, and the entropy fit in Fig. 4. Our energy estimate $`e^0=1.4007\pm 0.0085`$ is consistent with the previous MC estimate $`e^0=1.407\pm 0.008`$ as well as with the transfer matrix result $`e^0=1.4024\pm 0.0012`$. Our entropy estimate $`s^0=0.0709\pm 0.006`$ is also consistent with the MC estimate $`s^0=0.071\pm 0.007`$ as well as the transfer matrix result $`s^0=0.0701\pm 0.005`$. For the two-dimensional Ising spin glass system, De Simone et al. use an exact algorithm based on the branch-and-cut technique to find the exact ground-states with system size up to $`50\times 50`$. They obtain the extrapolated result $`e^0=1.4022\pm 0.0003`$. When compared with Berg’s result, $`e^0=1.394\pm 0.007,s^0=0.081\pm 0.004`$, it seems that our method gives a more accurate estimate for ground-state energy and entropy for an infinite system. ## 4 Conclusions We have used a new approach to investigate the ground-state properties of the two-dimensional Ising spin glass. Compared with standard simulations, the advantage of our method is obvious. For the ergodicity problem encountered in standard simulations, our method behaves as well as Berg’s multicanonical ensemble method, while it is easier to be implemented compared with Berg’s method. Our method also generalize straightforwardly to multi-variate models without much effort in programming and theory. To find a true ground-state, we roughly need a CPU time of order $`L^6`$. It is the same with Lawler’s exact algorithm . Up to size $`50\times 50`$, De Simone’s algorithm also needs a time of order $`L^6`$. But it is not clear whether his algorithm can be efficiently implemented for 3D systems. However our method can also be easily applied to 3D spin glass system. If one is just interested in finding the ground-states, there are also other optimized algorithms. Chen’s learning algorithm is fast in finding the ground-states compared with most algorithms, but it is not a general one. Thermodynamic quantities cannot be obtained with this algorithm. We believe that the approach we present in this paper is useful in studying the thermodynamics as well as ground-state properties for spin glass systems. It also can be applied to other models because of its generality.
warning/0004/math-ph0004012.html	ar5iv	text	# References Sergey P. Novikov<sup>1</sup><sup>1</sup>1Math Department and IPST, University of Maryland, College Park, MD 20742-2431, USA and Landau Institute for Theoretical Physics, Moscow 117940, Kosygina 2; phone in USA 301-4054836(o), fax in USA 301-3149363, e-mail novikov@ipst.umd.edu; this work is supported by the NSF Grant DMS9704613 1.Classical and Modern Topology. 2.Topological Phenomena in Real World Physics According to the opinion of the Ancient Greeks, the famous real and mythical founders of Mathematics and Natural Philosophy like Pythagoras, Aristotle and others, in fact, borrowed them from the Egyptian and Middle East civilizations. However, what had been told before in the hidden mysteries Greek scientists transformed into written information acceptable for everybody. Exactly after that the development of science in the modern sense started and had already reached a very high level 2000 years ago. Therefore you may say that the free exchange of information and making it clear for people have been the most important discoveries of Greeks. I would say it is the basis of our science now. As you will see, any violation of this fundamental rule does serious harm to our science and inevitably leads to its decay. 1.Classical and Modern Topology. Prehistory. First fifty years of Topology. The first important topological ideas were observed by famous mathematicians and physicists like Euler, Gauss, Kelvin, Maxwell and their pupils, during the XVIIIth and XIXth Centuries. As everybody knows, it was Poincare’ who really started Topology as a branch of Mathematics in the late XIXth Century. Many top class mathematicians participated in the development of Topology in the first half of our century. A huge number of mutually connected fundamental notions were invented: degree of maps and singularities of vector fields, homotopy and homology groups, differential forms and smooth manifolds, the fundamental idea of transversality, the simplicial/cell(CW) and singular complexes as tools for studying topological invariants, braids, knot invariants and 3-manifolds, coverings, fibre bundles and characterisic classes and many others. Deep connections with Qualitative Analysis, Calculus of Variations, Complex Geometry and Dynamical Systems were established in this period. Combinatorial Group Theory and Homological Calculus started from topological sources. A great new field of topological objects unknown to the classical mathematics of the XIXth Century appeared finally in the 1940s. At that time this new area was known to few number of mathematicians only. However there was a very high density of really outstanding scientists among them. 1950s and 60s: Golden Age of Classical Topology. The fundamental set of algebraic ideas unifying all these branches of mathematics appeared in the 40s; a new era started about 1950. Spectral sequences of fibre bundles, sheafs, highly developed homological algebra of the groups, algebras and modules, Hopf algebras and coalgebras, were invented and heavily used for the calculation of topological invariants needed for the solution of the fundamental problems of topology. Let me point out that, in many cases, it was a completely new type of calculations based on the deep combination of the very general ”categorial properties” of these quantities with very concrete geometric, algebraic or analytical study of a completely new type. In the previous period, people even had no dreams as to how they could be calculated. Regular methods were built to calculate homotopy groups, for example. It was one of the most difficult problems of topology. A lot of them were computed completely or partially including the homotopy groups of spheres, Lie groups and homogeneous spaces. The topologically important cobordism rings were computed and used in many topological investigations. The famous signature formula for differentiable manifolds was discovered. It has an innumerous number of applications in the topology of manifolds. Besides that, this formula played a key role in the proof of the so-called Riemann-Roch theorem in Algebraic Geometry and later in the study index, the famous homotopy invariant of Fredholm operators. The mutual influence of Topology and Algebraic Geometry during that period led to the broad extension of the ideas of homology: the extraordinary (co)homology theories like K-theory and cobordisms appeared. They brought a new type of technic to topology with many applications. Representation theory and complex geometry of manifolds deeply unified with homological algebra and Hopf algebras. The technic of formal groups appeared here. It has been applied in particular for the improvement of the calculations of stable homotopy groups of spheres. As everybody knows, during this period topology solved the most fundamental problems in the theory of multidimensional smooth manifolds: Nontrivial differentiable structures on multidimensional spheres were discovered on the basis of the results of algebraic topology combined with a new understanding of the geometry of manifolds and bundles. The multidimensional analog of the Poincare Conjecture and H-cobordism theorem were proved. Counterexamples to the so-called ”Hauptvermunung der Topologie” were found. A classification theory for the multidimensional smooth (and for PL-manifolds as well) was completely constructed. The role of the fundamental group in this theory led to the development of a new branch of algebra: the algebraic K-theory. Topological invariance of the most fundamental characteristic classes was finally proved. The so-called ”Annulus Conjecture” was proved. No matter how elementary these results can be formulated, nobody has succeeded to avoid the use of a whole bunch of results and tools of algebraic and differential topology in the proof. The classification theory for the immersions of manifolds was constructed. The theory of multidimensional knots was constructed. Several classical problems of the theory of 3-manifolds also were solved during that period: the so-called Dehn’s program was finished after a 50 year break; the algorithm for recognizing the trivial knot in three-space has been theoretically constructed as a part of the deep understanding of the structure of 3-manifolds and the surfaces in them. As a by-product of topology, the fundamental breakthrough in the topological understanding of generic dynamical systems was reached. A new great period started in this area. Qualitative theory of foliations has been constructed with especially deep results for 3-manifolds. As a summary, I would like to add one more very important characteristic of the topological community in the golden age of classical topology: All important works have been carefully checked. If some theorem had not been proved, it immediately became known to everybody. So you can find a full set of proofs in the literature. Unfortunately, a full set of textbooks covering all these developments (1950-1970) has not been written yet. Many modern textbooks are written in a very absract way. Even if they cover some pieces formally, it is more difficult to read them than the original papers. Let me recommend to you the Encyclopedia article written exactly for the exposition of these ideas. 1970s: Period of decay. In my opinion, the period of the 1970s can be characterized as a period of decay for classical topology. There are many indications for that. Several leading scientists left topology for the new areas like algebra and number theory, riemannian and symplectic geometry, dynamical systems and complexity theory, functional analysis and representations, PDEs, and different branches of mathematical/theoretical physics…. It is certainly a good characterization of the community if it could generate such a flux of scientists in many different areas and bringing to them completely new ideas. Anyway, this community dispersed. What can we say about the topological community after that? First of all, some important new ideas appeared in the 70s (like localization technic in homotopy topology, the nicely organized theory of the rational homotopy type, hyperbolic topology of 3-manifolds). However, a huge informational mess was created in the 1970s. Let me point out that a series of fundamental results of that period was not written, with full proof, until now. Let me give you a list: Sullivan’s Haupvermutung theorem was announced first in early 1967. After the careful analysis made by Bill Browder and myself in Princeton of the first version in May 1967 (before publication), his theorem was corrected: a necessary restriction on the 2-torsion of the group $`H_3(M,Z)`$ was missing. This gap was found and restriction was added. Full proof of this theory has never been written and published. Indeed, nobody knows whether it has been finished or not. Who knows whether it is complete or not? This question is not clarified properly in the literature. Many pieces of this theory were developed by other topologists later (they used sometimes different ideas). Nobody has unified them until now. Indeed, these results were used by many others later. In particular, the final Kirby-Siebenmann classification of topological multidimensional manifolds therefore is not proved yet in the literature. The second story is the theory of Lipshitz structures on the manifolds. In the mid-seventies Sullivan distributed a preprint containing the idea how to prove existence and uniqueness of such structures on the manifold $`N^n,n4`$. This idea obviously included (for the uniqueness) the direct use of the Annulus Conjecture (and therefore of all ideas and technic needed in the proof of topological invariance of the rational Pontryagin Classes inside). Proof of the Lipshitz Theory has never been published. Indeed, many years later, already in the 1990s, some brilliant younger scientists developed a very nice theory of Fredholm (elliptic) operators on Lipshitz manifolds. As a corollary, they claimed that a new proof of topological invariance of rational Pontryagin classes has been obtained from Analysis (it was a problem posed by Singer in the 60s). Young scientists made a ”logical circle” believing in the classical results. Nobody told them that corresponding theorems have never been proved. How could it happen? This funny story shows the modern state of information in the topological community. Another informational mess has been created in 3D Hyperbolic topology. This beautiful area was started by Thurston in the mid-70s. For many years people could not find out what was proved here. In this area the situation has been finally resolved: it has been aknowledged that these methods lead to the proof of the original claim (the so-called Geometrization Conjecture) only for the special class of Haken manifolds. The Geometrization Conjecture means more or less that (in the case of closed 3-manifolds) the fundamental group can be realized as a discrete subgroup acting in the 3D Hyperbolic space if trivial necessary conditions are satisfied: all its abelian subgroups are cyclic and $`\pi _2=0`$. However, it is difficult to find out who actually proved this theorem? It seems for me that the younger mathematicians who managed to finish this program did not receive proper credit. I would like to mention that this kind of informational mess has happened since 1970 not only in topology. For example, the famous results of KAM in the three-body problem known since the early 60s were found recently unproved. It was announced for the first time at the Berlin Congress last year. In this case, some works supposedly containing full proof were published in the first half of the 60s. Does this mean that nobody accually read them for at least 30 years? Do you think that algebra is better? Let me tell you as a curious remark that all works of the Steklov Institute (i.e. Shafarevich’s) school in algebraic number theory, algebraic geometry and theory of finite $`p`$-groups awarded by the highest (Lenin and State) prizes in the former Soviet Union since 1959, did not contain full proof. The gaps in the proofs were found many years later. Not all these gaps were really deep. However, some of these authors knew their mistakes many years before they became publicly known and could not correct them. They managed to fulfill gaps after many years , using much later technical achievements made by other people. Does it mean that in the corresponding time, despite many public presentations, nobody in fact read these great works? Can we say that all proofs are known now in all these cases? There are much worst cases in modern algebra indeed. How many of you know that the so-called classification of simple finite groups did not exist as a mathematical theorem until now? In this case we can even say that in fact (as a few number of real experts have known since 1980) no one work existed claiming that this problem was finished in this work. All public opinion has been based only on the ”New York Times Theorem” for the past 20 years. 1980s and 90s: Period of recovery. The role of Quantum Field Theory. It became clear already in the late 70s that modern quantum field theory started to generate new ideas in topology. It gave several new alternative ways to construct topological invariants: Path integral for the metric-independent actions on manifolds was used for the first time. The famous self-duality equation appeared first in the works of physicists. It was applied in the 80s for the solution of fundamental topological problems in the theory of 4-manifolds. Quantum string theory brought in the early 80s new deep results in the theory of the classical Fuchsian groups and moduli spaces. At first physicists (like t’Hooft and Polyakov) were not interested very much in such by-products of their activity. They always said that they were doing physics of the real world, not pure mathematics. However the next wave of brilliant physicists (like Witten, Wafa and others) started to solve problems of pure mathematics. Such purely topological subjects like the Morse theory and cobordism theory associated with action of compact groups on manifolds, were developed in the 80s from the completely new point of view. Symplectic Topology reached a very high level in the late 80s. We are facing now impressive development of Contact Topology. Certainly Quantum Theory brought new beautiful ideas. Besides that, the fundamental new invariants of knots were discovered in the 80s by the topologists who came from functional analysis and theory of $`C^{}`$ algebras. These invariants also received quantum treatment in the late 80s. The beautiful connection of the specific Feinmann diagrams with surfaces was borrowed from physics literature. It became a very effective tool for the solution of several topological problems. Unfortunately, only a few number of mathematicians learned this technic and started to apply it in topology. I know only Singer, Konzevich and a very small number of others. Even if you will add here the names of pure mathematicians who learned this with the intention to do real physics, this list will increase inessentially. I do not count here people who were trained originally in the physics community. A large number of them moved into pure mathematics with the intention to prove rigorous theorems about the models serving (in their opinion) as an idealization of theoretical physics. They call this area Mathematical Physics, but not everybody agrees with such a definition of mathematical physics. This community does not do topology. I would like to make a remark here concerning a beautiful work of Konzevich calculating certain Chern numbers on the punctured moduli spaces of Riemann Surfaces through the special solution to the KdV hierarchy. This folmula has been known as a Witten Conjecture. You have to specify for this some compactification of the moduli spaces of punctured Riemann surfaces, otherwise it makes no sense. Konzevich accually proved this formula for one specific (”Strobel-Penner”) compactification in 1991. What about the standard Deligne-Mumford compactification? Konzevich claimed in 1992 in his work in Inventiones that it is true. However, no proof has been presented until now. So this problem is open. There was a mistakable statement about this at the Berlin Congress. Let me point out that the physics community did not create any informational mess in topology. According to their training tradition, theoretical work produces Conjectures which should be proved only by some kind of experiment. Starting to do beautiful nonrigorous mathematics, they do not claim that they ”proved” something. They are saying that they ”predicted this fact”. In the case of pure mathematics, the final proof done by pure mathematicians these people may treat as an ”experimental confirmation”. In the past ten years several deep results have been obtained in the 4D topology. We cannot say this about 3D topology: quantum invariants here created some sort of ”invariantology”: a lot of people are constructing topological invariants but no one new topological result has been obtained for almost 10 years. Indeed, these ideas look beautiful in some cases. In my opinion, new deep results will appear after better understanding of the relationship of new invariants with classical topology. Topological Phenomena in Real World Physics Topological ideas in physics in the period of the early 80s. I spent about 10 years learning different parts of Modern Theoretical Physics in the 60s and 70s. After joining the physics community (i.e. Landau school) in the early 70s I found out that most physicists did not know at all the new areas of mathematics like topology, dynamical systems and algebraic geometry, including analysis on Riemann surfaces. The quantum people knew some extracts from the group theory and representations because they needed it in Solid State Physics as well as in Elementary Particles Theory since the 1960s. A lot of them knew something about Riemannian Geometry because of the Einsteinian General Relativity. However, these people had already heard something about the new mathematics of the XXth century and badly wanted to find its realization in physics. You have to take into account that between them there was a great number of extremely talented people at that time with very good training in practical mathematics. In some cases I was able to help physicists (like Polyakov, Volovic and some others) to learn and to use topology in the 70s. I worked this period in General Relativity (Homogeneous Cosmological Models) and Periodic KdV Theory with my pupils and collaborators. We found completely nonstandard applications of Dynamical Systems and Algebraic Geometry in these areas. However, until the late 70s I did not produce any new topological ideas. My very first topological work in physics was made in 1980 (see ). I started to use in the spectral theory of the Schrodinger operators in periodic lattice and magnetic field the idea of transversality applied to the families of Hermitian matrices or elliptic operators on the torus. This idea led to the discovery of the series of topological invariants, Chern Numbers of Dispersion Relations. They are well-defined for the generic operators only. The classical Spectral Theory in mathematics never considered such quantities because they are not defined for every operator with prescribed analytical properties of coefficients. The ideology of transversality is important here. This work was not understood by my colleagues-physicists at that time (the vice-editor of JETP did not want to publish it as ”nonphysical”, so I published it in the math literature). People thought that the important integer-valued observable quantities in Solid State Physics may come from symmetry groups only. Indeed, the Integral Quantum Hall phenomenon was discovered soon. Some famous theoretical physicists rediscovered my mathematical idea after that. It is certainly a sum of the Chern classes of dispersion relations below the Fermi level. My next topological discovery was made in the joint work with student I.Schmelzer in 1981, dedicated to the very special problem of classical mechanics and hydrogynamics (see ). I immediately realized its value for modern theoretical physics, as well as for mathematics, and developed this idea in several directions in the same year (). The series of work in the Theory of Normal Metals which I am going to discuss today, is also one of by-products of that discovery. Doing the Hamiltonian factorization procedure for the top systems on the phase spaces like $`T^{}(SO_3)`$ by the action of $`S^1`$, you are coming to the systems mathematically equivalent to the motion of the charge particle on the 2-sphere. This sphere is equipped by some nontrivial Riemannian metric. What is important and has been missed by the good experts in analytical mechanics like Kozlov and Kharlamov is that the effective magnetic field like Dirac monopole appears here for the nonzero values of the ”area integral” associated with $`S^1`$-action. It means precisely that the magnetic flux along the sphere is nonzero. The reason for this is that the symplectic (Poisson) structure after factorization is topologically nontrivial. In terms of modern symplectic geometry, the magnetic field is equivalent to the correction of the symplectic structure. This fact is not widely known in the geometric community even now. The appearance of the topologically nontrivial symplectis structures after $`S^1`$-factorization of symplectic manifolds was independently discovered and formulated in geometric, nonphysical terminology in 1982 in the beautiful work for different goals (calculating of integrals). It has been realized in that the action functional for such systems is in fact a closed 1-form on the spaces of loops. These functionals have been immediately generalized for higher dimensions, to the spaces of mappings $`F`$ of $`q`$-manifolds in some target space $`M`$ where a closed $`q+1`$-form is given instead of magnetic field. We are coming finally to the action functional well-defined as a closed 1-form on the mapping spaces $`F`$. The topological quantization condition for such actions was formulated in 1981 as a condition that this closed 1-form should define an integral cohomology class in $`H^1(F,Z)`$. It is necessary and sufficient for the Feinmann amplitude to be well-defined as a circle-map $$\mathrm{exp}\{iS/h\}:FS^1$$ For the case $`q=1`$ the original Dirac requirement was based on a different idea: the magnetic field should be a Chern class for the line bundle whose space of sections should serve as a Hilbert space of states for our Quantum Mechanics. Therefore it should be integral in $`H^2(M,Z`$). In pure topology and in the Calculus of Variations these ideas led to the construction of the Morse-type theory for the closed 1-forms on the finite- and infinite-dimensional manifolds. Let me refer to the last publication of the present author (with P.Grinevich) in this direction where the survey of results and problems is discussed. I would like to point out that for the compact symplectic manifolds the action functional for any nontrivial Hamiltonian system is multivalued. The cohomological class of symplectic form cannot be trivial here. I do not know of such cases in real physics where the symplectic manifold is compact. However, even in the community of symplectic geometers nobody paid attention to such properties of action functional until the 90s. After that I started to think about different aspects of the Hamiltonian Theory where the class of one-valued functions naturally can be extended to the class of all closed 1-forms. For every symplectic (Poisson) manifold $`M`$ with $`H^1(M)0`$ we may consider Hamiltonian Systems generated by the closed 1-form $`dH`$ where the function $`H`$ is multivalued. Instead of energy levels $`H=const`$ we have to consider nontrivial codimension 1 foliation $`dH=0`$ with Morse (or Morse-Bott) singularities. We are coming to the topological problems of studying such foliations. It has been posed in . Several participants of my seminar (A.Zorich, Le Tu Thang, L.Alania) have made very important contribution to the study of this subject. Interesting quasiperiodic structure appears here. It is not revealed fully in my opinion (see references and discussions in the article ). Multivalued Hamiltonians in real physics. I started to look around in 1982 asking the following question: can you find such systems in real physics where Hamiltonian or some other important integral of motion is multivalued (i.e. $`dH`$ is well-defined as a closed 1-form)? Much later people realized that in the theory of the so-called Landau-Lifshitz equation (which is a well-known physical integrable system with zero-curvature representation elliptic in the spectral parameter) the momentum is a multivalued functional. At that time (1982) I found only one such system describing motion of the quantum (”Bloch”) electron in the single crystal D-dimensional normal metal (D=1,2,3) under the influence of the homogeneous magnetic field $`B`$. We are working here with one-particle approximation for the system of Fermi particles whose temperature is low enough. For the zero temperature our electrons fill in all one-particle quantum Bloch states $`\psi _p`$ below the so-called ”Fermi Level” $`ϵϵ_F`$. Its value depends on the number of electrons in the system. It is the intrinsic characteristic of our metal. The index $`p`$ here may be considered finally as a point in the torus $`T^D`$ defined by the reciprocal lattice dual to the crystallographic one $$pT^D,T^D=R^3/\mathrm{\Gamma }^{}$$ There is a Morse function $`ϵ(p):T^DR`$ (dispersion relation) such that the domain $`ϵϵ_F`$ in the torus $`T^D`$ is filled in by Bloch electrons. Its boundary $`ϵ=ϵ_F`$ is a closed surface $`M_FT^D`$ for $`D=3`$. We call it Fermi Surface. It is homologous to zero in the group $`H_2(T^3,Z`$. For finite but very small temperature all essential events are happening nearby the Fermi Surface. Add now a homogeneous magnetic field to our system (i.e. put metal in the magnetic field $`B`$). Nobody succeeded in constructing a suitable well-founded theory for the exact description of electrons in the magnetic field and lattice. Irrational phenomena appear in the spectral theory of Schrodinger operators and destroy all geometric picture. However, since the late 50s physicists have used some sort of adiabatic approximation which they call ”semiclassical”. Let me warn you that this approximation has nothing to do with the standard understanding of semiclassical approximation. We take dispersion relation $`ϵ(p)`$ as a function on the torus $`T^3`$ extracted from the exact solution of the one-particle Schrodinger operator in the lattice without magnetic field. We consider a phase space $`T^3\times R^3`$ with coordinates $`p_i,x^j,i,j=1,2,3`$ and Poisson bracket of the form: $$\{x^j,x^k\}=0,\{x^j,p_k\}=\delta _k^j,\{p_j,p_k\}=B_{jk}$$ where $`B_{jk}(x)`$ are components of the magnetic field $`B`$ treated as a 2-form. Our space $`R^3`$ is Euclidean, so we can treat magnetic field as a vector $`B`$ with components $`B^j`$. Take now the function $`ϵ(p)`$ as a Hamiltonian. It generates through the Poisson structure above a Hamiltonian system in the phase space $`T^3\times R^3`$. For the homogeneous (i.e., constant) magnetic field we can see that our phase space projects on the torus $`T^3`$ with Poisson bracket $`\{p_j,p_k\}=B_{jk}`$. This Poisson bracket has a Casimir (Annihilator) $`C_B(p)=ϵ^{ijk}p_iB_{jk}=B^ip_i`$. This Casimir is multivalued: it is defined by the closed 1-form $`\omega _B=_iB^idp_i`$ on the torus. As you will see, this is the main reason for the appearence of nontrivial topological phenomena in this problem. Our Hamiltonian $`H=ϵ(p)`$ depends on the variable $`p`$ only. Therefore all important information can be extracted from the Hamiltonian system on the 3-torus with Poisson bracket defined by the magnetic field. The electron trajectories for the low temperature can be described as a curves in this torus such that $$ϵ(p)=ϵ_F,C_B(p)=const$$ However, the levels of the Casimir on Fermi surface are in fact leaves of foliation given by the closed 1-form restricted on the Fermi surface $$\omega _B\|_{M_F}=\underset{i}{}B^idp_i\|_{M_F}=0$$ Some people in ergodic theory studied in fact the most generic ergodic properties of ”foliations with transversal measure” on the Riemann surfaces. In a sense, our situation is a partial case of that. However, our picture in 3-torus is nongeneric in that sense. We cannot apply any results of that theory. We have to work with foliations obtained in the 3-torus by this special procedure only. Our use of word ”generic” here is resticted by that reqiurement. As we shall see, ergodicity is a nongeneric property within this physically realizable subclass of foliations 2-surfaces given by the closed 1-form. What is interesting is that ergodic examples exist in our picture but they occupy a measure zero subset on the sphere of directions of the magnetic fields (if generic Fermi surface is fixed). As I realized in 1982 (see ), this picture leads to nontrivial 3-dimensional topology, and I posed it as a purely topological problem to my students. The first beautiful topological observation was made by A.Zorich for the magnetic fields closed to the rational one. After new discussion and reconsidering all conjectures (see ), I.Dynnikov made a decisive breakthrough in the topological understanding of this problem for the generic directions of magnetic fields (see ). S.Tsarev constructed in 1992 the first nontrivial ergodic examples, later improved by Dynnikov see ). However, several years passed before some physical results were obtained (see the first remark about the possibility of that in my article ). We made a series of joint works with A.Maltsev (see ) dedicated to physical applications. Essentially, we borrowed topological results from the works of Zorich and Dynnikov. However, the needs of applications required that we not apply their theorems directly, but extract the key points from the proofs and reformulate them. So the modern topological formulations of these results are by-products of these works with applications (see the most modern survey in ). Let me formulate here our main physical results and after that explain the topological background and generalizations. This picture has been extensively used in solid state physics since the late 50s. The leading theoretical school in that area has been the Kharkov-Moscow school of I.Lifshitz and his pupils, like M.Azbel, M.Kaganov, V.Peschanski, A.Sludskin and others. You may find all proper quotations to physics literature in the survey article . The following fundamental Geometric Strong Magnetic Field Limit was formulated by that school (and fully accepted later by the physics community): All essential phenomena in the conductivity of normal metals in strong magnetic field should follow from the geometry of the dynamical system described above. How to understand this principle? You have to take into account that this picture certainly will be destroyed by the ”very strong” magnetic field where quantum phenomena (of the magnetic origin) are important. It should happen for such magnetic fields that magnetic flux through the elementary lattice cell is comparable with quantum unit. However, the lattice cell in solid state physics is so small that you need for that magnetic field the order of magnitude $`B10^8Gauss`$ or $`B10^4t`$ where 1t is equal to $`10^4Gauss`$. Therefore we are coming to the conclusion that even for the ”real strong” magnetic fields like $`10^2t`$ this picture still works well. For our goal we need to consider such metals that Fermi surface is topologically nontrivial. It means precisely that the imbedding homomorphism of fundamental groups $$\pi _1(M_F)\pi _1(T^3)=Z^3$$ is onto. As people have known already for many years, the noble metals like copper, gold, platinum and others satisfy to this requirement. Probably the very first time this property was found was by Pippard in 1956 for copper. Many other materials with really complicated Fermi surfaces are known now. By definition, the electron orbit is compact if it is periodic and homotopic to zero in $`T^3`$. Therefore it remains compact on the covering surface in $`R^3`$, where $`R^3`$ is a universal covering space over the torus $`T^3`$. All other types of trajectories will be called noncompact. Normally all pictures in physics literature are drawn in $`R^3`$, but everybody knows that quasimomentum vectors $`p_1,p_2R^3`$, such that $`p_1p_2`$ belongs to the reciprocal lattice $`\mathrm{\Gamma }^{}`$, are physically identical. The Lifshitz group started to study this dynamical system about 1960 and made the first important progress. For example, Lifshitz and Peschanski found some nontrivial examples of noncompact orbits stable under the variation of the direction of magnetic field. It looks like nobody could understand them properly in the physics community at that time. It was several decades before this community started to understand the geometry of dynamical systems. The Lifshitz group was ahead of its time. They made some mistakes leading to wrong conclusions and investigations were stopped. You may find the detailed discussion in our survey article . Their mistakes have been found only now because they contradicted our final results describing the conductivity tensor. Our main results: Consider projection of the conductivity tensor on the direction orthogonal to magnetic field. This is a $`2\times 2`$ tensor $`\sigma _B`$. Applying any weak electric field $`E`$ orthogonal to $`B`$, we get current $`j`$. Its projection $`\sigma _B(E)`$ orthogonal to $`B`$ is only what is interesting for us now. We claim that for the strong magnetic field $`\|B\|\mathrm{}`$ of the generic direction in $`S^2`$ only two types of asymptotics are possible: Topologically Trivial Type: $$\sigma _B0,\|B\|\mathrm{}$$ More exactly, we have $`\sigma _B=O(\|B\|^1)`$ for the topologically trivial type. All directions with trivial type occupy a set $`U_0`$ of measure equal to $`\mu _0`$ on the two-sphere $`U_0S^2`$. Topologically Nontrivial Type: $$\sigma _B\sigma _B^0+O(\|B\|^1$$ Here $`2\times 2`$ tensor $`\sigma _B^0`$ is a nontrivial limit for the conductivity tensor. We claim that it has only one nonzero eigenvalue on the plane orthogonal to $`B`$. Let us describe the topological properties of this limiting conductivity tensor. It has exactly one eigen-direction $`\eta =\eta _B`$ with eigenvalue equal to zero. Consider any small variation $`B^{}`$ of the magnetic field $`B`$. For the new field $`B^{}`$ we have an analogous picture if perturbation is small enough. We have a new $`2\times 2`$ tensor $`\sigma _B^{}`$ with one zero eigen-direction $`\eta _B^{}=\eta ^{}`$. Our statement is that the plane $`a\eta +b\eta ^{},a,bR`$, generated by this pair of directions, is locally stable under the variations of magnetic field. This plane is integral (i.e. generated by two reciprocal lattice vectors). It contains zero eigen-directions $`\eta _{B^{\prime \prime }}`$ for all small variations of the magnetic field $`B`$. It can be characterized by 3 relatively prime integer numbers $`m=(m_1,m_2m_3)`$. This triple of integer numbers is a measurable topological invariant of the conductivity tensor. An open set of directions $`U_mS^2`$ with measure $`\mu _m`$ corresponds to this type. The total measure of all these types is full: $$\mu _0+\underset{mZ^3}{}\mu _m=4\pi $$ We started to look in the old experimental data obtained in the Kapitza Institute in the 60-s by Gaidukov and others (see references in ). They measured resistance for the single crystal gold samples in the magnetic field about 2t-4t following the suggestion of Lifshitz. Confirming the ideas of the Lifshitz group, several domains with nonisotropic behavior of conductivity were found and many suspicious ”black” dots (maybe domains of small size) on the sphere $`S^2`$. It is not hard to see even now that several larger domains in these data with nonisotropic conductivity should correspond to the simplest stable topological types like $`(\pm 1,0,0),(\pm 1,\pm 1,0),(\pm 1,\pm 1,\pm 1)`$ up to permutation in the natural basis of this cubic lattice. However, for good checking it would be nice to increase magnetic field to 20t-40t for a more decisive conclusion. The black dots either correspond to the smaller domains with larger values of the topological integers or to some ergodic regimes occupying measure zero set on the sphere. For the final decision these experiments should be repeated and increased about 10 times magnetic field and smaller temperature like $`10^2K`$. Let me explain now the topological background of these results. Consider the generic Morse function $`ϵ:T^3R`$ and its generic nonsingular level $`M_FT^3,ϵ=ϵ_F`$ in the torus and in the covering space $`M^{}R^3`$. We call the surface $`M^{}R^3`$ a periodic surface. Apply now generic magnetic field $`B`$ and make the following construction: Remove all nonsingular compact trajectories (NCT) from the periodic surface $`M^{}`$ and its image $`M_F`$ in the torus. The remaining part is exactly some surface with boundary if it is nonempty: $$M_F\backslash (NCT)=M_i$$ (i.e. Fermi surface minus all NCT is equal to the union of surfaces with boundary). We call these surfaces $`M_i`$ and their closure below the Carriers of Open Trajectories. All boundary curves are the separatrix type trajectories homotopic to zero in $`T^3`$. They bound 2-discs in the corresponding planes orthogonal to magnetic field $`B`$. Let us fill them by these discs in the planes. We get closed piecewise-smooth surfaces $`\overline{M}_i`$. We denote their homological classes by $`z_iH_2(T^3,Z)`$. We use the following extract from the proofs of the main theorems of Zorich and Dynnikov (see ; their theorems have not been formulated in that way, but you may extract these key points from the proofs): In the generic case all these homology classes are nontrivial and equal to each other up to sign $`0z_i=\pm zH_2(T^3,Z)`$ where $`z`$ is some indivisible class in this group. All these closed surfaces have a genus equal to 1. As you may see, this statement means in fact some kind of the ”Topological Complete Integrability” of our systems on the Fermi surfaces for the generic magnetic field. For obtaining our final result on the conductivity tensor, we need to use the Kinetic Equation for the quasiparticles based on Bloch waves nearby the Fermi level. This equation has been used a lot by solid state physicists for the past 30 years. For the small (but nonzero) temperature, strong magnetic field and apropriate general assumptions on the impurities, the motion of quasiparticles concentrates along the electron trajectories above. This fact leads to our conclusions. Despite the fact that this theory is considered a well established one already for many years in the physics community, any attempt to prove such things as the rigorous mathematical theorems would be a huge mess. As we see, our final conclusion is separated from all theorems by some gap which cannot be eliminated. Let me point out that it is always so. ”Rigorous proofs” in mathematical physics never prove anything in real world physics. What about nongeneric trajectories? Tsarev and Dynnikov constructed very interesting examples where genus of carriers of the open trajectories is larger than 1 (see). We call such cases stochastic. Sometimes we call them ergodic. There were some attempts to extract from their properties highly nontrivial asymptotics of the conductivity tensor in the strong magnetic field . However, these attempts need a better understanding of the properties of such trajectories. We have to answer the following questions: 1.How many directions of the magnetic field on the sphere $`S^2`$ admit ergodic trajectories? According to my conjecture, for the generic Fermi surface, this set of directions has a Hausdorf dimension not greater than some number $`a<1`$ on the sphere $`S^2`$. For the special Fermi surfaces $`ϵ=0`$ of the even functions like $`cosp_1+cosp_2+cosp_3=0`$, we expect to have ergodic trajectories for the set of directions with Hausdorf dimension like $`1<a<2`$. Dynnikov started to investigate this example in his Thesis and proved several general properties. Recently R.Deleo investigated such kinds of examples more carefully and performed more detailed calculations (). His results confirm our conjectures. However, the Hausdorf dimension of this set has been unknown in this example until now. 2.Which geometric properties does ”typical” ergodic trajectory have? According to the conjecture of Maltsev, these trajectories are typically the ”asymptotically self-similar” plane curves in the natural sense. His idea (if it is true) leads to the interesting unusual properties of the asymptotic conductivity tensor. Anyway, this problem is very interesting. Dynnikov investigated also the dependence of these invariants on the level $`ϵ_F`$ of the dispersion relation (see ). These results are useful for the right understanding of our conjectures. Multidimensional Generalizations. Consider the following problem: What can be said about topology of the levels $`f(x,y)=const`$ of the quasiperiodic functions with $`m`$ periods on the plane $`x,y`$? For the case $`m=3`$ this problem exactly coincides with our subject above: By definition, quasiperiodic function on the plane is a restriction on the plane $`R^2R^m`$ of the $`m`$-periodic function. Our space $`R^3`$ was a space of quasimomenta (more precisely, its universal covering). Our plane was orthogonal to the magnetic field. Can this theory be generalized to the case $`m>3`$? According to my conjecture, it can be generalized to the case $`m=4`$. I think that for small perturbations of the rational directions this theory can be generalized to any value of $`m`$. We consider now any 4-periodic function $`f:R^4T^4R`$ and pair of the rational directions $`l_1^0,l_2^0`$ corresponding to some lattice $`Z^4`$ in $`R^4`$. Let me formulate the following theorem. Theorem. There exist two nonempty open sets $`U_1,U_2`$ on the sphere $`S_3`$ containing the rational directions $`l_1^0,l_2^0`$ correspondingly such that: For every plane $`R_l^2R^4`$ from the family given by 2 equations $`l_1=const,l_2=const`$, the quasiperiodic functions $`f_l`$ have only the following two types of connectivity components of the levels $`f_l=const`$ on the plane $`R_l^2`$. 1.The connectivity component of the level is a compact closed curve on the plane. 2. The connectivity component of the level is an open curve lying in the strip of finite width between 2 parallel straight lines with the common direction $`\eta `$. This situation is stable in the following sense. After any small variations of the directions $`l_1U_1,l_2U_2`$, of function $`f`$ on $`T^4`$ or the level we still have such open component with direction $`\eta ^{}`$. For all possible perturbations this set of directions $`\eta ,\eta ^{},\mathrm{}`$ belong to some integral 3-hyperplane in $`R^4`$. This property can be formulated in terms of the integral homology class in the group $`H_3(T^4,Z)`$ and of the torical topology of the carriers of the open trajectories. The idea of the proof was recently published by the author in . We may reformulate this problem in terms of Hamiltonian systems. Let the constant Poisson Bracket $`B_{ij}`$ be given on the torus $`T^m`$ whose rank is equal to 2. Any Hamiltonian $`f`$ generates such systems whose trajectories are equal to the levels of $`f`$ on the planes. Our theorem means that in these cases this Hamiltinian system is Completely Integrable in the specific topological sense described above.
warning/0004/hep-ph0004050.html	ar5iv	text	# HIP-2000-18/TH Flat direction condensate instabilities in the MSSM ## 1 Introduction The scalar potential of the Minimal Supersymmetric Standard Model (MSSM) is a complicated function of 45 complex squark and slepton fields and 4 complex Higgs fields. A natural feature of this potential is the existence of flat directions, corresponding to linear combinations of these fields such that there are no renormalizable contributions to the scalar potential along the flat direction beyond the soft SUSY breaking terms . These flat directions may play a fundamental role in the cosmology of the MSSM, in particular as a natural source of the baryon asymmetry of the Universe via the Affleck-Dine mechanism , in which a baryon asymmetry is induced in a coherently oscillating condensate of squarks and sleptons. It has recently been realized that the cosmology of flat directions of the MSSM and Affleck-Dine baryogenesis may be more complicated than previously thought . If the flat direction scalar potential increases less rapidly than $`\varphi ^2`$, where $`\varphi `$ is the scalar field along the flat direction (”Affleck-Dine (AD) scalar”), then once coherent oscillations of the field begin, the condensate will have a negative pressure making it unstable with respect to spatial perturbations . These spatial perturbations, which generally arise during inflation as a result of quantum fluctuations of the AD scalar , will grow and go non-linear, fragmenting the condensate into lumps which eventually evolve into Q-balls of baryon number (”B-balls”) . The formation of Q-balls has recently been verified by lattice simulations . The cosmological evolution of the resulting Q-balls depends on the form of SUSY breaking. Q-balls from the AD condensate were first discussed in the context of gauge-mediated SUSY breaking , in which the flatness of the potential occurs above the mass of the messenger fields. It was later realized that Q-balls can also form in the more conventional gravity-mediated SUSY breaking models . In this case the potential is likely to be flatter than $`\varphi ^2`$ as a result of radiative corrections. The resulting Q-balls are not stable but can be very long-lived, decaying after the electroweak phase transition has occured . As a result Q-balls can protect the baryon asymmetry from the effect of lepton number violating interactions combined with sphaleron processes and, if they decay at a low enough temperature, can also act as a common source of baryons and dark matter neutralinos, so explaining the similarity of the number densities of baryons and dark matter particles when the dark matter particles have masses of the order of $`m_W`$ . They can also enhance the isocurvature density perturbations expected from AD baryogenesis in particular inflation models . The existence and cosmology of Q-balls in the MSSM with gravity-mediated SUSY breaking is dependent upon the details of the radiative correction to the flat direction scalar potential. The correction must be negative in sign relative to the SUSY breaking mass squared term in order for Q-balls to form, whilst the binding energy of the charges within the Q-ball and so the total charge and lifetime of the Q-ball depends upon the magnitude of the correction , as does the fraction of the total baryon number initially trapped within the Q-balls . In this letter we consider in detail the radiatively corrected scalar potential for a range of flat directions using the renormalization group. We will show that condensate collapse and Q-ball formation are almost a general feature of all the MSSM flat directions, with exceptions only for the cases of d=4 $`H_uL`$-direction and directions with large admixtures of stop. The existence of instabilities along the latter directions depend on the mass of the stop in a testable way, as will be discussed in the following. F- and D-flat directions of the MSSM have been classified and listed in . For gravity-mediated SUSY breaking the scalar potential along a flat direction has the form $$U(\mathrm{\Phi })m^2(1+K\mathrm{log}\left(\frac{\|\mathrm{\Phi }\|^2}{M^2}\right))\|\mathrm{\Phi }\|^2+\frac{\lambda ^2\|\mathrm{\Phi }\|^{2(d1)}}{M_{}^{2(d3)}}+(\frac{A_\lambda \lambda \mathrm{\Phi }^d}{dM_{}^{d3}}+h.c.),$$ (1) where $`m`$ is the conventional gravity-mediated soft SUSY breaking scalar mass term ($`m100\mathrm{GeV}`$), $`K`$ is a parameter which depends on the flat direction, and $`d`$ is the dimension of the non-renormalizable term in the superpotential which first lifts the degeneracy of the flat direction; it is of the form $`f=\lambda M_{}^{4d}\mathrm{\Phi }_1\mathrm{}\mathrm{\Phi }_d`$. We assume that the natural scale of the non-renormalizable terms is $`M_{}`$, where $`M_{}=M_{Pl}/\sqrt{8\pi }`$ is the supergravity mass scale. The equation of state for a field oscillating in a potential $`\varphi ^\gamma `$ reads as $`p=(\frac{2\gamma }{\gamma +2}1)\rho `$ so that Eq. (1) gives rise to the equation of state $`p=\frac{K}{2}\rho `$. With $`K<0`$ the pressure is negative and hence the AD condensate is unstable. The condensate fragments and eventually the lumps will evolve dynamically into the state of lowest energy, the Q-ball. If $`K>0`$ the condensate is stable and no Q-balls will form. Although a negative $`K`$ should be a generic feature of the MSSM, not all flat directions will have negative $`K`$ for all the values of the MSSM parameters. Moreover, as the actual value of $`K`$ dictates the dynamical evolution of the AD condensate and its fragmentation, it is of great interest to find out the precise value of $`K`$ for a given flat direction and for a range of the MSSM parameter values. $`K`$ can be computed from the RG equations, which to one loop have the form $$\frac{m_i^2}{t}=\underset{g}{}a_{ig}m_g^2+\underset{a}{}h_a^2(\underset{j}{}b_{ij}m_j^2+A^2),$$ (2) where $`a_{ig}`$ and $`b_{ij}`$ are constants, $`m_g`$ is the gaugino mass, $`A`$ is the A-term, $`h_a`$ the Yukawa coupling, and $`t=\mathrm{ln}M_X/\mu `$. The full RG equations are listed in and we do not reproduce them here. We assume unification at $`t=0`$ and neglect all other Yukawa couplings except the top Yukawa $`h_t(M_W)=1`$ (for definiteness, we choose $`tan\beta =1`$). We shall use quantities scaled by $`m`$ and denote $`m_g/m`$ at $`t=0`$ by $`\xi `$. The potential along the flat direction is then characterized by the amount of stop mixture (where appropriate), the values of $`\xi `$ and $`A`$, and in the special case of the d=4 $`H_uL`$ -direction, on the $`H_uH_d`$ -mixing mass parameter $`\mu _H`$. The mass of the AD scalar $`\varphi `$ is the sum of the masses of the squark and slepton fields $`\varphi _i`$ constituting the flat direction, $`m_S^2=_ap_i^2m_i^2,`$ where $`p_i`$ is the projection of $`\varphi `$ along $`\varphi _i`$, and $`p_i^2=1`$. The parameter $`K`$ is then given simply by $$K=\frac{m_S^2}{t}\|_{t=\mathrm{log}\mu }.$$ (3) To compute $`K`$, we have to choose the scale $`\mu `$. The appropriate scale is given by the value of the AD field when it first begins to oscillate at $`Hm`$. Let the value of the mean field be $`\varphi _0`$; the value of $`K`$ at this scale then determines whether the condensate is unstable or not. We may compute $`\varphi _0`$ from Eq. (1) by ignoring the radiative correction (i.e. setting effectively $`K=0`$) and minimizing the U(1) symmetric part of $`U`$ (or neglecting the A-term in Eq. (1)). One then finds $$\mu =\|\varphi _0\|=\left[\frac{m^2M_{}^{2(d3)}}{(d1)\lambda ^2}\right]^{\frac{1}{2(d2)}},$$ (4) where in what follows we assume for simplicity that $`\lambda =1`$. Table 1. The flat directions considered \| direction \| dimension \| $`\mathrm{mass}^2`$ \| \| --- \| --- \| --- \| \| $`H_uL`$ \| $`4`$ \| $`\frac{1}{2}(m_H^2+\mu _H^2+m_{\stackrel{~}{L}_i}^2)`$ \| \| $`uude`$ \| $`4`$ \| $`\frac{1}{4}(m_{\stackrel{~}{u}_i}^2+m_{\stackrel{~}{u}_j}^2+m_{\stackrel{~}{d}_k}^2+m_{\stackrel{~}{e}_l}^2);(ij)`$ \| \| $`QQQL`$ \| $`4`$ \| $`\frac{1}{4}(m_{\stackrel{~}{Q}_i}^2+m_{\stackrel{~}{Q}_j}^2+m_{\stackrel{~}{Q}_k}^2+m_{\stackrel{~}{L}_l}^2;(ij\mathrm{or}k)`$ \| \| $`(udd)^2`$ \| $`6`$ \| $`\frac{1}{3}(m_{\stackrel{~}{u}_i}^2+m_{\stackrel{~}{d}_j}^2+m_{\stackrel{~}{d}_k}^2);(jk)`$ \| \| $`(QLd)^2`$ \| $`6`$ \| $`\frac{1}{3}(m_{\stackrel{~}{u}_i}^2+m_{\stackrel{~}{d}_j}^2+m_{\stackrel{~}{Q}_k}^2)`$ \| The flat directions we have studied are listed in Table 1. There is one purely leptonic direction, $`H_uL`$, a purely baryonic one involving only the squarks, and directions which have both squark and lepton fields. In Fig. 1 we show the contours of $`K`$ for the d=4 $`uude`$ and $`QQQL`$ directions in the $`(A,\xi )`$ -plane; Fig. 2 is for the d=6 $`(udd)^2`$ and $`(QLd)^2`$ directions. These should be representative of all the other directions, too, except for $`H_uL`$. For $`\xi 𝒪(1)`$, typical value for $`K`$ is found to be about $`0.05`$. For all the squark directions with no stop, as long as $`h_b`$ and $`h_u`$ can be neglected, $`K`$ is always negative, and the contours of equal $`K`$ do not depend on $`A`$. This is evident from the RGEs Eq. (2). However, as far as flat directions are concerned, all squarks are equal, and having no stop mixture would appear rather unnatural. Therefore we have considered the effect of stop mixing in the squark directions, which results in $`A`$-dependence of the $`K`$-contours, as is depicted in Figs. 1 and 2. In the presence of stop mixing $`K<0`$ is no longer automatic even in the purely squark directions. The more there is stop, the larger value of $`\xi `$ is required for $`K<0`$. However, even for pure stop directions, positive $`K`$ is typically obtained only for relatively light gaugino masses with $`\xi <\mathrm{\hspace{0.33em}0.5}`$. In d=4 directions the effect of stop mixture is less pronounced than in the d=6 directions, as can be seen from Figs. 1 and 2. In contrast to the squark directions, $`K`$ was found to be always positive in the $`H_uL`$-direction. This is due to the fact $`H_uL`$ does not involve strong interactions which in other directions are mainly responsible for the decrease of the running scalar masses. The value of $`\mu _H`$ was chosen in such a way that for each value of $`A`$ and $`\xi `$, electroweak symmetry breaking is obtained at the scale $`M_W`$. Except for the $`H_uL`$-direction, the instability of the AD condensate is thus seen to depend on the amount of stop mixture in the flat direction. Since it would be natural to expect roughly equal mixtures of the different squark flavours in the flat direction scalar, this means that in principle instabilities, and hence Q-balls in the case of gravity mediated susy breaking, could be ruled out experimentally at LHC by measuring the mass of the gluino and the stop. To illustrate this, in Fig. 3 we show the the domains of positive and negative $`K`$ in the region of $`(m_{\stackrel{~}{t}},m_{\stackrel{~}{g}})`$ -plane corresponding to the range $`3<A<3`$ and $`0<\xi <2.5`$. The Figure is for the $`(udd)^2`$ direction with equal weight for all $`u`$-squarks. For most part $`K>0`$ and $`K<0`$ regions can be separated, although there is a small area below the upper $`K=0`$ contour where both values can be found. For a fixed $`m`$, the $`K=0`$ contour has endpoints which correspond to $`A=\pm 3`$; if one were to allow for a wider range in $`A`$, this would spread the region between the dashed lines towards the lower right-handed corner. Changing the value of $`m`$ would redefine the physical mass scale by a factor $`m/100\mathrm{GeV}`$. Thus measuring $`m_{\stackrel{~}{t}}`$ and $`m_{\stackrel{~}{g}}`$ would not alone be sufficient to determine the existence of instabilities. In addition, one needs the values of $`m`$ and $`tan\beta `$ which naturally will be measured if supersymmetry will be found. Very roughly, instability is found when $`m_{\stackrel{~}{g}}>m_{\stackrel{~}{t}}`$, although the exact condition should be checked case by case. In conclusion, we have shown that the MSSM scalar condensates in all but the $`H_uL`$ flat direction are unstable for a large part of the parameter space. Therefore the existence of Q-balls is a generic feature in all the models that incorporate both the MSSM and inflation. Moreover, as the existence of instabilities can in principle be ruled out by measuring the mass of the stop and the gluino, one may soon be able to subject the Affleck-Dine scenario to a direct test. ### Acknowledgements This work has been supported by the Academy of Finland under the contract 101-35224 and the PPARC (UK).
warning/0004/cond-mat0004170.html	ar5iv	text	# Multifractality in Time Series ## I introduction In this work we apply the concepts of multifractal physics to financial time series in order to characterize the onset of crash for the Standard & Poor’s 500 stock index. We shall present an example of dynamic phase transition in a simple economic model system based on a mapping with multifractality phenomena in random multiplicative processes and by applying former results obtained with a continuous probability theory for describing scaling measures. An attempt is made to characterize the presence of stock market crashes by solving for a price equation from a non-linear equilibrium model and showing how multifractal physics measures can be generated from this equation. We found that an ”analogous” specific heat $`C_q`$ of the S&P500 price data displays a shoulder to the right of the main peak as a function of time lags. For large time lags, $`C_q`$ resembles a classical phase transition at a critical point. We explain this dynamic phase transition by a mapping with multifractality phenomena in random multiplicative processes. Within this description the temporal price variations of a commodity displays features of an ”analogous” phase transition from inflated to devalued prices, when the excess demand is not linear in the asset price. An analytical expression for $`C_q`$ of the economic model system is derived. Finance and physics joint ”ventures” have attracted considerable interest in the literature for over many years . These efforts have allowed to pursue analogies between stock market dynamics and stochastic models commonly used in the statistical physics of complex systems . Such parallel analysis have been useful to best quantify and understand possible correlations in financial data by measuring the autocorrelation function and the power spectrum . Another notable example is the analogy with the scaling properties observed in turbulence . Parallels have also been drawn between a very simple theory of financial markets and the quantum gauge theory (see also comments in ). Other interesting studies for investment strategies in diversified portfolio stocks have been recently discussed in terms of products of random matrices and of multiplicative random walks . Realistic price fluctuations have also been found to emerge in the adaptation of a system (i.e., the traders) to complex environments and from the self-organization of a (closed) system of traders without external influences . Thus, the financial market dynamics is still an open subject for physicists and economists. In particular, the basic principles governing the origin of stock market crashes are far from well understood by both communities , specially in regard to the problem of world-wide market dynamics. Using the renormalization group theory , it has been proposed that this cooperative phenomenon could be the result of a critical phase transition (see also ). Motivated by these suggestions, it is then tempting to determine the constrains needed to understand and describe analytically an analogous phase transition from a new perspective. That is, by using the characterization of multifractal singularities where an analogy with thermodynamics has already been established . Multifractality was initially proposed to treat turbulence and, in recent years, it has been applied successfully in many different fields ranging from model systems such as Diffusion Limited Aggregation to physiological data such as heartbeat. The multifractal analysis of financial discrete sequences developed in our paper is another aspect in a relatively new topic of physics, the so-called econophysics. ## II Multifractal Characterization This paper will focus on scaling laws in financial returns, particularly on the absolute values of returns and their scatter as a function of time. Similarly to other studies , we are not attempting to predict a price drop or rise on a specific time on the basis of past records, but we are after a new characterization of the presence of stock crashes, which is the added value of our distinctive approach. This is important to mention since forecasting has been extensively studied in the econometrics literature by other techniques such as ARCH and multivariance models depending, e.g., on the data seasonal behavior and previous experiences of the forecasters. Interesting work on fractional processes of these types in econometrics is very well known and these are beyond the scope of our work. Important to mention is that recent theoretical work has considered the possibility of fractal nature of the absolute value of returns on exchange rates . Let us introduce next our new characterization of a financial time series $`x(t)`$ based on the well-known definitions used in multifractal physics. In Fig.1 we show the temporal $`x(t+T)x(t)`$ fluctuations of the S&P500 index (data available in ) for the 12 years period 2 January 1980 - 31 December 1992 as a function of the trading time lags $`1T220`$. From this figure, it can be observed that the peaked narrow fluctuations measured for $`T=1`$ spread out on increasing $`T`$. For $`T80`$, the maximum and minimun difference values of the so-called ”Black Monday” crash measured in 1987 (shown as the largest straight line for $`T=1`$ in the figure), become comparable to the valleys and picks differences of the relative S&P500 fluctuations over wide periods of time. This peculiar behavior is shown next to lead to an analogous thermodynamic phase transition when varying $`T`$. Let us consider the following measure over $`N`$ intervals $$\mu _i(T)=\frac{\|x(t+T)x(t)\|}{_{t=1}^N\|x(t+T)x(t)\|},$$ (1) with $`T`$ representing some finite trading time lags $`T`$. Clearly, the above relation can be viewed as a normalized probability measure with $`\mu _i>0`$. From $`\mu _i`$ in Eq.(1), we then construct the corresponding generating function $`Z`$, and its moments $`q`$, which follows the scaling $$Z(q,N)=\underset{T\mathrm{}}{lim}\underset{t=1}{\overset{N}{}}\mu _i(T)^qN^{\tau (q)}.$$ (2) To get a thermodynamic interpretation of multifractality (see, e.g., ), we divide the one-dimensional system of length $`L`$ into $`N`$ lines of length $`\mathrm{}`$; thus $`NL/\mathrm{}`$. We then associated this $`N`$ with the number of discrete $`x(t)`$ time sequences considered in Eq.(1) in order to relate $`T`$, $`L`$ and $`\mathrm{}`$ in the definition of the measure. For $`\mathrm{}/L0`$, the function $`\tau `$ relates to the generalized fractal dimensions $`(q1)D_q`$. Similar multifractal analysis has recently been performed for the energy dissipation field of turbulence , logistic maps and surface roughening . By following the thermodynamic formulation of multifractal measures, we can also derive an expression for the ”analogous” specific heat as follows $$C_q\frac{^2\tau (q)}{q^2}\tau (q+1)2\tau (q)+\tau (q1).$$ (3) For large $`T`$, we shall show that the form of $`C_q`$ resembles a classical phase transition at a critical point. We shall return to these equations in Section III. ## III Dynamical Model Similarly to in our approach there is only one stock. To study the price changes for one commodity it is necessary to derive a dynamical equation which results from the prevailing market conditions. The market is usually considered competitive so it self-organizes to determine the behaviour of prices. We assume here that all factors determining the demand $`D`$ and the supply $`Q`$ other than the asset price $`p`$ remain constant over time and denote these quantities in equilibrium with an asterisks (). In the following all variables are dimensionless. In a competitive market it is expected that the rate of price increase should be a functional of the excess demand function $`E(p)=D(p)Q(p)`$. Hence one writes $`dp/dtf[E(p)]`$ . Assuming that a commodity can be stored then, in general, the flow of demand does not equal the flow of $`Q`$ output. Hence stocks of the commodity (or product) build up when the flow of output exceeds the flow of demand and vice-versa. Then the rate at which the level of stocks $`S`$ changes can be approximated as $`dS/dt=Q(p)D(p)`$. From these relations, a price adjustment relation that takes into account deviations of the stock level $`S`$ above certain optimal level $`S_o`$ (to meet any demand reasonably quickly) is then given by $$\frac{dp}{dt}=\gamma \frac{dS}{dt}+\lambda (S_oS),$$ (4) where $`\gamma `$ (i.e., the inverse of excess demand required to move prices one unity ) and $`\lambda `$ are positive parameters. If $`\lambda =0`$, the price adjusts at a rate proportional to the rate at which stocks are either raising or running down. If $`\lambda >0`$, prices would increase when stock levels are low and rise when they are high (with respect to $`S_o`$). We shall assume here $`\lambda `$ to characterize a noise term. In our description, for each asset price $`p`$, we postulate simple non-linear forms for the quantities $`D`$ demanded and $`Q`$ supplied such that $`D(p)`$ $`=`$ $`d^{}+d_o[\mathrm{\hspace{0.33em}1}{\displaystyle \frac{\delta ^2}{2!}}(pp^{})^2+\mathrm{}](pp^{}),`$ (5) $`Q(p)`$ $`=`$ $`q^{}+q_o[\mathrm{\hspace{0.33em}1}{\displaystyle \frac{\delta ^2}{2!}}(pp^{})^2+\mathrm{}](pp^{}),`$ (6) where $`d_o`$, $`q_o`$ and $`d^{}=D(p^{})`$, $`q^{}=Q(p^{})`$ are arbitrary coefficients (related to material costs, wage rate, etc), $`p^{}`$ is an equilibrium price and $`\delta `$ is our order parameter. We write $`D`$ and $`Q`$ as a Taylor series expansion with the usual linear dependence (independent of $`\delta `$) plus a non-linear correction. Higher order terms $`𝒪(4)`$ are here neglected for small $`pp^{}`$. In the above we might also consider two different $`\delta `$s, but to reduce variables to a minimum we assume $`D`$ and $`Q`$ to vary similarly from linearity. To simplify notation we also define $$\beta _oq_od_o.$$ (7) In the context of a simple economic model , it is reasonable to assume that $`S_o`$ depends linearly on the demand; e.g., $`S_o=\mathrm{}_o+\mathrm{}D`$, with $`\mathrm{}_o`$ a constant and $`\mathrm{}`$ satisfying the constrain below. The postulated linear dependence of the optimal stock level $`S_o`$ on $`D`$ at equilibrium provides a complete economic model as in . Therefore, in equilibrium (where $`\frac{dp}{dt}\|_p^{}=0`$ and $`\frac{dS}{dt}\|_S^{}=0`$, so that demand equals supply and $`S=S_o`$), from the above we obtain $$d^{}q^{}=0,S^{}=\mathrm{}_o+\mathrm{}(d^{}+d_op^{}).$$ (8) And after some little algebra we find that the price of one strategic commodity is governed by the general equation $`{\displaystyle \frac{d^2p}{dt^2}}+(\gamma \beta _o\mathrm{}\lambda d_o)[1{\displaystyle \frac{3\delta ^2}{2!}}(pp^{})^2]{\displaystyle \frac{dp}{dt}}+`$ (9) $`\lambda \beta _o(pp^{})[1{\displaystyle \frac{\delta ^2}{2!}}(pp^{})^2]0.`$ (10) The linear case is for $`\delta =0`$. This leads to $`p(t)p^{}A_1\mathrm{cos}(t\sqrt{\lambda \beta _o})+A_2sin(t\sqrt{\lambda \beta _o})`$. ### A The Simple Case $`\mathrm{}\lambda d_o=\gamma \beta _o`$ To keep the mathematics simple we choose $`p^{}=0`$ and consider $`\mathrm{}`$ to satisfy $$\mathrm{}\frac{\gamma \beta _o}{\lambda d_o},$$ (11) The general Eq.(9) then reduces to $$\frac{d^2p}{dt^2}+\lambda \beta _op\frac{\delta ^2\lambda \beta _o}{2}p^30.$$ (12) This is our dimensionless price adjustment equation which gives rise to a burst as discussed later. When $`\delta 0`$ and $`[\lambda \beta _o,\delta ^2\lambda \beta _o/2]>0`$, it has the well-known kink solutions $$p(t)=\pm \frac{\sqrt{2}}{\delta }tanh(\sqrt{\frac{\lambda \beta _o}{2}}t).$$ (13) Clearly $`\beta _o`$ of Eq.(7) must be positive. Since in a free market economy the demand for a product (or commodity) fall when its price increases, then it is reasonable to assume $`d_o<0`$ in Eq.(5). As the price raises, the supply also increases; hence in general one also assumes $`q_o>0`$. These conditions yield $`\beta _o>0`$ as requested in Eq.(13) and also $`d_o\mathrm{}>0`$. In the case $`\delta <0`$ the above function for $`p(t)`$ displays a sudden decline around the equilibrium value $`p^{}`$ taken to be $`p(t^{})=p(0)=0`$. ### B The Case $`\mathrm{}\lambda d_o\gamma \beta _o`$ If we consider the case in which $`\mathrm{}\lambda d_o\gamma \beta _o`$, the price equation of Eq.(9) results in a Lienard-type of equation: $`p^{\prime \prime }+g_1(p)p^{}+g_o(p)=0`$ due to the presence of the $`\frac{dp}{dt}`$-term. What we present next is a brief discussion regarding its possible solutions. Again, we set $`p^{}=0`$. With the aid of the substitution $`\frac{dp}{dt}w(p)`$, so that $`\frac{d^2p}{dt^2}w(p)\frac{dw}{dp}`$, the Lienard equation can be reducible to an Abel equation of the second kind $`ww^{}=f_1(p)w+f_2(p)`$. For small $`\delta `$, the substitution $`w(p)=p^3K(p)+\frac{p}{4}(\mathrm{}\lambda d_o\gamma \beta _o)`$ leads to the Bernoulli equation with respect to $`p=p(K)`$: $$3K(p)[\frac{\delta ^2}{8}(\mathrm{}\lambda d_o\gamma \beta _o)K(p)]\frac{dp}{dK}=K(p)p+\frac{(\mathrm{}\lambda d_o\gamma \beta _o)}{4p},$$ (14) whose solution is $`p^2=[{\displaystyle \frac{\delta ^2}{8}}(\mathrm{}\lambda d_o\gamma \beta _o)K(p)]^{1/3}\{K_o+{\displaystyle \frac{(\gamma \beta _o\mathrm{}\lambda d_o)}{8}}\times `$ (15) $`\times ({\displaystyle \frac{1}{3K(p)}})_2^{2/3}F_1({\displaystyle \frac{1}{3}},{\displaystyle \frac{2}{3}},{\displaystyle \frac{5}{3}},{\displaystyle \frac{\delta ^2}{8}}{\displaystyle \frac{(\mathrm{}\lambda d_o\gamma \beta _o)}{K(p)}})\},`$ (16) with $`K_o`$ an arbitrary constant and $`{}_{2}{}^{}F_{1}^{}`$ the first, complex hypergeometric function as arising in many physical problems. It converges within the unit circle $`\|(\mathrm{}\lambda d_o\gamma \beta _o)/K(p)\|<2/\delta ^2`$. By examining these equations one recognizes that if $`\mathrm{}\lambda d_o\gamma \beta _o`$, other behaviour might appear for $`p(t)`$ (different from the one in Eq.(13)). However, it is important to note that such possible behaviour can essentially be found in the complex plane since the solutions of Eq.(15) are driven by the $`K^{2/3}`$-term and depend whether $`\gamma \beta _o`$ is greater or smaller than $`\mathrm{}\lambda d_o`$. A detailed analysis of such solutions is beyond the scope of this work. We only study here the simplest dynamical economic model where all its variables, including prices, are related to the demand and supply functions as seen in a free market. ## IV Model and Analogous Phase Transition Let us now identify the behaviour of $`p(t)`$ with an analogous phase transition as seen in multifractals. Following the analogy with critical phenomena in the time domain as proposed in , we also consider $`t`$ to be the relevant variable for the analysis of an possible existence of an analogous critical point. To derive a connection between our economic model and multifractality we first briefly review multifractal phenomena in Random Multiplicative Processes. ### A Multifractality in Random Multiplicative Processes Multifractality emerges in random multiplicative processes for a self-similar function $`\psi `$ such that is rescaled as $`\widehat{\psi }(x)=e^{L(1)x}\psi (x)`$, where $`L(1)`$ is the generalized Lyapunov exponent for the first moment of $`\psi `$ and $`x`$ is a space variable for $`N`$-disorder fluctuations on a unit interval . All information about these systems is embodied in the non-zero, positive $`\psi `$ measures. To analyse the analogy of multifractality with thermodynamics one then scales the moments $`q`$ of the functions $`\widehat{\psi }`$ with respect to segments $`ł1`$ as $$Z(q,ł)=\underset{N\mathrm{}}{lim}\underset{k=1}{\overset{ł}{}}\widehat{\psi }_{k,N}^qł^{\tau (q)},$$ (17) which defines the exponents $`\tau `$ and $`Z`$ is a formal partition function. It has been previously shown that for general random multiplicative processes , $`\tau `$ satisfies $$\tau (q)=(1q)\frac{1}{h}\{L(q)qL(1)\},$$ (18) where $`L(q)`$ $`=`$ $`\underset{N\mathrm{}}{lim}{\displaystyle \frac{1}{N}}\mathrm{ln}{\displaystyle \underset{k=1}{\overset{ł}{}}}\psi _{k,N}^qh,`$ (19) $`h`$ $`=`$ $`\underset{N\mathrm{}}{lim}{\displaystyle \frac{1}{N}}\mathrm{ln}ł^1.`$ (20) We shall use these findings to derive a connection between our economic model variables and multifractality phenomena. From this connection, we shall identify all the model quantities that might exhibit multifractality within the framework of a stochastic multiplicative process. These processes are known to generate power law probability density functions . ### B General Continuous Probability Theory We use next the simple continuous probability theory discussed in that allows to explore the genesis of an ”analogous” phase transition from the point of view of a nonlinear singularity spectrum equivalent to multifractals. A crucial feature of this formalism is to consider $`t\sqrt{\lambda \beta _o/2}`$ to be a continuous random variable. Then within the framework of general probability theory (see, e.g.,), the distribution function of this random variable, defined in a line and in terms of its probability distribution $`P`$, can be approximated as $`𝐏\{\zeta _1<\zeta ^{}\zeta _2\}=𝒢(\zeta _2)𝒢(\zeta _1)_{\zeta _1}^{\zeta _2}\varphi (\zeta ^{})𝑑\zeta ^{}`$, where $`\{\}`$ indicates the function interval and $`\varphi `$ is a uniform probability density that needs to be specified. Inspired by well-known results for the probability distribution function of real economic data with power-law tails , we assume $`\varphi (\zeta )(\varphi _0/2)[1(\delta /\sqrt{2})H(\zeta )]`$ such that $`\varphi (\zeta +\mathrm{})0`$ and $`\varphi (\zeta \mathrm{})\varphi _0>0`$, with $`H`$ given by the real solutions of a static, dimensionless Ginzburg-Landau-like equation. It resembles the spin-flip function of the well-known Glauber-Ising chain model. Using such solutions, we shall show next that it is possible to establish a relation with thermodynamics similarly to multifractality phenomena. It is at this point that we relate $`H(\zeta )`$ to the positive solutions $`p(t)`$ of Eq.(12) and map $`\zeta /\zeta _ot\sqrt{\lambda \beta _o/2}`$ within the framework of general probability theory. Following Ref., it is then straightforward to derive an expression for an ”analogous” specific heat, $`C_\zeta `$ , for our economic system. To obtain such an expression we evaluate first the integral of P over the range $`[\zeta _o,\zeta ]`$. Thus we have $$𝒢(\zeta )𝒢(\zeta _o)=\frac{\varphi _0}{2}_{\zeta _o}^\zeta \{1\frac{\delta }{\sqrt{2}}tanh(\frac{\zeta ^{}}{\zeta _o})\}𝑑\zeta ^{}\tau (\zeta ),$$ (21) which, in turn, defines the function $`\tau (\zeta )`$. The $`𝒢`$ functions satisfy $`𝒢(\zeta )>𝒢(\zeta _o)`$, or alternatively, $`\varphi (\zeta )>\varphi (\zeta _o)`$ (since $`𝒢/\zeta =\varphi (\zeta )`$ ), which is in accord with the above assumption for $`\varphi (\zeta \pm \mathrm{})`$. Let us see next how our analysis from general probability theory would capture multifractality. Similarly to the dielectric breakdown model or the Poisson growth model, where the local field is set proportional to the growth probability , we assume here that the continuous function $`\tau `$ (to be identified as a free energy) is proportional to the probability distribution $`𝐏`$ as in Eq.(21). After a little algebra, the above integral gives $$\tau (\zeta /\zeta _o)(1\zeta /\zeta _o)\tau (0)\frac{\delta \zeta _o\varphi _o}{2\sqrt{2}}\{\mathrm{ln}cosh(\zeta /\zeta _o)(\zeta /\zeta _o)\mathrm{ln}cosh(1)\},$$ (22) in which $`\tau (0)𝒢(0)𝒢(\zeta _o)=\frac{\zeta _o\varphi _o}{2}\{1+\frac{\delta }{\sqrt{2}}\mathrm{\Gamma }_\lambda \}`$, such that $`\mathrm{\Gamma }_\lambda \mathrm{ln}cosh(1)`$. We have assumed that $`\delta \mathrm{\Gamma }_\lambda /\sqrt{2}<<1`$, hence $`\delta <0`$. Up to this point we have not carried out any actual numeration or scaling of a particular fractal configuration or set as used in multifractal theory. However, to gain further insight and explain how the present general probability can mimics multifractal behaviour, let us now relate the above $`\tau `$ to a measure within a random multiplicative process. By mapping our results to such a measure, all other quantities that can be scaling and are directly derived from $`\tau `$, such as an analogous specific heat, will then follow. From a comparison between Eqs.(22) and (18) one can easily identify the following terms $`{\displaystyle \frac{1}{h}}`$ $``$ $`{\displaystyle \frac{\delta \zeta _o\psi _o}{2\sqrt{2}}},`$ (23) $`\tau (0)`$ $``$ $`1,`$ (24) $`L(\zeta /\zeta _o)`$ $``$ $`\mathrm{ln}\mathrm{cosh}(\zeta /\zeta _o).`$ (25) It is from this simple mapping between our results and those for random multiplicative processes that we can made an attempt to understand how the existance of the economic model order-parameter $`\delta `$ (and from it, the non-linearities in the demand and supply functions) leads to obtain multifractal-like behaviour. From the above mapping we deduce that if $`\delta 0`$, then the moments $`q`$ of the functions $`\widehat{\psi }`$ with respect to segments $`ł`$ would vanish since $`ł0`$. It follows that a linear economic model would never exhibit multifractal features since in this case $`\tau (q)=q1`$. Furthermore, it is important to note that a Lyapunov exponent of the type we derive resembles that of a random multiplicative process with $`\mathrm{}`$ described by the probability distribution $`P(\mathrm{})=(1/n)_{i=1}^n\delta (\mathrm{}\mathrm{}_i)`$ . By considering the actual definition of $`L(q)`$, the moments of the exponential measures in a random multiplicative process might well be related to our reduced variable as $`q\zeta /\zeta _ot\sqrt{\lambda \beta _o/2}`$. ### C Possible Multifractal Features To analyse multifractal features in our economic model, when mapped into a random multiplicative process as discussed above, we consider standard definitions: $`\tau (\zeta /\zeta _o)[(\zeta /\zeta _o)1]D_\zeta `$. According to such definitions , the function $`\tau `$ represents an ”analogous” free energy and $`D_\zeta `$ the multifractal dimension. From Eq.(22) it follows that $$D_\zeta \tau (0)+\frac{\delta \zeta _o\varphi _o}{2\sqrt{2}(1\zeta /\zeta _o)}\{\mathrm{ln}cosh(\zeta /\zeta _o)(\zeta /\zeta _o)\mathrm{ln}cosh(1)\},$$ (26) such that $`\zeta \zeta _o`$. From this relation we obtain $`D_{\zeta 0}=\tau (0)`$, and $`D_{\zeta +\mathrm{}}=D_{\zeta 0}\frac{\delta \zeta _o\varphi _o}{2\sqrt{2}}\{1+\mathrm{\Gamma }_\lambda \}`$; $`D_\zeta \mathrm{}=D_{\zeta 0}+\frac{\delta \zeta _o\varphi _o}{2\sqrt{2}}\{1\mathrm{\Gamma }_\lambda \}`$. If $`\zeta =\zeta _o`$, then $`D_{\zeta \zeta _o}=D_{\zeta 0}\frac{\delta \zeta _o\varphi _o}{2\sqrt{2}}\{\mathrm{\Gamma }_\lambda +tanh(1)\}`$. Complementary to $`\tau `$ we also define $$\alpha (\zeta /\zeta _o)\frac{}{(\zeta /\zeta _o)}\tau (\zeta )D_{\zeta 0}\frac{\delta \zeta _o\varphi _o}{2\sqrt{2}}\{\mathrm{\Gamma }_\lambda +tanh(\zeta /\zeta _o)\}.$$ (27) It can be easily shown that $`\alpha _{max}\alpha (\zeta /\zeta _o\mathrm{})=D_{\zeta /\zeta _o\mathrm{}}`$, and $`\alpha _{min}\alpha (\zeta /\zeta _o+\mathrm{})=D_{\zeta /\zeta _o+\mathrm{}}`$. Also according to multifractality phenomena, a possible analogy with thermodynamics can be established by relating $`\tau `$ to $`f(\alpha )(\zeta /\zeta _o)\alpha (\zeta /\zeta _o)\tau (\zeta /\zeta )`$ via a Legendre transformation. From this analogy, where $`f`$ is athe analogous ”entropy”, our analytical expression for the analogous ”specific heat” of the economic system becames $$C_\zeta \frac{^2\tau }{(\zeta /\zeta _o)^2}\frac{\delta \zeta _0\varphi _o}{2\sqrt{2}}sech^2(\zeta /\zeta _0).$$ (28) ## V Discussion Let us see next how the type of behaviour given in Fig.1 influences the ”analogous” specific heat function $`C_q`$ defined in Eq.(3) and how it might characterize the onset of a crash for a real stock market index. In Fig.2 we plot the ”analogous” specific heat $`C_q`$ of the S&P500 index for four different time lags $`T=1,10,30,80`$. The $`+++`$ curve is for the 1984-1988 data, $`xxx`$ is for 1982-1990 and $`ooo`$ is for 1980-1992. For time lags $`T80`$, we find that the main peak of our numerical $`C_q`$ resembles a classical (first-order) physics phase transition at a critical point given by the main peak position. The peak turns symmetric around the value $`q=1`$. Surprisely, this ”analogous” specific heat $`C_q`$ of the S&P500 index also displays a shoulder to the right of the main peak as a function of smaller time lags. Clearly, on decreasing $`T`$, the presence of the shoulder is a consequence of the large, temporal $`x(t+T)x(t)`$ fluctuations in this regime. We note that such peculiar behaviour for a double peaked specific heat function is known to appear in the Hubbard model within the weak-to-strong coupling regime . The relation in Eq.(2) requires $`T\mathrm{}`$ where the shoulder tends to vanish. It is this feature that make us believe that a large crash for an stock market index can be characterized by an ”analogous” specific heat which resembles a classical phase transition at a critical point as studied in multifractal physics. We now turn to the results of our theoretical economic approach. For $`T=80`$, the full line in Fig.2 represents theoretical results for $`C_\zeta `$ from Eq.(28) by choosing $`\zeta q+1`$ to fit the main peak position. It can be seen that the theoretical $`C_\zeta `$ curve resembles the phase transition features of the real S&P500 economic data. Our simple model for one commodity shade light into the main observed features regarding a possible analogous phase transition that occurs when the excess demand becomes non-linear (c.f., cubic $`p`$-term in Eq.(12)) in terms of the price for one commodity. We believe the width difference between both curves, i.e. the analytical $`C_\zeta `$ and the estimated $`C(q)`$ function from historical S&P500 data, is due to the fact that the S&P500 price index is made of large-capitalization stocks representing a ”basket” or portfolio of commodities. For large time lags, there is a sharp peak that resembles the quantitative signals measured in multifractals . From this feature we presume the existence of an analogous, say, critical point $`\zeta ^{}`$ above which inflated prices for one strategic commodity might be found. The maximum and minimum values of $`\alpha `$ in Eq.(27) (for more details see also ) allows for the existence of a critical point $`\zeta ^{}`$ above which the infinite hierarchy of phases can be found, but below which a single phase appears characterized by $`\alpha _{max}`$. It resembles a classical phase transition at a critical point . Of course, the analogy between multifractality and a thermodynamic phase transition as discussed here does not imply that the economic system has a phase transition. What we have shown, as a direct consequence of the $`p^3`$-term in Eq.(12), is that prices can became inflated prior to equilibrium (i.e., $`t<0`$ by convention), whereas after a sudden crash prices might devalue. The greater $`\|\delta \|`$-values are taken, the smaller the price reduction becomes. If $`p>0`$ the amount supplied exceeds demand and stocks accumulate whereas if $`p<0`$ the stocks deplete. Also by tuning $`\delta 0`$ in Eq.(13), we are able to predict that prices can decrease (or increase if $`A_2>0`$) monotonically with $`C0`$ for all $`\zeta `$. It is a well-known fact the hill behaviour of the generalized dimensions $`D_q`$ for $`q<0`$ when using the box-counting method as in the present work . Thus, we check if $`D_q`$ is sufficiently smooth for $`C_q`$ to be meaningful. In Fig.3 we represent $`D_q`$ of the S&P500 index for the 1980-1992 period within the range $`q0`$ for the time lags $`T=1,10,30,80,120`$. The total number of sequence points analysed include 1495 points for the different $`C_q`$ curves, 1500 points for $`D_q`$ and 9864 points for the S&P500 1980-1992 data set. We also fit this non-seasonal data by standard exponential smoothing techniques, and from the fitting our error estimates for the $`D_q`$ curves is found to be less than $`3\%`$. Our results for $`D_q`$ at negative $`q`$ (not shown) follow a typical convergent behaviour as can be seen, for example, in . From Fig.3 it can be seen that, when increasing $`T>30`$, $`D_q`$ is fully multifractal-like and for $`T>120`$ it becomes flatter (i.e., uniform measure) and tends to one independently of $`q`$, so multifractality becomes smaller. For lower values of $`T<30`$ we find a non-monotonous decreasing behaviour of $`D_q`$, conceivable within the double peaked form of $`C_q`$ displayed in Fig.2, which relates to the presence of the onset of crash for the S&P500 stock index in Fig.1. Using this data set, we also estimated the multifractality strength of the time sequence by considering the limit $`1D_q\mathrm{}`$ for the few different time lags displayed in Fig.3. We find that that this quantity does not follow a power-law scaling for $`T`$ as in citeRom97 for values $`T1`$. This is a consequence of the complex network of trading interactions comprised in our non-linear forms for the supply and demand functions. The present choice for a excess demand function of the form $`E(p)(pp^{})\frac{\delta ^2}{2}(pp^{})^3`$, with $`\delta 0`$ plays a key role in our description. We have $`d^{}=q^{}`$, hence the second price derivative $`\frac{d^2E}{dp^2}3\beta _o\delta ^2p`$ (or that of $`D`$ and $`Q`$) is price dependent. It is such a behaviour of $`E`$ (independently of the sign of $`\delta `$), that leads to obtain an abrupt fluctuation in the price dynamics and characterize multifractality phenomena. Our expressions for the demand and supply functions of Eq.(5) are justified as follows. As seen in Fig.4, the (commonly used) linear $`p`$-dependence for $`D`$ and $`Q`$ and our assumed non-linear form for the these functions display similar behaviour when $`\|\delta p\|<<1`$. Even more important, this figure depicts the fact that as price falls, the quantity demanded for a commodity can increase in agreement with one of the basic principles of economy. Since the demand curve can indeed change in a number of ways which may not be at all obvious, similar tails to our $`D`$ and $`Q`$ functions has also been previously hypothesized in . In real world, exceptions to the general law of demand can take place making the $`D`$ curve to increase upwards from low- to high-prices. However these exceptions, which include goods of conspicuous consumption -as, e.g., certain articles of jewelry- are not very important . Theoretically, this simply would mean to set $`d_o>0`$ in our $`D(p)`$ function independently of the order parameter $`\delta `$. The upward dependence of $`D`$ can occur as soon as there is speculation (as in precious stones mentioned but also for market prices). The idea is that when price increases, investors buy, because they hope the price will keep climbing. This is called a ”trend-following” investment strategy. On the other hand, our choice for $`Q`$ (with $`q_o>0`$) also follows the typical behaviour observed in a competitive market (where no individual producer can set his own desired price). That is, the higher the price, the higher the profit, then the higher the supply. It is also important to mention that the present choice for $`D`$ and $`Q`$, and their linear $`p`$-dependence both lead to the same linear relation for $`D`$ vs. $`Q`$, namely $`(Dd^{})/d_o=(Qq^{})/q_o`$. Such a linear behaviour founds frequent use in applied economics. Furthermore, it is well known in econometrics that demand functions are somewhat abstract quantities since all of these data are taken to refer to possible events at just one moment of time . In particular, consumer behaviour (i.e. tastes, desires, $`\mathrm{}`$) can shift a (concave up or concave down) $`D`$ curve (which may resemble the tails of our $`D`$ and $`Q`$ functions). Also, typical demand and supply of capital look like a step function underlying forces toward interest rates . Usually additional hypothetical information is needed to make up realistic demand relationships (e.g., consumer interviews). In principle, one might also determinate $`D`$ out of many data sources from different economic sectors using standard statistical techniques for multiple regression time series analysis. But the effectiveness of such an approach varies case by case. In view of all of this, our non-linear approach for $`D`$ and $`Q`$ may well be placed for simulating market situations leading to a sudden decline around equilibrium for the price of a commodity. To this end, we add that a commodity price function displaying an inflection point at a characteristic frequency has also been theoretically discussed in . Another important test for our $`D`$ and $`Q`$ expressions arises by considering possible aggregated changes in conditions of demand (or, supply). Because of such aggregated changes, e.g. due to buyers’ income and scale of preferences, it is always difficult to estimate how much of the change in $`D`$ is due to price alone. To know what these effects are upon $`D`$, and obtain the degree of responsiveness of demand to price variations, it is necessary to study the point price elasticity of demand curves, defined as $`\epsilon d\mathrm{log}D/d\mathrm{log}p`$ . This quantity is said to be elastic (or flat) if $`\epsilon >1`$ or inelastic if $`\epsilon <1`$ depending on the demand schedule. By assuming a linearly decreasing $`p`$-dependence for $`D`$ one obtains the expression $`\epsilon =\frac{1}{1+(d^{}/d_op)}`$. And if $`\delta 0`$ one gets $`\epsilon \frac{13(\delta p)^2/2}{[1(\delta p)^2/2]+(d^{}/d_op)}`$. Hence we find that a transition from an inelastic to an elastic demand curve appears at $`p\frac{d^{}/d_o}{1(\delta p/2)^2}`$ Whereas by considering the common linear dependence of $`D(p)`$ one finds $`pd^{}/d_o`$. Therefore, both approaches lead similar results for the elasticity of demand (or supply) if $`\delta 0`$. ## VI Concluding Remarks We have found that within the framework of multifractal physics, the ”analogous” specific heat of the S&P500 discrete price index displays a shoulder to the right of the main peak for low values of time lags. On decreasing $`T`$, the presence of the shoulder is a consequence of the peaked, temporal $`x(t+T)x(t)`$ fluctuations in this regime. For large time lags ($`T>80`$), we have found that $`C_q`$ displays typical features of a classical phase transition at a critical point according to multifractal physics. Our simple continuos model for one commodity mimics the main observed features of $`C(q)`$ for large trading time lags. We believe the width difference between the analytical $`C_\zeta `$ and estimated $`C(q)`$ curves from historical S&P500 data is due to the fact that the S&P500 index comprises many commodities. From these results we conclude that an analogous phase transition might occurs when the excess demand becomes non-linear (c.f., cubic $`p`$-term in Eq.(12)) in terms of the commodity price. We have assumed that there is only one stock in which a commodity can be stored. The market has been considered competitive so it self-organizes to determine the behaviour of prices. All factors determining $`D`$ and $`Q`$ other than $`p`$ are assumed to remain constant over time. There exists a price adjustment relation that takes into account deviations of the stock level $`S`$ above certain optimal level $`S_o`$ characterized by a noisy $`\lambda `$ parameter. We have postulated simple non-linear forms for the quantities $`D`$ demanded and $`Q`$ supplied, with $`\delta `$ the order parameter, and have neglected higher order terms in their expansion on $`pp^{}`$. We have followed the context of other simple economic models and assumed that the optimal stock level $`S_o`$ depends linearly on the demand (with $`\mathrm{}`$ being the slope). We have considered $`\mathrm{}`$ to satisfy the constrain in Eq.(11), relating the economic model variables: $`\gamma `$, $`\lambda `$, $`q_o`$ and $`d_o`$. We have $`d^{}=q^{}`$ and set the equilibrium price $`p^{}`$ to 0. We have identified the behaviour of $`p(t)`$ with an analogous phase transition as seen in multifractals, by considering $`t`$ to be the relevant variable and using a simple continuous probability theory. We have related $`t\sqrt{\lambda \beta _o/2}`$ to a continuous random variable and have related its probability distribution function to our solutions for $`p(t)`$ given in Eq.(13). Our definitions for the analogous thermodynamic variables have been done according to definitions used in multifractal physics and by mapping to multifractality phenomena in random multiplicative processes. Of course the scenario of a transition from, say, inflated to devalued price changes in the time domain is pure speculation. However a great deal of relevant information has been extracted from the present continuous approach which, essentially, does relay on $`\delta `$ only. Our description presumes the existence of a stationary probability function $`𝐏`$ as in Eq.(21). We have assumed $`𝐏`$ to be proportional to the continuous function $`\tau `$ (identified as a free energy) similarly to the dielectric breakdown model or the Poisson fractal growth models . We add that this type of approximation is also used when modeling earthquakes, where the probability function of the total number of relaxations (size) of the earthquakes is set proportional to the energy release during an earthquake . Within the framework of general probability theory (see, e.g., ), a continuous arbitrary random variable can have an associated probability distribution (if it is discrete) or a positive probability density (if it is continuously distributed). The later is required to be differentiable. These can then be related by an integral equation of the type used in this work to then derive the analogous thermodynamics equations. To do this we have proceed as follows. We have used the characterization of multifractal singularities -where an analogy with thermodynamics has been established in the literature- to derive, and associate a meaning to our analogous quantities. The main point to understand is that we have not carried out any direct numeration or scaling of particular fractal configurations or sets, but we have evaluated the integral of the probability distribution $`𝐏`$ (for the continuous random variable $`\zeta `$) assumed to be related with our $`p(t)`$ function and mapped the results for $`\tau `$ to those of a random multiplicative process. Our equivalent definitions for the thermodynamic variables, derived by solving such integral, followed from the convention used in multifractal physics in the sense that $`\tau `$ must not be a linear function of $`\zeta `$ (see, e.g., ). Alternative approaches can also be found in . In all theses cases, analogous thermodynamic quantities follow from the Legendre transform of $`\tau (\zeta )`$. As discussed in , the concept of a phase transition in multifractal spectra was first found in the study of logistic maps, julia sets and other simple systems. Evidence was then found for a phase transition in more complex random systems such as diffusion limited aggregation. The condition for an analogy between multifractality and thermodynamic phase transition is that the analogous free energy ($`\alpha `$ in our notation) undergoes a quite sharp jump near a critical value $`\zeta ^{}`$. For values of $`\zeta <\zeta ^{}`$, the analogous free energy $`\tau `$ is dominated by the maximum energy term $`\alpha _{max}`$ and a singular behaviour of the specific heat is found at this point. A well-defined analogous entropy function f($`\alpha `$) also suggests the existence of an analogous phase transition. We have shown that the present kinetic description satisfies such properties obeying the constrains for $`\mathrm{}`$ in Eq.(11). This condition simplifies the analysis and corresponds to the case in which the term for the first derivative of $`p`$ with respect to $`t`$ (i.e., $`\frac{dp}{dt}`$) is absent in the price adjustment Eq.(12). Application of multifractal analysis to discrete 1D time sequences as derived, for example, from the cellular automaton model of a rice model is no new (see, e.g., ). Nevertheless, to our best knowledge, we have related multifractal physics to financial time series for the first time. Our main contribution has been the analysis of the ”analogous” specific heat (or second derivatives) $`C_q`$ of the data sequence, in conjunction with the analytical form derived from our proposed one-stock model, which suggest typical features of a classical physics phase transition at a critical point. The double peaked form of $`C_q`$ is a consequence of the presence of the onset of crash for the S&P500 stock index. Our work also differs from previous formalisms of multifractality of time series in that we have analytically characterized multifractal singularities and its thermodynamics interpretation. The suggested non-linear analytical forms for the supply and demand functions of a commodity lead us to derive theoretically the observed features of a classical phase transition using the simplest economic model. We believe all these novel aspects of the topic could stimulate further investigations on this direction and can be important to open beneficial discussions in the field of econophysics. ### Acknowledgments The author gratefully acknowledges discussions via e-mail with Prof. D. Sornette.
warning/0004/cond-mat0004272.html	ar5iv	text	# Spin chirality induced by the Dzyaloshinskii-Moriya interaction and the polarized neutron scattering ## Abstract We discuss the influence of the Dzyaloshinskii-Moriya (DM) interaction in the Heizenberg spin chain model for the observables in the polarized neutron scattering experiments. We show that different choices of the parameters of DM interaction may leave the spectrum of the problem unchanged, while the observable spin-spin correlation functions may differ qualitatively. Particularly, for the uniform DM interaction one has the incommensurate fluctuations and polarization-dependent neutron scattering in the paramagnetic phase. We sketch the possible generalization of our treatment to higher dimensions. Since the works by Dzyaloshinskii and Moriya , the antisymmetric spin exchange interaction plays an important role in the physics of condensed matter. Being introduced for the explanation of the weak ferromagnetism in antiferromagnets without center of inversion, the Dzyaloshinskii-Moriya (DM) interaction is found nowadays in various problems of magnetism and statistical physics. Being the relativistic effect, the magnitude of DM interaction, $`D`$, is generally expected to be small in comparison with the usual symmetric superexchange, $`J`$. In some compounds, however, this interaction can attain a sizeable value. For instance, one has $`D/J=0.18`$ in the hexagonal perovskite CsCuCl<sub>3</sub> and $`D/J0.05`$ in copper benzoate. Remarkably, the DM interaction in the latter compounds takes place in the quasi-one-dimensional spin subsystems. On this reason, we are primarily concerned below with the one-dimensional (1D) situation of a quantum spin chain. Generally, the DM interaction between two spins, $`𝐒_{1,2}`$, is written as $`𝐃(𝐒_1\times 𝐒_2)`$ with an axial DM vector $`𝐃`$. In a chain, $`𝐃`$ may spatially vary both in direction and magnitude, however, the symmetry arguments usually rule out most of the possibilities and confine the theoretical discussion to two principal cases. The first one is the uniform DM interaction, $`𝐃=const`$ over the system. The second case is the staggered DM interaction, with antiparallel $`𝐃`$ on adjacent bonds. Among the other studies, we should mention the discussion of the $`XY`$ spin chains with randomly distributed values of $`D`$ and the growth models with imaginary uniform $`D=i\lambda `$ leading to non-Hermitian Hamiltonian. A model of $`XY`$ spin chain with a ternary DM interaction was introduced and solved recently. In a present paper, we deal mostly with two above cases, uniform and staggered DM interaction. We consider also a model of a non-ideal lattice, where one finds, say, an almost uniform situation with one possible DM value, $`𝐃`$, taking place on a chain fragment of average length $`l_1`$, and another value, $`𝐃`$, on a chain fragment of length $`l_2`$, while $`l_{1,2}1`$. The situation is then described in probabilistic terms. We show that, being the situation uniform, staggered or random, the spectrum of the Heisenberg chain with the DM interaction is equivalent to one of XXZ spin model and is computed exactly for spin 1/2. Of course the observable susceptibilities can differ crucially, as we demonstrate below. Therefore we extend the previous result by Alcaraz and Wreszinski, that the Heisenberg 1D model with the uniform DM interaction is reduced to XXZ spin exchange model and is exactly solvable. For the uniform and almost uniform DM interaction, we show that the observable spin-spin correlation function possess an incommensurate structure. This incommensurability phenomenon was noted previously for the $`XY`$ spin chain model and uniform antisymmetric spin interaction. We show further that in this case the spin susceptibility tensor $`\chi ^{\alpha \beta }`$ acquires an antisymmetric part. This leads to the appearance of the polarization-dependent part of the neutron scattering cross-section, which makes possible the direct observation of the direction and the value of the DM vector $`𝐃`$. Our results are applicable in the absence of the long-range magnetic order in the system. They also can be generalized towards higher-dimensional situation, as discussed below. Therefore our treatment may provide an explanation of the earlier experiments in cubic ferromagnet MnSi, where the DM-induced incommensurability of the magnetic fluctuations and polarization dependence of the neutron scattering were observed both below and above the Curie ordering temperature. We consider the spin chain Hamiltonian of the form $``$ $`=`$ $`{\displaystyle \underset{l=1}{\overset{L}{}}}(J𝐒_l𝐒_{l+1}+𝐃_l[𝐒_l\times 𝐒_{l+1}])`$ (1) with AF Heisenberg coupling $`J`$ and the Dzyaloshinskii-Moriya term $`𝐃_l`$. We choose the vector $`𝐃_l`$ to be directed along the $`z`$axis. We observe that $``$ is simplified upon a canonical transformation $`e^{iU}e^{iU}`$ with $$U=\underset{l=1}{\overset{L}{}}\alpha _lS_l^z,\alpha _l=\underset{i=1}{\overset{l1}{}}\mathrm{tan}^1(D_i/J),$$ (2) and $`\alpha _1=0`$. Note that this transformation works for all values of $`S`$. The periodic boundary conditions (BC) require $`\alpha _{L+1}=0\mathrm{mod}\mathrm{\hspace{0.17em}2}\pi `$ which relation is not generally satisfied. However, in the thermodynamic limit $`L\mathrm{}`$ the influence of BC can be negliged. Introducing $`S_j^\pm =S_j^x\pm iS_j^y`$, one can easily see that $`\stackrel{~}{S}_l^\pm `$ $``$ $`e^{iU}S_l^\pm e^{iU}=S_l^\pm e^{i\alpha _l},\stackrel{~}{S}_l^z=S_l^z,`$ (3) and our choice of the coefficients $`\alpha _l`$ removes the antisymmetric part of the Hamiltonian, $`[\stackrel{~}{𝐒}_l\times \stackrel{~}{𝐒}_{l+1}]`$. We consider below two principal possibilities: the “uniform” situation $`D_l=J\mathrm{tan}\delta `$ and the “staggered” one $`D_l=(1)^lJ\mathrm{tan}\delta `$. In both cases the Hamiltonian is reduced to the $`XXZ`$ model : $``$ $`=`$ $`{\displaystyle \underset{l=1}{\overset{L}{}}}(J^x(\stackrel{~}{S}_l^x\stackrel{~}{S}_{l+1}^x+\stackrel{~}{S}_l^y\stackrel{~}{S}_{l+1}^y)+J\stackrel{~}{S}_l^z\stackrel{~}{S}_{l+1}^z)`$ (4) with $`J^x=\sqrt{J^2+D_l^2}`$ independent of $`l`$. Our subsequent discussion is based on the observation, that the DM interaction results in two effects for the observable susceptibility of the system. First effect is the modification of the spectrum, as seen in the equivalent Hamiltonian (4). The appearance of the ”easy-plane” anisotropy ($`J^x>J`$), however, does not lead to a gap in the spectrum. The exact solution of (4) for $`S=1/2`$ shows that the correlation functions $`\stackrel{~}{S}_l^x\stackrel{~}{S}_m^x\|lm\|^\nu `$ and $`\stackrel{~}{S}_l^z\stackrel{~}{S}_m^z\|lm\|^{1/\nu }`$ with $`\nu =1\|\delta \|/\pi `$. Since the value of DM exchange is expected to be small, $`\|D_l\|J`$, the long-distance decay of the above correlation functions is described to a good accuracy by the “isotropic” Heisenberg situation, with $`\nu =1`$. The second effect of DM interaction for the observables is the explicit dependence of the relation (3) between the new and old spin variables on the values of $`D_l`$. We focus our attention below on the latter effect, which leads to the qualitative changes in the experimentally observable susceptibilities. The two-time Green’s function for the operators $`A`$ and $`B`$ is defined as $`\chi _{AB}(t)=i\theta (t)[A(t),B]`$ where $`[\mathrm{},\mathrm{}]`$ stands for a commutator and $`\theta (t)=1`$ at $`t>0`$. Upon the ”twist” $`e^{iU}`$ the $`z`$-component of spin operators remains unchanged, and one has for the longitudinal $`zz`$ susceptibility $`\chi _{lm}^{zz}(t)=i\theta (t)[S_l^z(t),S_m^z]=i\theta (t)[\stackrel{~}{S}_l^z(t),\stackrel{~}{S}_m^z]𝒢_{lm}^{}(t)`$. Therefore the observable $`\chi _{lm}^{zz}(t)`$ has a commensurate antiferromagnetic modulation. The expressions for the transverse spin susceptibility are more complicated. It is convenient to introduce the matrix $$\chi _{lm}^{}(t)=i\theta (t)\left[\begin{array}{cc}[S_l^x(t),S_m^x],& [S_l^x(t),S_m^y]\\ [S_l^y(t),S_m^x],& [S_l^y(t),S_m^y]\end{array}\right],$$ (5) in the initial system (1). In the simpler ”twisted” system (4) we have $`i\theta (t)[\stackrel{~}{S}_l^x(t),\stackrel{~}{S}_m^x]=i\theta (t)[\stackrel{~}{S}_l^y(t),\stackrel{~}{S}_m^y]𝒢_{lm}^{}(t)`$ and $`[\stackrel{~}{S}_l^x(t),\stackrel{~}{S}_m^y]=[\stackrel{~}{S}_l^y(t),\stackrel{~}{S}_m^x]=0`$. Returning back to quantities $`[S_l^\alpha (t),S_m^\beta ]`$ with the use of (3), we get $$\chi _{lm}^{}(t)=𝒢_{lm}^{}(t)\left[\begin{array}{cc}\mathrm{cos}\alpha _{l,m},& \mathrm{sin}\alpha _{l,m}\\ \mathrm{sin}\alpha _{l,m},& \mathrm{cos}\alpha _{l,m}\end{array}\right]$$ (6) with $`\alpha _{l,m}=\alpha _l\alpha _m`$. For later comparison, it is worth to consider first the case of the staggered DM interaction. We have $`D_l=(1)^lJ\mathrm{tan}\delta `$ and $`\alpha _{l,m}=((1)^l(1)^m)\delta /2`$. In this case we write $`\mathrm{cos}\alpha _{l,m}=\mathrm{cos}^2(\delta /2)+(1)^{lm}\mathrm{sin}^2(\delta /2)`$ and $`\mathrm{sin}\alpha _{l,m}=((1)^l(1)^m)(\mathrm{sin}\delta )/2`$. Clearly, the off-diagonal components $`\chi _{lm}^{xy}(t)`$ of the matrix (6) do not depend on the difference $`(lm)`$ only and the two-momenta Fourier transform $`A(q,q^{},\omega )=𝑑t_{lm}e^{iqliq^{}mi\omega t}A_{lm}(t)`$ should be introduced. Then we obtain the off-diagonal components in the form $`\chi ^{xy}(q,q^{},\omega )=(𝒢^{}(q,\omega )𝒢^{}(q\pi ,\omega ))_\tau \delta (qq^{}\pi +\tau )`$ where $`_\tau `$ stands for the sum over all vectors $`\tau =2\pi n`$ of the reciprocal lattice. However, apparently in all physical observables one finds the symmetrized form of the susceptibility ($`q=q^{}`$) and the off-diagonal terms in the matrix $`\chi ^{}`$ vanish. Therefore in the case of staggered DM interaction one is left with the diagonal component of the matrix $`\chi ^{}(q,\omega )`$ of the form : $`\chi ^{}(q,\omega )`$ $`=`$ $`𝒢^{}(q,\omega )\mathrm{cos}^2(\delta /2)`$ (8) $`+{\displaystyle \frac{1}{2}}[𝒢^{}(q\pi ,\omega )+𝒢^{}(q+\pi ,\omega )]\mathrm{sin}^2(\delta /2)`$ We see that the regions of the AF and the ferromagnetic fluctuations are mixed in the observable susceptibility. A consequence of this feature is the anomalous temperature behavior of the uniform static susceptibility $`\chi ^{xx}(0,0)=\chi ^{yy}(0,0)`$ for the AF chain (see also the discussion after Eq.(5.3) in the original Moriya’s paper ). It is known that in the Heisenberg $`S=1/2`$ chain one has $`𝒢^{}(0,0)J^1`$ and $`𝒢^{}(\pi ,0)T^{2+\nu }`$, therefore $`\chi ^{zz}(0,0)J^1`$ and $`\chi ^{xx}(0,0)`$ $`=`$ $`\chi ^{yy}(0,0)J^1[const+\delta ^2(J/T)^{1+\|\delta \|/\pi }]`$ (9) The Eq. (9) has a simple physical meaning. Indeed, in the considered case the operator $`U`$ “cants” the local coordinate frames by an angle $`\pm \delta /2`$. It leads effectively to the non-compensated spin $`\mathrm{\Delta }S=\delta /4`$ in the $`xy`$ plane. Being the spins $`\mathrm{\Delta }S`$ free, it would then lead to the Curie law for the susceptibility $`\chi \mathrm{\Delta }S^2/T`$. The 1D character of the interacting spin system results in the nontrivial exponent in the $`T`$dependence of this term (cf. also ). At the same time, the above transfer of the spectral weight, Eq. 8, is apparently negligible to be observed in the neutron scattering experiments. On the other hand, for the “uniform” DM interaction $`D_l=D`$ we have $`\alpha _{l,m}=(lm)\delta `$. It results in the incommensurability of the transverse spin correlations. Fourier transforming Eq. (6), we obtain $`\chi ^{}(q,\omega )`$ $`=`$ $`{\displaystyle \frac{1}{2}}𝒢^{}(q+\delta ,\omega )\left[\begin{array}{cc}1,& i\\ i,& 1\end{array}\right]`$ (15) $`+{\displaystyle \frac{1}{2}}𝒢^{}(q\delta ,\omega )\left[\begin{array}{cc}1,& i\\ i,& 1\end{array}\right]`$ Let us discuss the physical consequences of this expression. Evidently, $`\chi ^{xx}(0,0)const`$ in this case and the presence of the “uniform” DM interaction is not revealed by the measurements of the temperature dependence of the uniform static susceptibility. Much more interesting are the implications of (15) for the polarized neutron scattering experiments. The basic quantity here is neutron scattering cross-section, which is connected to the Green’s function $`\chi ^{\alpha \beta }(q,\omega )`$ of the spin system. It is convenient to write $`\chi ^{\alpha \beta }=\chi _S^{\alpha \beta }i\chi _A^{\alpha \beta }`$, with the symmetric and antisymmetric tensors, $`\chi _S^{\alpha \beta }`$ and $`\chi _A^{\alpha \beta }`$, respectively. Up to fundamental constants, we have : $`{\displaystyle \frac{d^2\sigma (q,\omega )}{d\mathrm{\Omega }d\omega }}`$ $``$ $`N(\omega )[Im\chi _S^{\alpha \beta }(q,\omega )(\delta ^{\alpha \beta }\widehat{Q}^\alpha \widehat{Q}^\beta )`$ (17) $`+Im\chi _A^{\alpha \beta }(q,\omega )ϵ_{\alpha \beta \gamma }\widehat{Q}^\gamma (\widehat{𝐐}𝐏_0)],`$ where $`N(\omega )`$ is the Planck function, the unit vector $`\widehat{𝐐}=𝐐/Q`$ is directed along the neutron’s momentum transfer $`𝐐`$, $`ϵ_{\alpha \beta \gamma }`$ is totally antisymmetric tensor, $`𝐏_0`$ is the incident neutron’s polarization and $`q`$ is the on-chain projection of $`𝐐`$. From (17) we see that if the whole crystal is characterized by Dzyaloshinskii vector $`𝐃`$ (uniform situation, Eq. 15), then the polarization-dependent part of cross-section is non-zero and is given by $`{\displaystyle \frac{d^2\sigma _1}{d\mathrm{\Omega }d\omega }}`$ $``$ $`(𝐃\widehat{𝐐})(\widehat{𝐐}𝐏_0)Im{\displaystyle \frac{𝒢^{}(q+\delta ,\omega )𝒢^{}(q\delta ,\omega )}{2D}}.`$ (18) A certain subtlety should be discussed here. Under the parity transformation we have $`qq`$, $`𝐃𝐃`$, $`𝐏_0𝐏_0`$. At the first glance $`\delta \delta `$ and thus $`d^2\sigma _1/d\mathrm{\Omega }d\omega `$ changes the sign, as it should not be. An inspection of (2), (6) (15) shows however that the quantity $`\delta `$ in (18) appears as the differential of $`\alpha _l`$. The latter object is the sum of the phases $`\delta _j`$ over the bonds $`j`$ to the left of $`l`$. Hence we have under the parity transformation $`\alpha _l\alpha _l+const`$ and $`\delta \delta `$ in (18), which restores the desired property of the cross-section. The contribution of the symmetric part of $`\chi ^{\alpha \beta }`$ to $`d^2\sigma (q,\omega )/d\mathrm{\Omega }d\omega `$ in the considered case of uniform $`𝐃`$ is two-fold. One still has the commensurate fluctuations of spin components $`S^z𝐃`$, with a peak at the AF position. At the same time, the DM interaction splits the AF peak related to transverse fluctuations into two peaks of the weight $`1/2`$, Eq. 15. The relative weights of these two structures depend on the direction of $`𝐐`$ as follows $`{\displaystyle \frac{d^2\sigma _2(q,\omega )}{d\mathrm{\Omega }d\omega }}`$ $``$ $`(1+\widehat{Q}_z^2)Im{\displaystyle \frac{𝒢^{}(q+\delta ,\omega )+𝒢^{}(q\delta ,\omega )}{2}}`$ (20) $`+(1\widehat{Q}_z^2)Im𝒢^{}(q,\omega ).`$ Note that at low temperatures, the incommensurate peaks have more singular behavior according to our discussion after Eq. (4). An important thing to be stressed here is the following. It is known that the incommensurate long-range magnetic structures may arise due to the competing interactions in the spin system. In this case one expects that all three diagonal components of the spin susceptibility $`\chi ^{\alpha \alpha }`$ are peaked in the paramagnetic region at the same incommensurate wave-vector. The off-diagonal components of $`\chi ^{\alpha \beta }`$ are absent. This is fairly different from the picture described above, Eqs. 18, 20. Hence the experimental observation of the incommensurability phenomenon in the paramagnetic phase, accompanied by the polarization dependence of the neutron scattering cross-section could serve as an indication to the presence of the uniform DM interaction. Remarkably, the value and the direction of the pseudo-vector $`𝐃`$ can be, in principle, determined this way. In reality, however, the macroscopic sample is rarely uniform and it should be expected to split to domains with different directions of $`𝐃`$. To account for this situation, it is instructive to analyze a model where the value of the DM interaction $`D_l`$ takes randomly two values $`\pm J\mathrm{tan}\delta `$. Consider first the oversimplified case when $`D_l=0`$ and $`D_lD_m=0`$ for $`lm`$, here $`\mathrm{}`$ denotes averaging over the realizations. In this case the spectrum is still defined by Eq.(4), and $`\chi ^{zz}`$ is given by the above expression. At the same time, one can easily show that the averaged susceptibility $`\chi ^{}`$ has a diagonal form and exhibits an exponential decay of correlations : $$\chi _{lm}^{}(t)=𝒢_{lm}^{}(t)\mathrm{exp}(\|lm\|/l_{})$$ (21) with the correlation length $`l_{}=1/\mathrm{ln}(\mathrm{cos}\delta )\delta ^2`$. Now consider a more realistic situation when one still has $`\delta _j=\pm \delta `$, but the signs of $`\delta _j`$ on the adjacent bonds are correlated. The diagonal and off-diagonal parts of the matrix (6) are given, respectively, by the real and imaginary part of the average $$\mathrm{exp}i\alpha _{l,m}\underset{\{\delta _j\}}{}p(\delta _1,\mathrm{},\delta _L)\mathrm{exp}(i\underset{k=m}{\overset{l1}{}}\delta _k).$$ (22) We assume that the joint distribution function $`p(\delta _1,\mathrm{},\delta _n)`$ has a Markovian character, $`p(\delta _1,\mathrm{},\delta _n)=p(\delta _1,\mathrm{},\delta _{n1})\widehat{p}(\delta _n\|\delta _{n1})`$. In this physically important case we arrive at the dichotomous Markovian noise $`\delta _j`$ with a discrete “time” $`j`$. We set $`\delta _j=\delta d`$ which defines the on-site (“equilibrium”) probability as $`p_0=(\frac{1+d}{2},\frac{1d}{2})`$. The matrix $`\widehat{p}(\delta _n\|\delta _{n1})`$ satisfies the “conservation laws” for the total and equilibrium probabilities, $`(1,1)\widehat{p}(\delta _n\|\delta _{n1})=(1,1)`$ and $`\widehat{p}(\delta _n\|\delta _{n1})p_0=p_0`$, respectively. These equalities fix $`\widehat{p}(\delta _n\|\delta _{n1})`$ in the form $`\widehat{p}=\left[\begin{array}{cc}1x(1d),& x(1+d)\\ x(1d),& 1x(1+d)\end{array}\right]`$ for all $`n`$, which corresponds to the following correlator on the adjacent sites : $`\delta _j\delta _{j+1}\delta _j^2=\delta ^2(1d^2)(12x)`$. The latter equalities mean that the absence of correlations corresponds to $`x=1/2`$ and the correlation lengths for positive and negative sequences of $`\delta _j`$ are $`1/l_{1,2}=x(1d)`$. Introducing the matrix $`𝒟=diag(e^{i\delta },e^{i\delta })`$, the quantity (22) is represented as a product $`(1,1)(𝒟\widehat{p})^{lm1}𝒟p_0`$, which is evaluated using the multiplication rules for the Pauli matrices. After straightforward, though tedious, calculation, we obtain the average (22) in general form, which is somewhat simplified in two principal cases of small $`\delta `$ : i) $`\delta x1`$ and ii) $`\delta x1`$. In the first case, keeping the terms of order of $`\delta ^2`$, we have $`e^{i\alpha _{m+n,m}}`$ $``$ $`\mathrm{exp}[n(i\delta d\delta ^2a_1)+\delta ^2a_2]`$ (23) $``$ $`\delta ^2a_2(12x)^n\mathrm{exp}[(n1)(i\delta d+\delta ^2a_1)]`$ (24) here $`a_1=(1x)(1d^2)/(2x)`$ and $`a_2=(1d^2)(12x)/(4x^2)`$. We see that the incommensurability wave vector is defined by the average on-bond value $`\delta d`$. Note that Eqs. (8), (21) are recovered at $`d=0`$, $`x=1`$ and $`d=0`$, $`x=1/2`$, respectively. When $`x\delta 1`$, we come to a more complicated situation. We have in this case $`e^{i\alpha _{m+n,m}}`$ $``$ $`e^{xn}\left[\mathrm{cosh}bn+{\displaystyle \frac{\mathrm{sinh}bn}{b}}(x+i\delta d)\right]`$ (25) with $`b=\sqrt{x^2+2ix\delta d\delta ^2}`$. We return to Eq. (15) at $`x=0`$ and $`d=1`$. It should be stressed that at $`d=0`$ one has $`Ime^{i\alpha _{m+n,m}}=0`$ both in Eqs. (23), (25) and in general case. It corresponds to the fact that the off-diagonal components of susceptibility $`\chi ^{}`$, Eq. (6), vanish in the system with zero average Dzyaloshinskii vector $`𝐃`$. As a result, the observable cross-section is polarization-independent. The position of maximum of the transverse spin fluctuations, though, may be incommensurate one for the almost uniform DM interaction, as seen from (25) at $`d=0`$ and $`x0`$. Now we discuss the possible generalization of our approach to a higher dimensional case. Consider a planar system with spins $`𝐒_{lm}`$ labeled by two indices. The interaction between spins takes place in two directions, and we write the corresponding quantities as $`J_{lm}^{(\alpha )}`$ and $`𝐃_{lm}^{(\alpha )}`$, with $`\alpha =x,y`$. For simplicity we consider the case when the vectors $`𝐃_{lm}^{(\alpha )}`$ lie along one direction, with possible variation in their magnitude. We introduce then two angles, $`\delta _{lm}^x=\mathrm{tan}^1(D_{lm}^{(x)}/J_{lm}^{(x)})`$ and $`\delta _{lm}^y=\mathrm{tan}^1(D_{lm}^{(y)}/J_{lm}^{(y)})`$. We are interested to arrive to a symmetrized Hamiltonian, similar to (4), by making the transformation with $`U=_{l,m=1}^L\alpha _{lm}S_{lm}^z`$. One can show that this transformation is possible if and only if $`\delta _{l,m+1}^x\delta _{lm}^x=\delta _{l+1,m}^y\delta _{lm}^y`$. In this case $`\alpha _{lm}`$ is uniquely defined and may be written in a form $`\alpha _{lm}=_{j=1}^{l1}\delta _{j,1}^x+_{j=1}^{m1}\delta _{l1,j}^y`$. The latter relations are not surprizing when we note that in the continuum limit they read as $`\times \stackrel{}{\delta }_𝐫=0`$, while $`\stackrel{}{\delta }_𝐫=(\delta _𝐫^x,\delta _𝐫^y)=\alpha _𝐫`$. In particular, the two-dimensional situation with $`D_{lm}^{(x)}=D`$ and $`D_{lm}^{(y)}=0`$ allows this transformation. The calculation of observables and the generalization to a three-dimensional case are done in the way similar to the above one. In conclusion, we discuss the observables in the spin system with the antisymmetric DM interaction. We show that in one spatial dimension the exactly found spectrum of such problem may coincide for different choices of the parameters of DM interaction. Despite this fact, the observable spin-spin correlation functions may differ crucially, and we discuss this feature with the application to the neutron scattering experiments. In particular, the incommensurability and the polarization dependence of the neutron scattering may be used for the determination of the value of the uniform DM interaction. We thank S.L. Ginzburg, D.R. Grempel, A. Gukasov, D. Petitgrand, V.P. Plakhty for useful discussions and comments. This work was supported by Russian State Program for Statistical Physics (Grant VIII-2), RFBR Grant No. 00-02-16873, and the Russian Program ”Neutron Studies of Condensed Matter”.
warning/0004/nucl-th0004025.html	ar5iv	text	# Statistical Model Description of 𝐾⁺ and 𝐾⁻ Production between 1 - 10 𝐴⋅GeV ## Abstract The excitation functions of $`K^+`$ and $`K^{}`$ mesons in heavy ion collisions are studied within a statistical model assuming chemical and thermal equilibrium with exact strangeness conservation. At low incident energies the associate production of kaons, i.e. the production of a $`K^+`$ together with a hyperon and the production of a $`K^{}`$ together with a $`K^+`$, implies specific features: different threshold energies and different dependences of $`K^+`$ and $`K^{}`$ yields on baryon number density. It is shown that the experimentally observed equality of the $`K^+`$ and $`K^{}`$ rates at energies $`\sqrt{s}\sqrt{s_{th}}0`$ is due to a crossing of the two excitation functions. Furthermore, the independence of the $`K^+`$ to $`K^{}`$ ratio on the number of participating nucleons observed at 1 and 10 $`A`$GeV is consistent with this model. , , Central heavy ion collisions at relativistic energies present an ideal tool to study nuclear matter at high densities and high temperatures. However, these collisions are complex and in order to interpret the results, two strategies are commonly used: (i) to describe the time evolution of the collisions using transport models and (ii) to use a statistical concept assuming thermal and chemical equilibrium and common freeze out parameters for all particles. In this Letter the second procedure, the statistical concept, will be followed. Of special interest here is the production of $`K^+`$ and $`K^{}`$ below and above the respective $`NN`$ thresholds. The experimental results have attracted much interest as the measured $`K^+`$ to $`K^{}`$ ratios in heavy ion collisions differ strongly from the ratios obtained in $`NN`$ reactions . These findings have lead to the proposal that in heavy ion collisions the “effective masses” of $`K^+`$ and $`K^{}`$ are changed as predicted for dense nuclear matter. The aim of this Letter is to discuss the $`K^+`$ and $`K^{}`$ production within a statistical model. This model describes the condition at freeze out using masses of free particles. The production of strange particles has to respect strangeness conservation. The attempts to describe the measured particle ratios including strange hadrons at AGS and SPS using a strangeness fugacity $`\lambda _S`$ is quite successful . However, the usual grand-canonical treatment is not sufficient, if the number of strange particles is small . This requires exact strangeness conservation which is done in the statistical model using the canonical formulation of strangeness conservation . Consequently, the abundance of $`K^+`$ mesons is suppressed since together with each $`K^+`$ also another strange particle, e.g. a $`\mathrm{\Lambda }`$ hyperon is produced via $`NNN\mathrm{\Lambda }K^+`$. And for $`K^{}`$ the corresponding channel is $`NNNNK^{}K^+`$. While the pion multiplicity per $`A_{part}`$ is approximately given by a simple Boltzmann factor (neglecting resonance contributions and isospin asymmetry), $$\frac{M_\pi }{A_{part}}\mathrm{exp}\left(\frac{E_\pi }{T}\right),$$ (1) the multiplicity of positively charged kaons is given by $$\frac{M_{K^+}}{A_{part}}\mathrm{exp}\left(\frac{E_{K^+}}{T}\right)\left[g_\mathrm{\Lambda }V\frac{d^3p}{(2\pi )^3}\mathrm{exp}\left(\frac{(E_\mathrm{\Lambda }\mu _B)}{T}\right)\right],$$ (2) with the temperature $`T`$, the baryo-chemical potential $`\mu _B`$, the degeneracy factors $`g_i`$, the volume $`V`$ (see ) and the energies $`E_i`$ of the particles $`i`$ and integrating over momentum $`p`$. The formula above, simplified for demonstration purpose, neglects higher order terms in $`V`$ , quantum statistics and other processes leading to the production of $`K^+`$. The corresponding formula for $`K^{}`$ production is similar, but does not depend on $`\mu _B`$, $$\frac{M_K^{}}{A_{part}}\mathrm{exp}\left(\frac{E_K^{}}{T}\right)\left[g_{K^+}V\frac{d^3p}{(2\pi )^3}\mathrm{exp}\left(\frac{E_{K^+}}{T}\right)\right].$$ (3) ¿From Eqs. (1) - (3) it is obvious that the exact strangeness conservation implies a reduction of $`K^+`$ and $`K^{}`$ yields as compared to the values calculated without exact strangeness conservation . In addition, since the volume in Eqs. (2) - (3) is proportional to the number of participants $`A_{part}`$, the $`K^+`$ and $`K^{}`$ multiplicities are expected to rise (for low $`T`$ and small $`V`$) quadratically with $`A_{part}`$ while $`M_\pi `$ increases linearly with $`A_{part}`$. These properties are in remarkable agreement with the experimental observations . The measured yields (or particle ratios) can be described in this statistical concept by lines in the $`T`$ and $`\mu _B`$ plane. All particle ratios measured around 1 $`A`$GeV (besides $`\eta /\pi _0`$) intersect within the experimental errors reflecting common values for $`T`$ and $`\mu _B`$ for all particles at freeze out . Surprisingly, even the measured $`K^+/K^{}`$ ratio fits into this picture and the calculated ratio does not depend on the choice of the volume $`V`$. However, the $`A_{part}`$ dependence enters in the ratios of strange to non-strange particles, e.g. in $`K^+/\pi ^+`$. Figure 1 shows the $`K^+/K^{}`$ ratios measured by the KaoS Collaboration , by the FRS Group and by the E866/E917 Collaboration at AGS as a function of $`\sqrt{s}`$. To obtain the theoretical results, shown as dashed line, we start from the universal freeze-out curve suggested in . Together with the measured systematics for the pion multiplicities, relations for $`T`$ and $`\mu _B`$ as functions of $`\sqrt{s}`$ are obtained . Within this approach, the $`K^+/K^{}`$ ratios are given as a dashed line in Fig. 1. The observed rise towards low incident energies reflects the fact that the two kaon species have different threshold energies due to their associate production. The $`K^+/K^{}`$ ratios measured in heavy ion collisions by the KaoS Collaboration show that the $`K^{}`$ yield compared to the $`K^+`$ cross section is much higher than expected from $`NN`$ collisions . This is especially evident, if the kaon multiplicities are plotted as a function of $`\sqrt{s}\sqrt{s_{th}}`$ where $`\sqrt{s_{th}}2m_N`$ is the energy needed to produce the corresponding particles taking into account the mass of the produced partners ($`\sqrt{s_{th}(K^+)}2m_N`$ = 0.67 GeV, $`\sqrt{s_{th}(K^{})}2m_N`$ = 0.987 GeV). The measured $`K^+`$ and $`K^{}`$ yields in heavy ion collisions are about equal for $`\sqrt{s}\sqrt{s_{th}}0`$ while the $`K^+`$ yield in $`NN`$ collision exceeds the $`K^{}`$ yields by a factor of 10 – 100 close to threshold. In Fig. 2 we show in the upper part the multiplicities of $`K^+`$ and $`K^{}`$ divided by $`A_{part}`$ as a function of $`\sqrt{s}\sqrt{s_{th}}`$ over a large energy range from SIS up to AGS. The full and dashed lines refer to the statistical model results for $`K^{}`$ and $`K^+`$ respectively. At values of $`\sqrt{s}\sqrt{s_{th}}`$ less than zero the two excitation functions cross. They differ at AGS energies by a factor of five which is in good agreement with the result for central collisions of Au+Au at 10.8 $`A`$GeV . The model calculations depend on the choice of the system, here Ni+Ni collisions. At SIS energies, only inclusive measurements for Ni+Ni are available. The values for $`K^+`$ are from Ref. . The results for $`K^{}`$ are from Ref. and corrected for the angular distribution . $`A_{part}`$ is chosen as $`A`$ which is based on estimates from the mean $`A_{part}`$ for $`K`$ production as kaons originate more from central collisions. At AGS energies the choice of the system has little influence which allows to plot the results for Au+Au collisions at 10.8 $`A`$GeV as well. This figure evidences that the similarity of the $`K^+`$ and $`K^{}`$ yield observed around 1 – 2 $`A`$GeV arises from the difference in the rise of the two excitation functions. This difference can be understood by the approximate formulae given in Eqs. (2) and (3). The density of $`K^+`$ contains the term $`E_\mathrm{\Lambda }\mu _B`$ while the $`K^{}`$ density has $`E_{K^+}`$ in the exponent. As these two values are different, the excitation functions, i.e. their variation with $`T`$, exhibit different slopes. Furthermore, Eqs. (2) and (3) evidence that for a low temperature $`T`$ and a small volume $`V`$ the dependence of the $`K^+`$ and $`K^{}`$ multiplicity on $`A_{part}`$ is quadratic which is in very good agreement with data . Small variations from the quadratic dependence can occur due to a change of $`T`$ and $`\mu _B`$ with $`A_{part}`$ . It is interesting to note that also hydrodynamical models predict a variation of $`K^+`$ and $`K^{}`$ with $`A_{part}^2`$ . Transport models, on the other hand, show an increase with $`A_{part}^\alpha `$ where $`\alpha `$ is approximately 1.4 - 1.6 . In these models the kaons are produced in multi-step processes which are more likely in central collisions where the density is higher. As already mentioned, the statistical model predicts that the variation of the $`K^+`$ and of the $`K^{}`$ yields with $`A_{part}`$ are equal. Hence, for a given collision the $`K^+/K^{}`$ ratio is expected not to vary with centrality. Indeed, this is in accordance with the data for Au+Au collisions at 10.2 $`A`$GeV . Figure 3 shows the results together with the prediction of the statistical model. It has even been observed at SIS energies for Ni+Ni collisions at 1.93 $`A`$GeV . This is remarkable as the $`K^+`$ production is above and the $`K^{}`$ production below their respective $`NN`$ thresholds. In summary, the statistical model using exact strangeness conservation is able to describe most of the measured particle ratios from SIS up to SPS energies. Within this framework the equality of $`K^+`$ and $`K^{}`$ multiplicities as a function of $`\sqrt{s}\sqrt{s_{th}}`$ is a consequence of two excitation functions with different slopes crossing at values $`\sqrt{s}\sqrt{s_{th}}`$ below zero. This model is also able to describe the dependence of the kaon yields on $`A_{part}`$ being quadratic around 1 $`A`$GeV. This effect fades away with increasing incident energy. The $`K^+/K^{}`$ ratio is predicted in the considered model as being independent of $`A_{part}`$ and this is, indeed, observed from SIS up to AGS energies. The statistical model presented here uses a unique freeze out for all particles. Detailed experimental studies on pion production show evidence for a time evolution of the pion emission with high-energy pions being emitted earlier . Such effects, however, are not visible on the level of total particle multiplicities since these involve integrals over the whole phase space. Despite the apparent success of the statistical model of particle production under the assumption of thermal and chemical equilibration and using masses of free particles, the present understanding of hadronic interactions contradicts chemical equilibrium for strange particles This discrepancy seem to put into question our present understanding of interactions at the high densities reached in heavy ion collisions. Indeed, already at and above twice nuclear matter densities, nucleons are hardly “free” individual particles. This “in-medium” effect clearly deserves further studies. K.R. acknowledges the partial support of the State Committee for Scientific Research (KBN).
warning/0004/hep-th0004102.html	ar5iv	text	# 1 Introduction ## 1 Introduction In the harmonic superspace approach to the study of correlation functions of gauge-invariant operators in four-dimensional superconformal field theories (with $`N=2`$ or $`N=4`$), the operators that can be investigated most easily are either products of hypermultiplets ($`N=2`$) or products of the $`N=4`$ field strength superfield $`W`$. It was originally conjectured in that the constraints of analyticity in the internal space (H-analyticity), generalised chirality (G-analyticity) and superconformal symmetry might be strong enough to determine correlation functions of such operators, at least for sufficiently low charges (the charge being proportional to the number of fields in the product). Although this has turned out not to be the case for generic values of the charges (see for a detailed discussion of $`N=2`$ four-point functions in this approach), it is true for certain special values of the charges. Correlation functions for these special values of the charges, called extremal correlators, were discussed in AdS supergravity in and , and the results found there, namely that these correlators are products of free propagators, were verified at one-loop order in perturbative field theory in and later to all orders, for four points, using harmonic superspace methods . It was further shown in that these simple results also hold for certain “next-to-extremal” values of the charges, and this has subsequently been confirmed by an AdS calculation . Related works investigate the existence of other non-renormalised SCFT correlators and the special structure of “near-extremal” correlation functions, where quantum corrections do occur . In this article we extend the results of to $`n`$ points for both extremal and next-to-extremal correlators. These results hold for arbitrary $`N=2`$ superconformal gauge theories. In we began by studying four-point correlators directly using the ideas mentioned in the first paragraph above, namely H-analyticity, G-analyticity and superconformal symmetry. However, it is in fact quicker to employ the reduction formula first discussed in the SCFT context in . This formula relates the derivative of an $`n`$-point correlator with respect to the coupling to an $`(n+1)`$-point correlator with one integrated insertion of the on-shell action; it was used in to give a proof of the non-renormalisation theorem for two- and three-point functions for arbitrary operators corresponding to Kaluza-Klein multiplets in AdS in $`N=4`$ SYM theory. Strictly speaking, this proof, and indeed the discussions in and in the current paper on extremal correlators, depends on the assumption that no contact terms need to be taken into account in the integrated insertion despite the fact that there is an integration which could involve coincident points. In principle rather special terms of this type could affect the derivation of these results. Some contact terms have been observed in perturbation theory , but these were of a type that do not affect the results we are interested in establishing. A search using superconformal symmetry failed to find any dangerous ones and there is no evidence that they exist from perturbation theory calculations . There is thus no indication so far that contact terms will affect the validity of the non-renormalisation theorems, despite fears to the contrary . In $`N=2`$ SYM theory the reduction formula, when applied to a correlation function of $`n`$ hypermultiplet composites, relates its derivative with respect to the coupling to an $`(n+1)`$-point correlator with an insertion of the chiral operator $`\mathrm{tr}(W)^2`$, where $`W`$ is the $`N=2`$ field strength superfield. The idea is then to show that this mixed $`(n+1)`$-point function vanishes when the charges of the original correlator are either extremal, which means that the charge of the first operator is equal to the sum of all of the others, or next-to-extremal, in which case the charge of the first operator equals the sum of all the others minus two. From this we conclude that the derivative of the $`n`$-point correlator with respect to the coupling vanishes, and so there can be no corrections to the lowest-order free-field expression. For the extremal case this expression is unique while there are several possibilities for the next-to-extremal case. We shall begin with a discussion of $`N=2`$ theories in the harmonic superspace formalism of GIKOS . This is followed by a discussion of the same topic in coordinate form, using the formalism of . ## 2 Extremal and next-to-extremal $`n`$-point correlators in $`N=2`$ harmonic superspace We shall be interested in gauge-invariant operators constructed from the hypermultiplets of the theory and the Yang-Mills field strength multiplet. Each hypermultiplet has charge 1 under the $`U(1)`$ isotropy group of $`SU(2)`$ which defines the sphere which is adjoined to ordinary superspace to form harmonic superspace. The hypermultiplet composites are analytic with respect to the sphere (H-analytic) and with respect to half of the spinorial covariant derivatives of superspace (G-analytic). We shall denote an $`n`$-point function of such operators by $`<p_1\mathrm{}p_n>`$. The Yang-Mills field strength superfield $`W`$ is chiral and we shall use only the operator $`\mathrm{tr}(W)^2`$ which occurs in the Yang-Mills part of the action. We shall denote an $`(n+1)`$-point function with $`n`$ hypermultiplets and one insertion of $`\mathrm{tr}(W)^2`$ by $`<0p_1\mathrm{}p_n>`$. In the extremal case the charges satisfy $`p_1=p_2+p_3+\mathrm{}p_n`$, while in the next-to-extremal case we have $`p_1=p_2+p_3+\mathrm{}p_n2`$. The reduction formula (see for a derivation in $`N=2`$ superspace) states that $$\frac{}{\tau }<p_1\mathrm{}p_n>d^4x_0d^4\theta _0<0p_1\mathrm{}p_n>$$ (1) where $`(x_0,\theta _0)`$ are the coordinates at the chiral point and $`\tau `$ is the usual complex Yang-Mills coupling constant. The points on the LHS are taken to be non-coincident and, as we remarked in the introduction, it is assumed that there are no contact contributions to the integral on the RHS. The function $`<0p_1\mathrm{}p_n>`$ is chiral at point 0 and G-analytic at points $`1,\mathrm{},n`$, it has the corresponding superconformal properties and carries positive charges $`p_1,p_2,\mathrm{},p_n`$ at points $`1,2,\mathrm{},n`$; it is also H-analytic, $$D_r^{++}0p_1p_2\mathrm{}p_n=0,r=1,\mathrm{},n\text{if point 0 }\mathrm{}\text{ point }n.$$ (2) In addition, it has the $`R`$-weight of 4 left-handed $`\theta `$’s, as required by the chiral superspace integral at point 0 in the reduction formula. We shall prove the following result: A sufficient condition for the vanishing of $`<0p_1\mathrm{}p_n>`$ is $$p_1>p_2+\mathrm{}+p_n4.$$ (3) This means that the extremal and next-to-extremal correlators are ruled out. Note that we shall not attempt to find out the most general conditions for the vanishing of $`<0p_1\mathrm{}p_n>`$. This would require a detailed study of the structure of the nilpotent covariants which is difficult to carry out; it is not clear that this would lead to any new results of a similar kind. To prove this, consider first the leading term in $`<0p_1\mathrm{}p_n>`$. To construct it we need a set of odd variables which are invariant under $`Q`$ supersymmetry and under the shift-like part of $`S`$ supersymmetry . $`Q`$ supersymmetry obviously suggests to use the combinations $$\theta _{0r}^\alpha =\theta _0^{i\alpha }u_{ri}^+\theta _r^{+\alpha },\delta _Q\theta _{0r}^\alpha =0,r=1,\mathrm{},n.$$ (4) Then we can form the following two cyclic combinations of three $`\theta _{0r}^\alpha `$: $$(\xi _{12r})_{\dot{\alpha }}=(12)\rho _{r\dot{\alpha }}+(2r)\rho _{1\dot{\alpha }}+(r1)\rho _{2\dot{\alpha }},r=3,4,\mathrm{},n$$ (5) where $$\rho _{r\dot{\alpha }}=x_{0r}^2(x_{0r}\theta _{0r})_{\dot{\alpha }}$$ (6) and $`x_{0r}x_{L0}x_{Ar}`$ are translation-invariant and $`(rs)u_r^{+i}u_{si}^+`$ are $`SU(2)`$-invariant combinations of the space-time and harmonic coordinates, correspondingly. It is now easy to check that $`\xi _{12r}`$ are completely shift-invariant, i.e., $$\delta _{Q+S}\xi _{12r}=O(\theta ^2).$$ (7) Here one makes use of the harmonic cyclic identity $$(rs)t_i+(st)r_i+(tr)s_i=0.$$ (8) Note the choice we have made: point 1 which carries the highest charge according to (3) is one of the two common points in the set of independent variables (5). Now, taking into account the required $`R`$-weight of the correlator $`<0p_1\mathrm{}p_n>`$, we can write down its leading term in the following form: $`0p_1\mathrm{}p_n=`$ $`{\displaystyle \underset{a,b,c,d=3}{\overset{n}{}}}\xi _{12a}\xi _{12b}\xi _{12c}\xi _{12d}F_{abcd}^{p_14\|p_24\|\mathrm{}p_a1\mathrm{}p_b1\mathrm{}p_c1\mathrm{}p_d1\mathrm{}\|p_n}(x,u)+O(\theta ^5\overline{\theta })`$ (9) where the Lorentz indices have been suppressed. The coefficient function $`F`$ depends on the space-time and harmonic variables and carries $`U(1)`$ charges to match those of the correlator and of the nilpotent prefactor. We want to study the consequences of the H-analyticity condition (2). It turns out that in order to prove (3) it will be sufficient to look at the terms not containing $`\theta _0`$, $`\theta _1^+`$ and $`\theta _2^+`$. In this case the variables (5) become very simple: $$\xi _{12a}(12)\theta _a^+$$ (10) (the space-time factor is of no importance). Consequently, eq. (9) is reduced to $$0p_1\mathrm{}p_n\underset{a,b,c,d=3}{\overset{n}{}}\theta _a^+\theta _b^+\theta _c^+\theta _d^+f_{abcd}^{p_1\|p_2\|\mathrm{}p_a1\mathrm{}p_b1\mathrm{}p_c1\mathrm{}p_d1\mathrm{}\|p_n}$$ (11) where $$f_{abcd}(12)^4F_{abcd}.$$ (12) Note that each of the coefficients $`f_{abcd}`$ is associated to a single and unique nilpotent structure $`\theta _a^+\theta _b^+\theta _c^+\theta _d^+`$. This means that we have to impose the H-analyticity condition (2) on each of the $`f_{abcd}`$’s independently. We recall that H-analyticity for a harmonic function of positive charge simply means that it is a polynomial in the harmonics of degree equal to the charge. Our functions (12) have to be $`SU(2)`$ invariant polynomials in the $`n`$ sets of harmonics. Then it becomes clear that, taking into account the restriction (3) on the charges, it is not possible to match $`p_1`$ copies of $`u_{1i}^+`$ with the remaining harmonics $`u^+`$ at points $`2,3,\mathrm{},n`$, so all the coefficients must vanish. We now turn to nilpotent invariants of subleading order, i.e. those involving $`\overline{\theta }^+`$’s in their expansion. The simplest one of them has 5 $`\theta `$’s and one $`\overline{\theta }`$. In the left-handed sector we can still use the $`Q`$ and $`S`$supersymmetry shift-invariant variables $`\xi _{12a}`$. In principle, in the right-handed sector we should employ the analogous shift-invariant variables made out of four $`\overline{\theta }^+`$’s (see ). However, they are very complicated and have a rather non-trivial harmonic dependence. Fortunately, we do not need them in the present context. It is sufficient to use variables which are only $`Q`$ supersymmetric: $$\overline{\theta }_{12a}=(12)\overline{\theta }_a+(2a)\overline{\theta }_1+(a1)\overline{\theta }_2.$$ (13) This means that the term of order $`5+1`$ will contain more coefficients compared to the true $`Q`$ and $`S`$covariant one, but our condition (3) turns out sufficient to eliminate all of them. Indeed, the general form of such a term is $$\underset{a,\mathrm{},f=3}{\overset{n}{}}\xi _{12a}\xi _{12b}\xi _{12c}\xi _{12d}\xi _{12e}\overline{\theta }_{12f}F_{abcdef}^{p_16\|p_26\|\mathrm{}p_a1\mathrm{}p_f1\mathrm{}\|p_n}.$$ (14) Once more, we are only interested in terms not containing $`\theta _0`$, $`\theta _1^+`$ and $`\theta _2^+`$, so (14) can be reduced to $$\underset{a,\mathrm{},f=3}{\overset{n}{}}\theta _a^+\mathrm{}\theta _e^+\overline{\theta }_f^+f_{abcdef}^{p_1\|p_2\|\mathrm{}p_a1\mathrm{}p_f1\mathrm{}\|p_n}$$ (15) where $$f_{abcdef}(12)^6F_{abcdef}.$$ (16) It then becomes clear that the following restriction on the charges: $$p_1>p_2+\mathrm{}+p_n6$$ (17) will be sufficient to kill all such coefficients. In fact, this condition clearly follows from the extremal or next-to-extremal constraints on the charges, so that the result holds to this order. The generalisation to higher-order nilpotents is obvious and so the result is established. ## 3 The coordinate approach The above result can also be derived in the coordinate formalism, as we shall now sketch. We are interested in the correlator $$G=<0p_1\mathrm{}p_n>$$ (18) where we have one insertion of the chiral operator $`\mathrm{tr}(W^2)`$ at point $`0`$ and $`n`$ hypermultiplets of with charges $`p_r`$ at the other $`n`$ points, $`1,\mathrm{}n`$. The Ward identity reads $$\left((V_0+2\mathrm{\Delta }_0)+\underset{r}{}(V_r+p_r\mathrm{\Delta }_r)\right)G=0$$ (19) where $`V_0`$ and $`V_r`$ are the superconformal Killing vectors in chiral superspace and analytic superspace respectively and $`\mathrm{\Delta }_0`$ and $`\mathrm{\Delta }_r`$ are the corresponding weight functions. These Ward Identities have been written out in detail elsewhere and we refer the reader to the literature for the details . In the present context we shall only need to use supersymmetries, dilations, R-symmetry and internal ($`SL(2)`$) transformations (for simplicity, we work in complex spacetime). The main difference with the harmonic formalism is that the internal transformations now appear as (holomorphic) conformal transformations of $`\text{}P^1`$. Thus we have internal translations, dilations and “conformal boosts”. The internal dilations correspond to the $`U(1)`$ transformations in the harmonic formalism. Initially, the space on which we are working has coordinates $`(x_0^{\alpha \dot{\alpha }},\theta _0^{\alpha 1}\theta ^\alpha ,\theta _0^{\alpha 2}\phi ^\alpha )`$ and $`(x_r^{\alpha \dot{\alpha }},\lambda _r^\alpha ,\pi _r^{\dot{\alpha }},y_r)`$ where $`x_0`$ and $`x_r`$ are chiral and analytic $`x`$’s respectively and the $`y_r`$’s are local coordinates on the $`n`$ copies of $`\text{}P^1`$. For translations, internal translations and Q-supersymmetries the weight functions vanish. These transformations are $`\delta x_0^{\alpha \dot{\alpha }}`$ $`=`$ $`B^{\alpha \dot{\alpha }}\overline{ϵ}_1^{\dot{\alpha }}\theta ^\alpha \overline{ϵ}_2^{\dot{\alpha }}\phi ^\alpha `$ $`\delta x_r^{\alpha \dot{\alpha }}`$ $`=`$ $`B^{\alpha \dot{\alpha }}ϵ^{\alpha 1}\pi _r^{\dot{\alpha }}\overline{ϵ}_2^{\dot{\alpha }}\lambda _r^\alpha `$ $`\delta y_r`$ $`=`$ $`B`$ $`\delta \theta ^\alpha `$ $`=`$ $`ϵ^{\alpha 1}`$ (20) $`\delta \phi ^\alpha `$ $`=`$ $`ϵ^{\alpha 2}+B\theta ^\alpha `$ $`\delta \lambda _r^\alpha `$ $`=`$ $`ϵ^{\alpha 2}ϵ^{\alpha 1}y_r`$ $`\delta \pi _r^{\dot{\alpha }}`$ $`=`$ $`\overline{ϵ}_1^{\dot{\alpha }}+\overline{ϵ}_2^{\dot{\alpha }}y_r`$ where the $`ϵ`$ parameters are the supersymmetry parameters and the $`B`$’s are the translational parameters. It is easy to solve the Ward identities for these transformations to eliminate four spinorial coordinates, one $`x`$ and one $`y`$. Doing this, one finds that $`G`$ can be taken to be a function of the following variables: $`(x_{0r}^{\alpha \dot{\alpha }},\zeta _r^\alpha ,\pi _{12r}^{\dot{\alpha }},y_{1r})`$. Here, we use $`y_{rs}=y_ry_s`$ and similarly for other coordinate differences, although the $`x`$ difference variables require a nilpotent correction. Thus we have $`x_{0r}^{\alpha \dot{\alpha }}`$ $`=`$ $`x_0^{\alpha \dot{\alpha }}x_r^{\alpha \dot{\alpha }}\theta ^\alpha \pi _r^{\dot{\alpha }}+{\displaystyle \frac{1}{(n1)}}{\displaystyle \underset{s,sr}{}}{\displaystyle \frac{\zeta _r^\alpha \pi _{rs}^{\dot{\alpha }}}{y_{rs}}}`$ $`\zeta _r^\alpha `$ $`=`$ $`y_r\theta ^\alpha +\lambda _r^\alpha \phi ^\alpha `$ (21) $`\pi _{12r}`$ $`=`$ $`y_{12}\pi _r^{\dot{\alpha }}+y_{r1}\pi _2^{\dot{\alpha }}+y_{2r}\pi _1^{\dot{\alpha }}`$ $`y_{1r}`$ $`=`$ $`y_1y_r`$ Note that we have $`n`$ $`x`$’s and undotted spinor coordinates, $`(n2)`$ dotted spinor coordinates and $`(n1)`$ $`y`$ coordinates left. We should remark that the choice of correction terms for the $`x`$ difference variables is not unique but this is not an issue which will be relevant for the ensuing discussion. The weights of these coordinates under (dilations, internal dilations, R) are as follows $`x`$ $`:`$ $`(1,0,0)`$ $`y`$ $`:`$ $`(0,1,0)`$ (22) $`\zeta `$ $`:`$ $`({\displaystyle \frac{1}{2}},{\displaystyle \frac{3}{2}},1)`$ $`\pi `$ $`:`$ $`({\displaystyle \frac{1}{2}},{\displaystyle \frac{3}{2}},1)`$ For R symmetry $`\mathrm{\Delta }_0=2R`$, so that G has R-weight four. Schematically it must depend on the odd variables in the following way $$G\zeta ^4(1+\mathrm{power}\mathrm{series}\mathrm{in}(\zeta \pi ))$$ (23) The idea is then to carry out the analysis order by order in odd variables. Since we shall only proceed to the second order when the first order vanishes this means that we can ignore the correction terms in $`x_{0r}`$ and that we can simplify the remaining S-supersymmetry transformations considerably. The only non-trivial simplified S-supersymmetry transformation with dotted parameters is $$\delta \zeta _r^\alpha =(\overline{\eta }_{\dot{\beta }}^2+y_r\overline{\eta }_{\dot{\beta }}^1)x_{or}^{\alpha \dot{\beta }}$$ (24) It is easy to find variables invariant under these transformations since they are essentially translations. These invariant variables are $$\xi _{12r}^{\dot{\alpha }}=y_{12}(x_{0r}^1)_{\dot{\alpha }\alpha }\zeta _r^\alpha +y_{r1}(x_{02}^1)_{\dot{\alpha }\alpha }\zeta _2^\alpha +y_{2r}(x_{01}^1)_{\dot{\alpha }\alpha }\zeta _1^\alpha $$ (25) The $`\xi `$’s have weights $`(1/2,3/2,1)`$. Note that, in this approximation, all the other coordinates are invariant, so we may now work with the set $`(x_{0r},\xi _{12r},\pi _{12r},y_{1r})`$. The leading term in $`G`$, $`G_0`$, say, can be written as $$G_0=\underset{abcd}{}\xi _{12a}\xi _{12b}\xi _{12c}\xi _{12d}F^{abcd}(x,y)$$ (26) where the Lorentz indices have been suppressed. To complete the analysis we need only consider the remaining internal symmetry transformations; we can forget about $`x`$ altogether. Internal dilational symmetry implies that $`F`$ has weight $`(_{r=1}^np_r)12`$. (The chiral function $`\mathrm{\Delta }_0=0`$ for this while $`\mathrm{\Delta }_r=1/2`$). The internal “conformal boosts” act in the following way $`\delta y_{rs}`$ $`=`$ $`C(y_r+y_s)y_{rs}`$ $`\delta \xi _{12r}`$ $`=`$ $`C(y_1+y_2+y_r)\xi _{12r}`$ (27) $`\delta \pi _{12r}`$ $`=`$ $`C(y_1+y_2+y_r)\pi _{12r}`$ where $`C`$ is the parameter. If we define an operator $`D`$ by $$D=\left(\frac{\delta y}{C}\frac{}{y}+\frac{\delta \xi }{C}\frac{}{\xi }+\frac{\delta \pi }{C}\frac{}{\pi }\right)$$ (28) where the sum is over these coordinates with the variations given in the previous equation, then the corresponding Ward identity implies that $$DG=(\underset{r=1}{\overset{n}{}}p_ry_r)G$$ (29) since $`\mathrm{\Delta }_0=0`$ and $`\mathrm{\Delta }_r=Cy_r`$, up to nilpotent terms arising from the variation of the $`x`$’s. The same equation holds (exactly) for $`G_0`$. It makes it clear that we can identify the internal dilation charges in the coordinate approach with the $`U(1)`$ charges in the harmonic formalism. From (29) we find $$\underset{abcd}{}\xi _{abcd}\left(DF^{abcd}+(4y_1+4y_2+y_a+\mathrm{}y_d)F^{abcd}\right)=(\underset{r=1}{\overset{n}{}}p_ry_r)\underset{abcd}{}\xi _{abcd}F^{abcd}$$ (30) where we have used the abbreviation $`\xi _{abcd}\xi _{12a}\xi _{12b}\xi _{12c}\xi _{12d}`$. Now $`\xi _{abcd}`$ includes the term $`(y_{12})^4\lambda _a\lambda _b\lambda _c\lambda _d`$, and the only place such a term appears in the sum for given values of $`a,b,c,d`$ is as a term in the corresponding $`\xi _{abcd}`$. From this we conclude that, for all choices of $`a,b,c,d`$, $$DF^{abcd}=(\underset{r=1}{\overset{n}{}}p_r^{}y_r)F^{abcd}$$ (31) where $$p_r^{}=\{\begin{array}{cc}p_r4,\hfill & \text{ }r=1,2\hfill \\ p_r,\hfill & \text{ }r1,2,a,b,c,d\hfill \\ p_r1,\hfill & \text{ }r=a,b,c,d\hfill \end{array}$$ (32) In other words, any such $`F^{abcd}`$ is a conformal function with the above weights. Now $`G`$ is analytic in the internal coordinates. Analyticity of the term in $`G_0`$ in which $`F^{abcd}`$ is multiplied by $`(y_{12})^4\lambda _a\lambda _b\lambda _c\lambda _d`$ implies that $`F^{abcd}`$ cannot have any singularities in $`y_{rs}`$ for $`r<s,s3`$. It can therefore be written as a sum of terms of the form $$F^{abcd}=(y_{12})^{q_2}\mathrm{}(y_{1n})^{q_n}\stackrel{~}{F}^{abcd}$$ (33) where $`_{r=2}^nq_r=p_14`$ and where each $`\stackrel{~}{F}^{abcd}`$ depends only on the coordinates $`y_{rs}`$ with $`2r<sn`$. $`\stackrel{~}{F}^{abcd}`$ has charges $`\stackrel{~}{p}_r,`$ where $$\stackrel{~}{p}_r=\{\begin{array}{cc}0,\hfill & \text{ }r=1\hfill \\ p_r4q_r,\hfill & \text{ }r=2\hfill \\ p_rq_r,\hfill & \text{ }r1,2,a,b,c,d\hfill \\ p_rq_r1,\hfill & \text{ }r=a,b,c,d\hfill \end{array}$$ (34) By analyticity, $`\stackrel{~}{F}^{abcd}`$ must be a homogeneous polynomial in $`y_{rs},2r<sn`$, and so $`\stackrel{~}{p}_r0,r1`$. On summing these constraints we find $$p_14=\underset{r=2}{\overset{n}{}}q_r\underset{r=2}{\overset{n}{}}p_r8$$ (35) Hence there is no solution if $`p_1>_{r=2}^np_r4`$, which covers the extremal and next-to-extremal cases. Given that the lowest order term vanishes, at the next order we can write $$G_1=\underset{abcdef}{}\xi _{12a}\xi _{12b}\xi _{12c}\xi _{12d}\xi _{12e}\pi _{12f}F^{abcdef}$$ (36) Now since $`\pi _{12r}`$ transforms in the same way as $`\xi _{12r}`$ under conformal boosts it follows that we can repeat the above argument straightforwardly. The only difference is that the bound is stronger; there is no solution if $`p_1>_{r=2}^np_r6`$. Clearly the argument extends to all orders in this manner. Finally we remark that although the above results for $`N=2`$ cover a wider class of theories than $`N=4`$ they are in another sense more restricted. This is because the $`N=2`$ hypermultiplet operators have maximum spin $`1`$ while the single-trace analytic oper ators in $`N=4`$ have maximum spin $`2`$. It should in principle be possible to carry out a similar sort of analysis to the one given here for $`N=2`$ directly in $`N=4`$ harmonic superspace, but it would be considerably more complicated due to the fact that the internal symmetry group is significantly larger. On the other hand it is not difficult to show that the leading terms of correlators of such operators are completely determined, by group theory, from $`N=2`$ hypermultiplet “component” correlators in both the extremal and next-to-extremal cases. Thus the leading terms of these correlators are trivial and it seems very likely that this result will extend to all orders in an expansion with respect to the third and fourth spinorial coordinates. Acknowledgements: This work was supported in part by the British-French scientific programme Alliance (project 98074), by the EU network on Integrability, non-perturbative effects and symmetry in quantum field theory (FMRX-CT96-0012), by the grant INTAS-96-0308 and by PPARC through SPG 613.
warning/0004/gr-qc0004052.html	ar5iv	text	# 1 Introduction ## 1 Introduction There has been much interest in black hole solutions in $`p`$-brane theory - because of the possible resolution of various puzzles associated with quantum gravity. A growth of interest in classical $`p`$-brane solutions of supergravities of various dimensions is inspired by a conjecture that $`D=11`$ supergravity is a low-energy effective field theory of eleven-dimensional fundamental $`M`$-theory, which (together with so called $`F`$-theory) is a candidate for unification of five known ten-dimensional superstring models. Classical $`p`$-brane solutions may be considered as an instrument for investigation of interlinks between superstrings and $`M`$-theory. In this paper we consider generalized bosonic sector (without Chern-Simons terms) of supergravity theories in the form of a multidimensional gravitational model with several dilatonic scalar fields and differential forms of various ranks admitting an interpretation in terms of intersecting $`p`$-branes. As was shown in , for cosmological and static spherically symmetric space-times the equations of motion of such a model are reduced to the Euler-Lagrange equations for the so-called pseudo-Euclidean Toda-like Lagrange system. We reproduce this result in Sec. 2. Methods for integrating of pseudo-Euclidean Toda-like systems (see, and references therein) are based on a Minkowski-like geometry for the characteristic vectors determining the potential of the pseudo-Euclidean Toda-like system. If the characteristic vectors form an orthogonal set, then the pseudo-Euclidean Toda-like system is integrable. The corresponding $`p`$-brane models have been studied in ,,,. Here we apply these methods for integrating the $`p`$-brane model reducible to the algebraic generalizations of an open Toda chain. The characteristic vectors of such models may be interpreted as root vectors of the semi-simple Lie algebra. We consider the Lie algebras of the Cartan types $`A_r,B_r,C_r`$. Using the technique suggested by Anderson for solving the Toda chain’s equations of motion, in Sec. 3,4 we integrate the $`p`$-brane models. In the last section of the paper we examine the metric obtained for some particular exact solution appearing for a quite wide class of $`p`$-brane models. This solution describes the nonextremal black hole under some condition. The corresponding ADM-mass and the Hawking temperature of a black hole are calculated. ## 2 The general model Following the papers ,, we consider here a classical model of gravity theory with several dilatonic scalar fields $`\phi ^\alpha `$ and differential $`n_a`$-forms $`F_{M_1\mathrm{}M_{n_a}}^a`$ in (pseudo)-Riemannian space-time manifold $`\mathrm{M}`$ of dimension $`D`$. The action of the model reads $$S=\underset{\mathrm{M}}{}d^Dz\sqrt{\|g\|}\left(R[g]\underset{\alpha ,\beta =1}{\overset{\omega }{}}h_{\alpha \beta }g^{MN}_M\phi ^\alpha _N\phi ^\beta \underset{a\mathrm{\Delta }}{}\frac{\text{e}^{2\lambda _a(\phi )}}{n_a!}(F^a)^2\right),$$ (2.1) where $`ds^2=g_{MN}dz^Mdz^N`$ is the metric with Lorentzian signature on the manifold $`\mathrm{M}`$ ($`M,N=0,1\mathrm{},D1`$), $`\|g\|=\|det(g_{MN})\|`$. $`(h_{\alpha \beta })`$ is a symmetrical positively definite $`\omega \times \omega `$ matrix, $`\lambda _a(\phi )`$ is a linear combination of the scalar fields, i.e. $$\lambda _a(\phi )=\underset{\alpha =1}{\overset{\omega }{}}\lambda _{a,\alpha }\phi ^\alpha ,$$ (2.2) where $`\lambda _{a,\alpha }`$ are the coupling constants. Furthermore $`F^a={\displaystyle \frac{1}{n_a!}}F_{M_1\mathrm{}M_{n_a}}^adz^{M_1}\mathrm{}dz^{M_{n_a}}=dA^a,`$ (2.3) $`A^a={\displaystyle \frac{1}{(n_a1)!}}A_{M_1\mathrm{}M_{n_a1}}^adz^{M_1}\mathrm{}dz^{M_{n_a1}},`$ (2.4) $`(F^a)^2=F_{M_1\mathrm{}M_{n_a}}^aF_{N_1\mathrm{}N_{n_a}}^ag^{M_1N_1}\mathrm{}g^{M_{n_a}N_{n_a}}.`$ (2.5) The field $`A_{M_1,\mathrm{},M_{n_a1}}^a`$ may be called a gauge potential corresponding to the field strength $`F_{M_1,\mathrm{},M_{n_a}}^a`$. By $`\mathrm{\Delta }`$ we denote some finite set. The action (2.1) leads to the following equations of motion $`R_{MN}{\displaystyle \frac{1}{2}}g_{MN}R=T_{MN},`$ (2.6) $`[g]\phi ^\alpha ={\displaystyle \underset{a\mathrm{\Delta }}{}}{\displaystyle \frac{1}{n_a!}}\lambda _a^\alpha \text{e}^{2\lambda _a(\phi )}(F^a)^2,\alpha =1,\mathrm{},\omega ,`$ (2.7) $`_{M_1}[g](\text{e}^{2\lambda _a(\phi )}F^{a,M_1\mathrm{}M_{n_a}})=0,a\mathrm{\Delta }.`$ (2.8) The right side of the Einstein equations (2.6) looks as follows $`T_{MN}=T_{MN}[\phi ]+T_{MN}[F],`$ (2.9) where we denoted $`T_{MN}[\phi ]={\displaystyle \underset{\alpha ,\beta =1}{\overset{\omega }{}}}h_{\alpha \beta }\left(_M\phi ^\alpha _N\phi ^\beta {\displaystyle \frac{1}{2}}g_{MN}_P\phi ^\alpha ^P\phi ^\beta \right),`$ (2.10) $`T_{MN}[F]={\displaystyle \underset{a\mathrm{\Delta }}{}}{\displaystyle \frac{\text{e}^{2\lambda _a(\phi )}}{n_a!}}\left({\displaystyle \frac{1}{2}}g_{MN}(F^a)^2+n_aF_{MM_2\mathrm{}M_{n_a}}^aF_N^{a,M_2\mathrm{}M_{n_a}}\right).`$ (2.11) In (2.7),(2.8) we denoted the Laplace-Beltrami operator and covariant derivative with respect to the metric $`g_{MN}`$ by $`[g]`$ and $`_M[g]`$, respectively. The constants $`\lambda _a^\alpha `$ in (2.7) are introduced by $$\lambda _a^\alpha =\underset{\alpha ,\beta =1}{\overset{\omega }{}}h^{\alpha \beta }\lambda _{a,\beta },$$ (2.12) where $`(h^{\alpha \beta })`$ is the inverse matrix to $`(h_{\alpha \beta })`$. Consider the model introduced under the following assumptions. Let the $`D`$-dimensional space-time M be decomposed into the direct product of $`\text{IR}_+`$ (corresponding to a radial coordinate $`u`$), $`d_0`$-dimensional sphere $`S^{d_0}`$ ($`d_02`$), time axis $`\text{IR}_t`$ and $`(n1)`$ factor spaces $`M_2,\mathrm{},M_n`$, i.e. $$\mathrm{M}=\text{IR}_+\times S^{d_0}\times \text{IR}_t\times M_2^{d_2}\mathrm{}\times M_n^{d_n},n2.$$ (2.13) The metric on M is assumed correspondingly to be $$ds^2=\text{e}^{2\gamma (u)}du^2+\text{e}^{2x^0(u)}d\mathrm{\Omega }_{d_0}^2\text{e}^{2x^1(u)}dt^2+\underset{i=2}{\overset{n}{}}\text{e}^{2x^i(u)}ds_i^2,$$ (2.14) where $`u`$ is the radial coordinate, $`d\mathrm{\Omega }_{d_0}^2=g_{m_0n_0}^0(y_0)dy_0^{m_0}dy_0^{n_0}`$ is the line element on $`d_0`$-dimensional unit sphere, $`t`$ is the time coordinate, $`ds_i^2=g_{m_in_i}^i(y_i)dy_i^{m_i}dy_i^{n_i}`$ is the positively definite metric on the $`d_i`$-dimensional factor space $`M_i`$, $`\gamma (u),x^2(u),\mathrm{},x^n(u)`$ are scalar functions of the radial coordinate $`u`$. Herein, for reasons of simplicity, only Ricci-flat spaces $`M_2,\mathrm{},M_n`$ are assumed (i.e. the components of the Ricci tensor for the metrics $`g_{m_in_i}^i`$ are zero). It is useful to consider $`S^{d_0}`$ and $`\text{IR}_t`$ as factor spaces $`M_0`$ and $`M_1`$, respectively. So we put $`M_0S^{d_0},`$ (2.15) $`M_1\text{IR}_t,d_1=1.`$ (2.16) We split the coordinates on M into the following ranges: $$(z^0,z^1,\mathrm{},z^{d_0},z^{d_0+1},\mathrm{},z^{Dd_n},\mathrm{},z^{D1})=(u,y_0^1,\mathrm{},y_0^{d_0},t,\mathrm{},y_n^1,\mathrm{},y_n^{d_n}).$$ (2.17) We introduce the following $`d_i`$-forms on M $$\tau _1=dt,\tau _i=\sqrt{det(g_{m_in_i}^i)}dy_i^1\mathrm{}dy_i^{d_i},i=0,2,\mathrm{},n.$$ (2.18) Clearly, the canonical projection $`\widehat{p}_i:\mathrm{M}M_i`$ of $`\tau _i`$ provides with the volume form of $`M_i`$. In order to construct the $`p`$-brane worldvolumes we introduce submanifolds of the following type $$M_I=M_{i_1}\times \mathrm{}\times M_{i_r},$$ (2.19) where $$I=\{i_1,\mathrm{},i_r\},i_1<\mathrm{}<i_r,$$ (2.20) is any ordered non-empty subset of natural numbers $`2,\mathrm{},n`$. Let $`\mathrm{\Omega }_0`$ be the set of all such elements including the empty set, i.e. $$\mathrm{\Omega }_0=\{\mathrm{},\{2\},\{3\},\mathrm{},\{n\},\{2,3\},\mathrm{},\{2,3,\mathrm{},n\}\}.$$ (2.21) By definition, put $$\overline{I}\{2,\mathrm{},n\}I.$$ (2.22) In this paper we consider electrically charged $`p`$-branes with the following worldvolumes $$M_I^{(e)}=\text{IR}_t\times M_I,I=\{i_1,\mathrm{},i_r\}\mathrm{\Omega }_0.$$ (2.23) For empty $`I=\mathrm{}`$ we put $`M_I^{(e)}=\text{IR}_t`$. The dimension of $`M_I^{(e)}`$ is given by $$d(I)dimM_I^{(e)}=1+d_{i_1}+\mathrm{}+d_{i_r}.$$ (2.24) ($`d(I)=1`$ for $`I=\mathrm{}`$). The canonical projection $`\widehat{p}_I:\mathrm{M}M_I^{(e)}`$ of the following $`d(I)`$-form $$\tau (I)=dt\tau _{i_1}\mathrm{}\tau _{i_r}$$ (2.25) is the volume form of $`M_I^{(e)}`$. We put $`\tau (I)=dt`$ for $`I=\mathrm{}`$. In accordance with the terminology of $`p`$-brane theory an $`(n_a1)`$-form potential $$A^{(a,e,I)}=\mathrm{\Phi }^{(a,e,I)}(u)\tau (I),\text{rank}A^{(a,e,I)}n_a1=d(I),a\mathrm{\Delta },$$ (2.26) where $`\mathrm{\Phi }^{(a,e,I)}(u)`$ is a scalar function, describes an electrically charged $`p`$-brane ($`p=n_a2`$) with the worldvolume $`M_I^{(e)}`$. Moreover, the submanifold $`\text{IR}_+\times S^{d_0}\times M_{\overline{I}}`$ ($`\text{IR}_+\times S^{d_0}`$ for $`I=\{2,\mathrm{},n\}`$) is the so-called transverse space for this $`p`$-brane. The $`n_a`$-form field strength corresponding to $`A^{(a,e,I)}`$ was defined by (2.3) and may be written as $$F^{(a,e,I)}=d\mathrm{\Phi }^{(a,e,I)}(u)\tau (I)=\dot{\mathrm{\Phi }}^{(a,e,I)}(u)du\tau (I).$$ (2.27) The overdot means a derivative with respect to the radial coordinate $`u`$. An $`n_b`$-form field strength $$F^{(b,m,J)}=\text{e}^{2\lambda _b(\phi )}\left(d\mathrm{\Phi }^{(b,m,J)}(u)\tau (J)\right),J\mathrm{\Omega }_0,b\mathrm{\Delta },$$ (2.28) describes a $`p`$-brane ($`p=n_b1=Dd(J)2`$) with a magnetic-type charge. The submanifold $$M_J^{(m)}=S^{d_0}\times M_{\overline{J}}$$ (2.29) is a worldvolume of this $`p`$-brane. Clearly, $`M_J^{(m)}=S^{d_0}`$ for $`J=\{2,\mathrm{},n\}`$. By $``$ we denoted the Hodge operator on the manifold $`(\mathrm{M},g)`$, i.e. $$(F)_{M_1\mathrm{}M_{Dr}}=\frac{\sqrt{\|g\|}}{r!}\epsilon _{N_1\mathrm{}N_rM_1\mathrm{}M_{Dr}}F^{N_1\mathrm{}N_r}.$$ (2.30) In this paper we consider the so-called composite $`p`$-branes , i.e., by definition we put $$F^a=\underset{I\mathrm{\Omega }_{a,e}}{}F^{(a,e,I)}+\underset{J\mathrm{\Omega }_{a,m}}{}F^{(a,m,J)},$$ (2.31) where $`\mathrm{\Omega }_{a,e}\mathrm{\Omega }_0`$ is a subset (which may be empty) of all $`I\mathrm{\Omega }_0`$ such that $`d(I)+1=n_a\text{rank}F^{(a,e,I)}`$. Moreover, $`\mathrm{\Omega }_{a,m}\mathrm{\Omega }_0`$ is a subset (which may be empty) of all $`J\mathrm{\Omega }_0`$ such that $`Dd(J)1=dimM_J^{(m)}=n_a\text{rank}F^{(a,m,J)}`$. Evidently, $`\mathrm{\Omega }_{a,m}=\mathrm{}`$ for $`n_a=D1,D`$. We obtain the following non-zero components of the Ricci tensor for the metric (2.14) $`R_0^0=\text{e}^{2\gamma }\left({\displaystyle \underset{k=0}{\overset{n}{}}}d_k(\dot{x}^k)^2+\ddot{\gamma _0}\dot{\gamma }\dot{\gamma _0}\right),`$ (2.32) $`R_{n_k}^{m_k}=\left\{\delta _0^k(d_01)\text{e}^{2x^k}\left[\ddot{x}^k+\dot{x}^k(\dot{\gamma _0}\dot{\gamma })\right]\text{e}^{2\gamma }\right\}\delta _{n_k}^{m_k},`$ (2.33) where we denoted $$\gamma _0=\underset{k=0}{\overset{n}{}}d_kx^k.$$ (2.34) Indices $`m_k`$ and $`n_k`$ in (2.33) for $`k=0,\mathrm{},n`$ run over from ($`D_{l=k}^nd_l`$) to ($`D_{l=k}^nd_l+d_k1`$). We recall that $`D=1+_{k=0}^nd_k=dim\mathrm{M}`$. Under the above assumptions related to the $`F^a`$-fields and the metric (2.14) the Maxwell-like equations (2.8) and the Bianchi identities $`dF^a=0`$ have the following form, correspondingly $`{\displaystyle \frac{d}{du}}\left[\text{e}^{\gamma _0\gamma 2\sigma (I)+2\lambda _a(\varphi )}\dot{\mathrm{\Phi }}^{(a,e,I)}(u)\right]=0,I\mathrm{\Omega }_{a,e},a\mathrm{\Delta }`$ (2.35) $`{\displaystyle \frac{d}{du}}\left[\text{e}^{\gamma _0\gamma 2\sigma (J)2\lambda _a(\varphi )}\dot{\mathrm{\Phi }}^{(a,m,J)}(u)\right]=0,J\mathrm{\Omega }_{a,m},a\mathrm{\Delta },`$ (2.36) where $$\sigma (I)=d_1x^1+\underset{iI}{}d_ix^i.$$ (2.37) For empty $`I=\mathrm{}`$ we put $`\sigma (I)=d_1x^1`$. To denote $`F^a`$-fields and their potentials, it is useful the following collective index $$s=(a,v,I),I\mathrm{\Omega }_{a,v},v=e,m,a\mathrm{\Delta }.$$ (2.38) By $`S`$ we denote the set of all elements $`s`$, i.e. $$S=\underset{v=e,m}{}\left(\underset{a\mathrm{\Delta }}{}\{a\}\times \{v\}\times \mathrm{\Omega }_{a,v}\right).$$ (2.39) Integrating (2.35) and (2.36), we get $$\dot{\mathrm{\Phi }}^s(u)=Q_s\mathrm{exp}\left[\gamma \gamma _0+2\sigma (I_s)2\chi _s\lambda _{a_s}(\varphi )\right],s=(a_s,v_s,I_s)S,$$ (2.40) where $`\chi _s=+1,v_s=e,`$ (2.41) $`\chi _s=1,v_s=m.`$ (2.42) $`Q_s`$ are arbitrary constants. Let $`S_{}S`$ be a subset of all $`sS`$ such that $`Q_s0`$. To obtain the tensors $`T[F^a]_N^M`$ in a block-diagonal form, we put the following restriction: there are no elements $`(a,v,I),(a,v,J)S_{}`$ such that $$I=(IJ)\{i\},J=(IJ)\{j\},ij,d_i=d_j=1,$$ (2.43) where $`i,j=2,\mathrm{},n`$. Here the intersection $`IJ`$ may be empty. The total energy-momentum tensor of $`F^a`$-fields has a block-diagonal form whenever the restriction (2.43) is valid. Using (2.40), we present its non-zero components in the form $`T[F^a]_0^0={\displaystyle \frac{1}{2}}\text{e}^{2\gamma _0}{\displaystyle \underset{sS_{}}{}}Q_s^2\mathrm{exp}\left[2\sigma (I_s)2\chi _s\lambda _{a_s}(\varphi )\right],`$ (2.44) $`T[F^a]_{n_k}^{m_k}={\displaystyle \frac{1}{2}}\text{e}^{2\gamma _0}\left({\displaystyle \underset{sS_{}}{}}(2\delta _{kI_s}1)Q_s^2\mathrm{exp}\left[2\sigma (I_s)2\chi _s\lambda _{a_s}(\varphi )\right]\right)\delta _{n_k}^{m_k},`$ (2.45) where $$\delta _{kI}=\underset{iI}{}\delta _{ki}+\delta _{k1},I\mathrm{\Omega }_0,k=0,1,\mathrm{},n.$$ (2.46) We put $`\delta _{kI}=\delta _{k1}`$ for $`I=\mathrm{}`$. Evidently, $`\delta _{0I}=0`$ and $`\delta _{1I}=1`$ for any $`I\mathrm{\Omega }_0`$. We assume that the dilatonic scalar fields $`\phi ^\alpha `$ depend only on the radial coordinate $`u`$. Under this assumption the total energy-momentum tensor of the dilatonic scalar fields reads $$\left(T[\phi ]_N^M\right)=\frac{1}{2}\text{e}^{2\gamma }\left(\underset{\alpha ,\beta =1}{\overset{\omega }{}}h_{\alpha \beta }\dot{\phi }^\alpha \dot{\phi }^\beta \right)\text{diag}(1,\underset{d_0\mathrm{times}}{\underset{}{1,\mathrm{},1}},1,1,\mathrm{},1).$$ (2.47) The Einstein equations (2.6) can be written as $`R_N^M=T_N^MT\delta _N^M/(D2)`$. Further we employ the equations $`R_0^0R/2=T_0^0`$ and $`R_{n_k}^{m_k}=T_{n_k}^{m_k}T\delta _{n_k}^{m_k}/(D2)`$. Using (2.32),(2.33),(2.44)(2.45),(2.47), we obtain these equations in the form $$\frac{1}{2}\left(\underset{k,l=0}{\overset{n}{}}G_{kl}\dot{x}^k\dot{x}^l+\underset{\alpha ,\beta =1}{\overset{\omega }{}}h_{\alpha \beta }\dot{\phi }^\alpha \dot{\phi }^\beta \right)+V=0,$$ (2.48) $`\ddot{x}^k+(\dot{\gamma _0}\dot{\gamma })\dot{x}^k`$ $`=`$ $`\text{e}^{2\gamma }\{\delta _0^k(d_01)\text{e}^{2x^k}`$ (2.49) $`+`$ $`\text{e}^{2\gamma _0}{\displaystyle \underset{sS_{}}{}}\chi _s(\delta _{kI_s}{\displaystyle \frac{d(I_s)}{D2}})Q_s^2\mathrm{exp}[2\sigma (I_s)2\chi _s\lambda _{a_s}(\varphi )]\}`$ where we denoted $$G_{kl}=d_k\delta _{kl}d_kd_l,k,l=0,1,\mathrm{},n,$$ (2.50) $`V={\displaystyle \frac{1}{2}}\text{e}^{2\gamma }\left\{(d_01)d_0\text{e}^{2x^0}\text{e}^{2\gamma _0}{\displaystyle \underset{sS_{}}{}}Q_s^2\mathrm{exp}\left[2\sigma (I_s)2\chi _s\lambda _{a_s}(\varphi )\right]\right\}.`$ (2.51) Under the above assumptions equations (2.7) have the form $`\ddot{\phi }^\alpha +(\dot{\gamma _0}\dot{\gamma })\dot{\phi }^\alpha =\text{e}^{2\gamma 2\gamma _0}{\displaystyle \underset{sS_{}}{}}\chi _s\lambda _{a_s}^\alpha Q_s^2\mathrm{exp}\left[2\sigma (I_s)2\chi _s\lambda _{a_s}(\varphi )\right].`$ (2.52) It is not difficult to verify that equations (2.48),(2.49),(2.53) may be presented as the Euler-Lagrange equations obtained from the Lagrangian $$L=\text{e}^{\gamma _0\gamma }\left[\frac{1}{2}\left(\underset{k,l=0}{\overset{n}{}}G_{kl}\dot{x}^k\dot{x}^l+\underset{\alpha ,\beta =1}{\overset{\omega }{}}h_{\alpha \beta }\dot{\phi }^\alpha \dot{\phi }^\beta \right)V\right]$$ (2.53) viewed as a function of the generalized coordinates $`\gamma ,x^k,\phi ^\alpha `$. The equation $`L/\gamma =d(L/\dot{\gamma })/du`$ leads to the zero-energy constraint (2.48). After the gauge fixing: $`\gamma =F(x^k,\phi ^\alpha )`$ the equations (2.49),(2.52) may be considered as the Euler-Lagrange equations obtained from the Lagrangian (2.53) under the constraint (2.48). Further, we use the so-called harmonic gauge $$\gamma \gamma _0=\underset{k=0}{\overset{n}{}}d_kx^k.$$ (2.54) It is easy to check that the radial coordinate $`u`$ is a harmonic function in this gauge, i.e. $`[g]u=0`$. Let us introduce an $`(n+\omega +1)`$-dimensional real vector space $`\text{IR}^{n+\omega +1}`$. Denote by $`e_A,A=0,1,\mathrm{},n+\omega `$, the canonical basis in $`\text{IR}^{n+\omega +1}`$ ($`e_1=(1,0,\mathrm{},0`$) etc.). Define the following vectors: 1. The vector whose coordinates are to be found $$x(u)=\underset{A=0}{\overset{n+\omega }{}}x^A(u)e_A,\left(x^A(u)\right)=(x^0(u),\mathrm{},x^n(u),\phi ^1(u),\mathrm{},\phi ^\omega (u)).$$ (2.55) 2. The vector corresponding to the factor-space $`M_0S^{d_0}`$ with a non-zero Ricci tensor. Hereafter, we call it by the vector induced by the curvature of $`M_0`$ $$V_0=\underset{A=0}{\overset{n+\omega }{}}V_0^Ae_A=\frac{2}{d_0}e_0,\left(V_0^A\right)=2(\frac{1}{d_0},0,\mathrm{},0).$$ (2.56) 3. The vector induced by a $`p`$-brane $$U_s=\underset{A=0}{\overset{n+\omega }{}}U_s^Ae_A,\left(U_s^A\right)=2(\underset{n+1}{\underset{}{\delta _{kI}d(I_s)/(D2)}},\underset{\omega }{\underset{}{\chi _s\lambda _{a_s}^\alpha }}).$$ (2.57) A set of the vectors $`U_s`$ characterizes the space-time M, the configuration of $`p`$-branes and their couplings to dilatonic scalar fields. Further these vectors are called characteristic. Let $`<.,.>`$ be a symmetrical bilinear form on $`\text{IR}^{n+\omega +1}`$ such that $$<e_A,e_B>=\overline{G}_{AB},$$ (2.58) where we put by definition $`(\overline{G}_{AB})=\left(\begin{array}{cc}G_{kl}& 0\\ 0& h_{\alpha \beta }\end{array}\right).`$ (2.61) The form is nondegenerate, the inverse matrix to $`(G_{AB})`$ reads $`(\overline{G}^{AB})=\left(\begin{array}{cc}\frac{\delta ^{kl}}{d_k}+\frac{1}{2D}& 0\\ 0& h^{\alpha \beta }\end{array}\right).`$ (2.64) The form $`<.,.>`$ endows the space $`\text{IR}^{n+\omega +1}`$ with the metric, whose signature is $`(,+,\mathrm{},+)`$ . By the usual way we may introduce the covariant components of vectors. For the vectors $`V_0,U_s`$ the covariant components have the form $`V_{0,A}=2(\underset{n+1}{\underset{}{d_k\delta _{0k}}},\underset{\omega }{\underset{}{0,\mathrm{},0}}),`$ (2.65) $`U_{s,A}=2(\underset{n+1}{\underset{}{d_k\delta _{kI}}},\underset{\omega }{\underset{}{\chi _s\lambda _{a_s,\alpha }}}),`$ (2.66) The values of the bilinear form $`<.,.>`$ for $`V_0,U_s`$ look as follows $`<V_0,V_0>=4{\displaystyle \frac{d_01}{d_0}},`$ (2.67) $`<V_0,U_s>=0,sS,`$ (2.68) $`<U_s,U_s^{}>=4\left[d(I_sI_s^{}){\displaystyle \frac{d(I_s)d(I_s^{})}{D2}}+\chi _s\chi _s^{}{\displaystyle \underset{\alpha ,\beta =1}{\overset{\omega }{}}}h_{\alpha \beta }\lambda _{a_s}^\alpha \lambda _{a_s^{}}^\beta \right],`$ (2.69) where $`s=(a_s,v_s,I_s),s^{}=(a_s^{},v_s^{},I_s^{})S`$. A vector $`yR^{n+\omega +1}`$ is called time-like, space-like or isotropic, if $`<y,y>`$ has negative, positive or null values respectively. Vectors $`y`$ and $`z`$ are called orthogonal if $`<y,z>=0`$. It should be noted that the curvature induced vector $`V_0`$ is always time-like, while the $`p`$-brane induced vector $`U_s`$ admits any value of $`<U_s,U_s>`$. We mention that $`V_0`$ and $`U_s`$ are always orthogonal. Using the notation $`<.,.>`$ and the vectors (2.55)-(2.57), we represent the Lagrangian (2.53) and the zero-energy constraint (2.48) with respect to a harmonic time gauge in the form $`L={\displaystyle \frac{1}{2}}<\dot{x},\dot{x}>V,`$ (2.70) $`E={\displaystyle \frac{1}{2}}<\dot{x},\dot{x}>+V0,`$ (2.71) where the potential $`V`$ reads $$V=a^{(0)}\text{e}^{<V_0,x>}+\underset{sS_{}}{}a^{(s)}\text{e}^{<U_s,x>}.$$ (2.72) The following notation is used $`a^{(0)}={\displaystyle \frac{d_0(d_01)}{2}},a^{(s)}={\displaystyle \frac{1}{2}}Q_s^2,sS_{}.`$ (2.73) ¿From the mathematical viewpoint the obtaining of exact solutions in the $`p`$-brane model under consideration is reduced to integration of equations of motion for a system with $`(n+\omega +1)`$ degrees of freedom described by the Lagrangian of the form $$L=\frac{1}{2}<\dot{x},\dot{x}>\underset{\mu =1}{\overset{r}{}}a^{(\mu )}\text{e}^{<b_\mu ,x>},$$ (2.74) where $`x,b_\mu \text{IR}^{n+\omega +1}`$. It should be noted that the kinetic term $`<\dot{x},\dot{x}>`$ is not a positively definite quadratic form as there usually takes place in classical mechanics. Due to the pseudo-Euclidean signature $`(,+,\mathrm{},+)`$ of the form $`<.,.>`$ such systems may be called the pseudo-Euclidean Toda-like systems as the potential like that given in (2.74) defines the algebraic generalizations of the Toda chain . well-known in classical mechanics . ## 3 Integration of the $`p`$-brane model with linearly independent characteristic vectors We recall that $`S_{}S`$ is the subset of all $`sS`$ such that $`Q_s0`$. Define a bijection $`f:S_{}\{1,2,\mathrm{},r\}`$, where we denote by $`r`$ the cardinal number of $`S_{}`$, i.e. $$r=\|S_{}\|.$$ (3.1) Denote the natural number $`f(s)`$ corresponding to $`sS_{}`$ by the same letter $`s`$. The problem consists in integrating the equations of motion obtained from the Lagrangian (2.70) under the zero-energy constraint (2.71). Suppose the characteristic vectors $`U_s\text{IR}^{n+\omega +1}`$, induced by $`p`$-branes are linearly independent. Then $`rn+\omega `$. We introduce a basis $`\{f_A\}`$ in $`\text{IR}^{n+\omega +1}`$ in the following manner $`f_0={\displaystyle \frac{V_0}{<V_0,V_0>}},f_s={\displaystyle \frac{2U_s}{<U_s,U_s>}},s=1,\mathrm{},r,`$ (3.2) $`<f_A,f_p>=\delta _{Ap},A=0,\mathrm{},n+\omega ,p=r+1,\mathrm{},n+\omega .`$ (3.3) Notice that if $`r\|S_{}\|=n+\omega `$ then the basis $`\{f_A\}`$ does not contain the vectors $`f_p`$ with $`pr+1`$. We also mention that due to the relation (2.68) we get $`<f_0,f_s>=0`$ for $`s=1,\mathrm{},r`$. It is not difficult to prove that the vectors $`f_1,\mathrm{},f_r,f_{r+1},\mathrm{},f_{n+\omega }`$ must be space-like. Using the decomposition $$x(u)=q^0(u)f_0+\underset{s=1}{\overset{r}{}}[q^s(u)\mathrm{ln}C^s]f_s+\underset{p=r+1}{\overset{n+\omega }{}}q^p(u)f_p,$$ (3.4) we present the Lagrangian (2.70) and the constraint (2.71) in the form $`L=L_0+L_T+L_P,`$ (3.5) $`E=E_0+E_T+E_P0,`$ (3.6) where $`L_0={\displaystyle \frac{(\dot{q}^0)^2}{2<V_0,V_0>}}{\displaystyle \frac{d_0(d_01)}{2}}\text{e}^{q^0},`$ (3.7) $`E_0={\displaystyle \frac{(\dot{q}^0)^2}{2<V_0,V_0>}}+{\displaystyle \frac{d_0(d_01)}{2}}\text{e}^{q^0},`$ (3.8) $`L_T={\displaystyle \underset{s,s^{}=1}{\overset{r}{}}}{\displaystyle \frac{C_{ss^{}}}{<U_s,U_s>}}\dot{q}^s\dot{q}^s^{}+{\displaystyle \underset{s=1}{\overset{r}{}}}{\displaystyle \frac{2}{<U_s,U_s>}}\mathrm{exp}\left[{\displaystyle \underset{s^{}=1}{\overset{r}{}}}C_{ss^{}}q^s^{}\right],`$ (3.9) $`E_T={\displaystyle \underset{s,s^{}=1}{\overset{r}{}}}{\displaystyle \frac{C_{ss^{}}}{<U_s,U_s>}}\dot{q}^s\dot{q}^s^{}{\displaystyle \underset{s=1}{\overset{r}{}}}{\displaystyle \frac{2}{<U_s,U_s>}}\mathrm{exp}\left[{\displaystyle \underset{s^{}=1}{\overset{r}{}}}C_{ss^{}}q^s^{}\right],`$ (3.10) $`L_P=E_P={\displaystyle \frac{1}{2}}{\displaystyle \underset{p=r+1}{\overset{n+\omega }{}}}(\dot{q}^p)^2.`$ (3.11) We introduced, in (3.9),(3.10), the nondegenerate Cartan-type matrix $`(C_{ss^{}})`$ by the following manner $$C_{ss^{}}=2\frac{<U_s,U_s^{}>}{<U_s^{},U_s^{}>},s,s^{}=1,\mathrm{},r.$$ (3.12) The constants $`C^s`$ in the decomposition (3.4) are defined by $$C^s=\underset{s^{}=1}{\overset{r}{}}\left[\frac{<U_s^{},U_s^{}>}{4}Q_s^{}^2\right]^{C^{ss^{}}},s=1,\mathrm{},r,$$ (3.13) where $`(C^{ss^{}})`$ is the inverse matrix to $`(C_{ss^{}})`$. The Euler-Lagrange equations for $`q^{r+1}(u),\mathrm{},q^{n+\omega }(u)`$ read $`\ddot{q}^{r+1}(u)=\mathrm{}=\ddot{q}^{n+\omega }(u)=0`$. Integrating them, we get $`q^p(u)=a^pu+b^p,p=r+1,\mathrm{},n+\omega ,`$ (3.14) $`E_P={\displaystyle \frac{1}{2}}{\displaystyle \underset{p=r+1}{\overset{n+\omega }{}}}(a^p)^20,`$ (3.15) where the constants $`a^p,b^p`$ are arbitrary. For $`q^0(u)`$ we get the Liouville equation. The result of its integration reads $$\text{e}^{q^0(u)/2}=F_0(uu_0),$$ (3.16) where $`u_0`$ is an arbitrary constant. The function $`F_0`$ is defined by $$F_0(u)=\frac{\mathrm{sin}\left[\sqrt{\frac{2E_0}{d_0(d_01)}}(d_01)u\right]}{\sqrt{\frac{2E_0}{d_0(d_01)}}}.$$ (3.17) This representation implies $`F_0(u)`$ $`=(d_01)u,E_0=0,`$ (3.20) $`={\displaystyle \frac{\mathrm{sin}\left[\sqrt{\frac{2E_0}{d_0(d_01)}}(d_01)u\right]}{\sqrt{\frac{2E_0}{d_0(d_01)}}}},E_0>0`$ $`={\displaystyle \frac{\mathrm{sinh}\left[\sqrt{\frac{2\|E_0\|}{d_0(d_01)}}(d_01)u\right]}{\sqrt{\frac{2\|E_0\|}{d_0(d_01)}}}},E_0<0.`$ The equations of motion for $`q^1(u),\mathrm{},q^r(u)`$ look as follows $$\ddot{q^s}=\mathrm{exp}\left[\underset{s^{}=1}{\overset{r}{}}C_{ss^{}}q^s^{}\right],s,=1,\mathrm{},r.$$ (3.21) Using the transformation $$F_s(u)=\text{e}^{q^s(u)},$$ (3.22) we present the set of equations (3.21) in the form $$\dot{F}_s^2F_s\ddot{F}_s=F_s^2\underset{s^{}=1}{\overset{r}{}}(F_s^{})^{C_{ss^{}}}.$$ (3.23) The set of equations (3.21) proved to be completely integrable if $`(C_{ss^{}})`$ is the Cartan matrix of a simple complex Lie algebra. The general solutions for some algebras as well as some particular solution of the set (3.21) for quite a wide class of matrices $`(C_{ss^{}})`$ will be considered in the next sections. Here we suppose that the functions are known and the corresponding integral of motion (3.10) is calculated. Combining (3.2),(3.14),(3.16)(3.22), we present the decomposition (3.4) in the following form $$x(u)=\mathrm{ln}[F_0(uu_0)]\frac{2V_0}{<V_0,V_0>}\underset{s=1}{\overset{r}{}}\mathrm{ln}[C^sF_s(u)]\frac{2U_s}{<U_s,U_s>}+uQ+P,$$ (3.24) where vectors $`Q,P\text{IR}^{n+\omega +1}`$ are defined by $$Q=\underset{p=r+1}{\overset{n+\omega }{}}a^pf_p\underset{A=0}{\overset{n+\omega }{}}Q^Ae_A,P=\underset{p=r+1}{\overset{n+\omega }{}}b^pf_p\underset{A=0}{\overset{n+\omega }{}}P^Ae_A,$$ (3.25) Due to the assumptions (3.3) their coordinates $`Q^A,P^A`$ w.r.t. the canonical basis $`\{e_A\}`$satisfy the constraints $`<Q,V_0>=2{\displaystyle \underset{k=0}{\overset{n}{}}}Q^k(d_k\delta _{k0})=0,<P,V_0>=2{\displaystyle \underset{k=0}{\overset{n}{}}}P^k(d_k\delta _{k0})=0.`$ (3.26) $`<Q,U_s>={\displaystyle \underset{A=0}{\overset{n+\omega }{}}}Q^AU_{s,A}=0,<P,U_s>={\displaystyle \underset{A=0}{\overset{n+\omega }{}}}P^AU_{s,A}=0,s=1,\mathrm{},r.`$ (3.27) Finally, the exact solution can be summarized as follows. 1. The metric (2.14) in the harmonic time gauge (2.54) reads $`ds^2`$ $`={\displaystyle \underset{s=1}{\overset{r}{}}}\left[C^sF_s(u)\right]^{\frac{8d(I_s)}{(D2)<U_s,U_s>}}\{[F_0(uu_0)]^{\frac{2}{1d_0}}\text{e}^{2Q^0u+2P^0}[F_0^2(uu_0)du^2+d\mathrm{\Omega }_{d_0}^2]`$ (3.28) $``$ $`{\displaystyle \underset{s=1}{\overset{r}{}}}\left[C^sF_s(u)\right]^{\frac{8}{<U_s,U_s>}}\text{e}^{2Q^1u+2P^1}dt^2+{\displaystyle \underset{i=2}{\overset{n}{}}}{\displaystyle \underset{s=1}{\overset{r}{}}}\left[C^sF_s(u)\right]^{\frac{8\delta _{iI_s}}{<U_s,U_s>}}\text{e}^{2Q^iu+2P^i}ds_i^2\}.`$ 2. The dilatonic scalar fields are the following $$\phi ^\alpha (u)=\underset{s=1}{\overset{r}{}}\frac{4\chi _s\lambda _{a_s}^\alpha }{<U_s,U_s>}\mathrm{ln}[C^sF_s(u)]+uQ^{n+\alpha }+P^{n+\alpha },\alpha =1,\mathrm{},\omega .$$ (3.29) 3. For scalar functions $`\dot{\mathrm{\Phi }}^s(u)`$ we get $$\dot{\mathrm{\Phi }}^s(u)=Q_s\text{e}^{<U_s,x>}=Q_s\underset{s^{}=1}{\overset{m}{}}\left[C^s^{}F_s^{}(u)\right]^{C_{ss^{}}},s=1,\mathrm{},r.$$ (3.30) The corresponding $`F^a`$-field forms look as follows $$F^{(a_s,e,I_s)}=\dot{\mathrm{\Phi }}^{(a_s,e,I_s)}du\tau (I_s)$$ (3.31) for the electrically charged $`p`$-brane and $$F^{(a_s,m,I_s)}=Q_s\tau _0\tau _{i_1}\mathrm{}\tau _{i_c},\{i_1,\mathrm{},i_c\}=\overline{I_s},c=n_{a_s}d_0$$ (3.32) for the $`p`$-brane with magnetic-type charge. We put $`F^{(a_s,m,I_s)}=Q_s\tau _0`$ if $`\overline{I_s}=\mathrm{}`$. We stress that if $`r\|S_{}\|=n+\omega `$, then one must put $`Q^A=P^A=0,A=0,\mathrm{},n+\omega `$ in this solution. ## 4 General solutions for models associated with Lie algebras Now we list general solutions to the set of equations (3.23) for some special matrices $`(C_{ss^{}})`$. 1. $`(C_{ss^{}})=\text{diag}(2,\mathrm{},2)`$ is the Cartan matrix of the semi-simple Lie algebra $`A_1\mathrm{}A_1`$ of rank $`r`$. In this case the set of the characteristic vectors $`U_s`$ is orthogonal. $`F_s(u)\text{e}^{q^s(u)}={\displaystyle \frac{\mathrm{sin}[w_s(uu_{01})]}{w_s}},s=1,\mathrm{},r,`$ (4.1) $`E_T={\displaystyle \underset{s=1}{\overset{r}{}}}{\displaystyle \frac{2w_s^2}{<U_s,U_s>}},`$ (4.2) where $`w_s`$ are arbitrary constants, which may be real (including zero) or imaginary. 2. $`(C_{ss^{}})=(2\delta _{ss^{}}\delta _{s,s^{}+1}\delta _{s,s^{}1})`$ is the Cartan matrix of the simple Lie algebra $`A_rsl(r+1,C)`$. In this case all characteristic vectors $`U_s`$ are space-like with coinciding lengths, i.e. $$<U_s,U_s>U^2,s=1,\mathrm{},r.$$ (4.3) By the transformation $$q^sq^s\frac{\pi i}{2}m_s,$$ (4.4) where $$m_s=2\underset{s^{}=1}{\overset{r}{}}C^{ss^{}}=s(r+1s),$$ (4.5) we put the set of equations (3.21) into the form $$\ddot{q^s}=\mathrm{exp}\left[\underset{s^{}=1}{\overset{r}{}}C_{ss^{}}q^s^{}\right],s,=1,\mathrm{},r.$$ (4.6) These are precisely the $`A_r`$ Toda equations . Using the general solutions to these equations presented by Anderson , we obtain the following result $`F_s(u)=i^{s(r+1s)}{\displaystyle \underset{\mu _1<\mathrm{}<\mu _s}{\overset{r+1}{}}}v_{\mu _1}\mathrm{}v_{\mu _s}\mathrm{\Delta }^2(\mu _1,\mathrm{},\mu _s)\text{e}^{(w_{\mu _1}+\mathrm{}+w_{\mu _s})u},`$ (4.7) $`E_T={\displaystyle \frac{1}{2}}{\displaystyle \underset{\mu =1}{\overset{r+1}{}}}w_\mu ^2.`$ (4.8) where $`\mathrm{\Delta }^2(\mu _1,\mathrm{},\mu _s)`$ denotes the square of the Vandermonde determinant $$\mathrm{\Delta }^2(\mu _1,\mathrm{},\mu _s)=\underset{\mu _i<\mu _j}{}\left(w_{\mu _i}w_{\mu _j}\right)^2,\mathrm{\Delta }^2(\mu _1)1.$$ (4.9) The constants $`v_\mu `$ and $`w_m,\mu =1,\mathrm{},r+1`$, have to satisfy the following constraints: $$\underset{\mu =1}{\overset{r+1}{}}v_\mu =\mathrm{\Delta }^2(1,\mathrm{},r+1),\underset{\mu =1}{\overset{r+1}{}}w_\mu =0.$$ (4.10) The constants $`v_\mu ,w_\mu `$ are in general complex. There are additional constraints on them if one requires the functions $`F_s(u)`$ and the integral of motion (4.8) to be real. In (4.7) we used $`m_s=s(r+1s)`$ for $`A_r`$. 3. $`(C_{ss^{}})`$ is the following matrix $`C_{ss^{}}=\{\begin{array}{ccc}2\delta _{ss^{}}\delta _{s,s^{}+1}\delta _{s,s^{}1}\hfill & \mathrm{for}& s=1,\mathrm{},r,s^{}=1,\mathrm{},r1,\hfill \\ \delta _{ss^{}}\delta _{s,s^{}1}\hfill & \mathrm{for}& s=1,\mathrm{},r,s^{}=r.\hfill \end{array}`$ (4.13) The Cartan matrix of the simple Lie algebra $`B_rso(2r+1)`$ may be obtained from that given in (4.13) by multiplying the last column of $`(C_{sr})`$ by 2. In this case the general solution to the set of equations (3.21) may be obtained from the previous formulae (4.7), (4.8) as in . Notice that the equations (3.21) are symmetric under the following permutations $`q^sq^{r+1s}`$ for $`s=1,\mathrm{},r`$ if $`(C_{ss^{}})`$ if the Cartan matrix of $`A_r`$. This implies that there are solutions (4.7) with $`q^sq^{r+1s}`$ for $`s=1,\mathrm{},r`$. Moreover, this identification for $`r=4,6,8,\mathrm{}`$ leads to the $`(r/2)`$ equations of the form (3.21) with the matrix (4.13). Consequently the general solution of the equations (3.21) with the matrix (4.13) for $`r=r_0`$ may be obtained form (4.7) for $`r=2r_0`$ by putting additional constraints on the constants $`v_\mu ,w_\mu `$ providing with the identities $`F_s(u)F_{2r_0+1s}(u),s=1,\mathrm{},2r_0`$. 4. $`(C_{ss^{}})`$ is the Cartan matrix of the simple Lie algebra $`C_rsp(r,C)`$, i.e. $`C_{ss^{}}=\{\begin{array}{ccc}2\delta _{ss^{}}\delta _{s1,s^{}}\delta _{s1,s^{}}\hfill & \mathrm{for}& s=1,\mathrm{},r1,s^{}=1,\mathrm{},r,\hfill \\ \delta _{ss^{}}\delta _{s,s^{}1}\hfill & \mathrm{for}& s=r,s^{}=1,\mathrm{},r.\hfill \end{array}`$ (4.16) In this case the general solution to the set of equations (3.21) with $`r=r_0`$ may be obtained from (4.7) with $`r=2r_01`$ by putting additional constraints on the constants $`v_\mu ,w_\mu `$ providing with identities $`F_s(u)F_{2r_0s}(u),s=1,\mathrm{},2r_01`$. It stems ¿from the following property of the set (3.21): the identification $`q^sq^{2r_0s}`$ for $`s=1,\mathrm{},2r_01`$ reduces the set (3.21) with the Cartan matrix of $`A_{2r_01}`$ to the set (3.21) with the Cartan matrix of $`C_{r_0}`$ ($`r_02`$). ## 5 The particular solution describing black holes Suppose the nondegenerate Cartan-type matrix $`(C_{ss^{}})`$ satisfies the conditions $$m_s=2\underset{s^{}=1}{\overset{r}{}}C^{ss^{}}>0,s=1,\mathrm{},r.$$ (5.1) The conditions are valid for extremely large class of the $`p`$-brane models. For instance, the parameters $`m_s`$ are natural numbers if $`(C_{ss^{}})`$ is the Cartan matrix of a semi-simple Lie algebra $`G`$ . For $`G=A_r=sl(r+1,C)`$ the parameters $`m_s`$ are given by (4.5). Under the conditions (5.1) the set of equations (3.23) admits the following particular solution $$F_s(u)=a_s\left(\frac{\mathrm{sinh}[\overline{\mu }(uu_{01})]}{\overline{\mu }}\right)^{m_s},s=1,\mathrm{},r,$$ (5.2) where the constants $`a_s`$ are defined by $$a_s=\underset{s^{}=1}{\overset{r}{}}(m_s^{})^{C^{ss^{}}}$$ (5.3) and the constants $`\overline{\mu },u_{01}`$ are arbitrary. The corresponding to (5.2) integral of motion (3.10) has the form $$E_T=\underset{s,s^{}=1}{\overset{r}{}}\frac{C_{ss^{}}}{<U_s,U_s>}\frac{\dot{F}_s}{F_s}\frac{\dot{F}_s^{}}{F_s^{}}\underset{s=1}{\overset{r}{}}\frac{2}{<U_s,U_s>}\underset{s^{}=1}{\overset{r}{}}(F_s^{})^{C_{ss^{}}}=2\overline{\mu }^2\underset{s=1}{\overset{r}{}}\frac{m_s}{<U_s,U_s>}.$$ (5.4) For $`\overline{\mu }=0`$ the formulas (5.2),(5.4) read $`F_s(u)`$ $`=`$ $`a_s(uu_{01})^{m_s},s=1,\mathrm{},r,`$ (5.5) $`E_T`$ $`=`$ $`0.`$ (5.6) It is evident that (5.5) represents the polynomials in the radial coordinate $`u`$ if the parameters $`m_s`$ are natural numbers. As we have already mentioned, $`m_s`$ are natural numbers if, for instance, $`(C_{ss^{}})`$ is the Cartan matrix of a semi-simple Lie algebra. The formula (5.5) does not exhaust all possible polynomial solutions to the set (3.23). As far as we know, an explicit general form of the polynomial solution, which appears for the set of equations (3.23) with arbitrary natural numbers $`m_s`$ and vanishing integral of motion (5.4), is not found. There are only few examples in the literature. For instance, in all possible polynomial solutions were obtained for the matrix $`(C_{ss^{}})`$ supposed to be the Cartan matrix of the Lie algebras $`A_rsl(r+1,C)`$ with $`r=1,2,3`$. It is easy to check that an arbitrary polynomial solution to (3.23) under the condition $`E_T=0`$ may be obtained from (5.5) by adding some polynomial of lower degree to the leading term $`a_su^{m_s}`$. Now we use the particular solution (5.2) with $`\overline{\mu }>0`$ and its special form (5.5) corresponding to $`\overline{\mu }=0`$ for constructing non-extremal and extremal black holes, respectively. Consider the general form of exact solution (3.28)-(3.31) under the following additional assumptions 1. We put $$Q^A=\overline{\mu }\left(\frac{2V_0^A}{<V_0,V_0>}+\underset{s=1}{\overset{r}{}}\frac{2m_sU_s^A}{<U_s,U_s>}\delta _1^A\right),A=0,\mathrm{},n+\omega .$$ (5.7) One may verify that the conditions $`<Q,V_0>=0,<Q,U_s>=0,s=1,\mathrm{},r`$ are valid. Using the zero-energy constraint (3.6) and (5.4), we find the constant $`E_0`$ $$E_0=E_T\frac{1}{2}<Q,Q>=\frac{1}{2}d_0\overline{d}_0\mu ^2,$$ (5.8) where we denoted $$\mu =\overline{\mu }/\overline{d}_0,\overline{d}_0=d_01.$$ (5.9) 2. We take the parameters $`P^A`$ in the form $$\left(P^A\right)=\underset{s=1}{\overset{r}{}}\mathrm{ln}[C^sF_s(u_0)]\frac{2\left(U_s^A\right)}{<U_s,U_s>}u_0\left(Q^A\right)+\mathrm{ln}\epsilon _0(1,0,\underset{n1}{\underset{}{R^2,\mathrm{},R^n}},\underset{\omega }{\underset{}{0,\mathrm{},0}}),$$ (5.10) where $`\epsilon _0`$ is an arbitrary positive constant. The conditions $`<P,V_0>=0`$ lead to the following constraint on parameters $`R^2,\mathrm{},R^n`$ $$\underset{i=2}{\overset{n}{}}R^id_i=\overline{d}_0.$$ (5.11) Combining the conditions $`<P,U_s>=0,s=1,\mathrm{},r`$ and (3.13), we get $$Q_s^2=\frac{<U_s,U_s>}{4}\epsilon _0^{2_{i=2}^nd_i\delta _{iI_s}R^i}\underset{s^{}=1}{\overset{r}{}}[F_s^{}(u_0)]^{C_{ss^{}}},s=1,\mathrm{},r.$$ (5.12) 3. Now we consider the solution for $`u(u_0,+\mathrm{})`$ and introduce the following new radial coordinate $$R=\frac{R_0}{1\mathrm{exp}[2\overline{\mu }(uu_0)]}=\epsilon _0\left(\frac{2\mu }{1\mathrm{exp}[2\overline{\mu }(uu_0)]}\right)^{1/\overline{d}_0},R>R_0.$$ (5.13) The constant $`R_0`$ is defined by $$R_0=\epsilon _0(2\mu )^{1/\overline{d}_0}.$$ (5.14) Here we take the constant $`u_{01}`$ in (5.2) such that $`(u_0u_{01})>0`$. Moreover we introduce the constant $$R_{}=R_{u=2u_0u_{01}}=\epsilon _0\left(\frac{2\mu }{1\mathrm{exp}[2\overline{\mu }(u_0u_{01})]}\right)^{1/\overline{d}_0}>R_0.$$ (5.15) Finally, we obtain the metric $`ds^2`$ $`=`$ $`[1+\left({\displaystyle \frac{R_{}}{R}}\right)^{\overline{d}_0}\left({\displaystyle \frac{R_0}{R}}\right)^{\overline{d}_0}]^{_{s=1}^r\frac{8m_sd(I_s)}{(D2)<U_s,U_s>}}\{.{\displaystyle \frac{dR^2}{1(R_0/R)^{\overline{d}_0}}}+R^2d\mathrm{\Omega }_{d_0}^2`$ (5.16) $``$ $`\left[1+\left({\displaystyle \frac{R_{}}{R}}\right)^{\overline{d}_0}\left({\displaystyle \frac{R_0}{R}}\right)^{\overline{d}_0}\right]^{_{s=1}^r\frac{8m_s}{<U_s,U_s>}}\left(1\left({\displaystyle \frac{R_0}{R}}\right)^{\overline{d}_0}\right)dt^2`$ $`+`$ $`{\displaystyle \underset{i=2}{\overset{n}{}}}\epsilon _0^{2R^i}[1+\left({\displaystyle \frac{R_{}}{R}}\right)^{\overline{d}_0}\left({\displaystyle \frac{R_0}{R}}\right)^{\overline{d}_0}]^{_{s=1}^r\frac{8m_s\delta _{iI_s}}{<U_s,U_s>}}ds_i^2.\},`$ the dilatonic scalar fields $$\phi ^\alpha =\underset{s=1}{\overset{r}{}}\frac{4m_s\chi _s\lambda _{a_s}^\alpha }{<U_s,U_s>}\mathrm{ln}\left[1+\left(\frac{R_{}}{R}\right)^{\overline{d}_0}\left(\frac{R_0}{R}\right)^{\overline{d}_0}\right],\alpha =1,\mathrm{},\omega ,$$ (5.17) and the potential derivatives $`{\displaystyle \frac{d\mathrm{\Phi }^s}{dR}}=`$ $``$ $`\mathrm{sgn}(Q_s)2\overline{d}_0\sqrt{{\displaystyle \frac{m_s}{<U_s,U_s>}}\left(1\left({\displaystyle \frac{R_0}{R_{}}}\right)^{\overline{d}_0}\right)}`$ (5.18) $`\times `$ $`\epsilon _0^{_{i=2}^nd_i\delta _{iI_s}R^i}\left({\displaystyle \frac{R_{}}{R}}\right)^{\overline{d}_0}\left[1+\left({\displaystyle \frac{R_{}}{R}}\right)^{\overline{d}_0}\left({\displaystyle \frac{R_0}{R}}\right)^{\overline{d}_0}\right]^2.`$ The corresponding $`F^a`$-field forms may be obtained by (3.31),(3.32), where $$\|Q_s\|=2\overline{d}_0\sqrt{\frac{m_s}{<U_s,U_s>}\left(1\left(\frac{R_0}{R_{}}\right)^{\overline{d}_0}\right)}\epsilon _0^{_{i=2}^nd_i\delta _{iI_s}R^i}\left(\frac{R_{}}{\epsilon _0}\right)^{\overline{d}_0}.$$ (5.19) Then constants $`\overline{\mu },\epsilon _0,(u_0u_{01})`$ are independent. The constants $`\mu ,\overline{d}_0,R_0,R_{}`$ are defined by (5.9),(5.14),(5.15). The parameters $`R^2,\mathrm{},R^n`$ obey the relation (5.11). Now we analyze the particular solution (5.16)-(5.18). The metric (5.16) is asymptotically flat, i.e. $$\underset{R+\mathrm{}}{lim}ds^2=dR^2+R^2d\mathrm{\Omega }_{d_0}^2dt^2+\underset{i=2}{\overset{n}{}}\epsilon _0^{2R^i}ds_i^2.$$ (5.20) According to (5.11) all parameters $`R^2,\mathrm{},R^n`$ may be positive. Then, the constant scale factors $`\epsilon _0^{2R^i}`$ of internal spaces $`M_2,\mathrm{},M_n`$ are arbitrary small if $`\epsilon _0`$ is large enough. If $`\overline{\mu }>0`$ the particular solution describes a non-extremal black hole with the horizon at $`R=R_0`$. The active gravitational mass $`M_g`$ and the Hawking temperature $`T_H`$ of this black hole read $`2G_NM_g=R_{}^{\overline{d}_0}\left[\left(1\left({\displaystyle \frac{R_0}{R_{}}}\right)^{\overline{d}_0}\right){\displaystyle \underset{s=1}{\overset{r}{}}}{\displaystyle \frac{4m_sU_s^1}{<U_s,U_s>}}+\left({\displaystyle \frac{R_0}{R_{}}}\right)^{\overline{d}_0}\right],`$ (5.21) $`T_H={\displaystyle \frac{\overline{d}_0}{4\pi k_BR_{}}}\left({\displaystyle \frac{R_0}{R_{}}}\right)^{_{s=1}^r\frac{4m_s}{<U_s,U_s>}1},`$ (5.22) where $`G_N`$ and $`k_b`$ are Newton’s gravitational constant and Boltzmann’s constant, respectively. The solution (5.16)-(5.18) may be considered in the so-called extreme case, when $`\overline{\mu }=0`$ ($`R_0=0`$). It follows from general statements proved in that the point $`R=0`$ is a curvature singularity of the metric (5.16) with $`R_0=0`$ if $`T_H+\mathrm{}`$ as $`R_0+0`$. Then, the particular solution admits an extremal black hole only if $$\underset{s=1}{\overset{r}{}}\frac{4m_s}{<U_s,U_s>}1.$$ (5.23)
warning/0004/math-ph0004004.html	ar5iv	text	# 1 Introduction ## 1 Introduction Since the 1995-observations of Bose-Einstein condensation (BEC) in trapped alkali gases, an intense renewed interest is going on in the research of the physical properties and the nature of Bose condensed systems. In particular the interference pattern between two overlappping condensates has been measured, see e.g. and many other recent experimental settings and results. In this context of interference, the static and dynamic properties of the phase of the condensate are of major importance. This has been the subject of many theoretical studies all over the last years. As a primordial and old question, the very existence of the phase and/or the phase operator, comes into the picture again. One encounters continuous efforts to formulate the phase (operator) in the standard theory of BEC, which we could call the Bogoliubov-Hartree-Fock theory , in a system with a finite number of atoms (see e.g. ). One is constantly assuming that the condensation is occuring into a coherent state of the lowest energy mode of the system. Such a state fixes a well defined phase and amplitude, but should in stead exhibit inevitable fluctuations of both these quantities. Or, one is fixing the number of atoms in the system, i.e. the condensation takes place in a number state of Fock space, excluding any atom number fluctuations. Although these basic theoretical difficulties are now getting ripe in the minds of many researchers in the field, all kinds of procedures and tricks are permanently invented to wave away these difficulties. In this paper we take the point of view that these questions about the character of the quantum state into which the condensation occurs, and its major properties, are nevertheless of major importance. By now it is indeed well known that a condensate state is neither a ‘pure Fock’, nor a ‘pure coherent’ state in the strict mathematical sense, nor in the physical sense. As explained above, ‘simple’ coherent states or Fock states lead to annoying technical difficulties in order to describe and understand the essentials of many of the experimental challenging measurements on BEC which are constantly performed. After all, condensation is up to now, only clearly defined and generally accepted for homogeneous systems. Of course, we are aware of different tentatives to introduce decent thermodynamic limits for trapped gases. With all this knowledge in mind, we focus our attention here, not on the situation of BEC in trapped gases, but on the phenomenon of BEC for homogeneous systems, where one has a well defined thermodynamic limit, and where the occurence of BEC, accompanied by a spontaneous $`U(1)`$ \- symmetry breaking is well understood. Furthermore we take into account that the main entries of the theory of the Bose condensates and their interference patterns are the particle number fluctuations and the phase operator fluctuations. The main question is here, can one define rigorously a phase operator fluctuation and a particle number operator fluctuation of the condensate? The answer is proved to be positive. It is based on the notion of fluctuation operator which was introduced in a mathematically rigorous framework some time ago . We realise however that these results did not reach so far the majority of the theoretical physics community. The aim for introducing the notion of fluctuation operator, was precisely to study the quantum effects on the level of the fluctuations. We applied this theory of quantum fluctuation operators already in order to derive exactly rigorous results on BEC for the Bogoliubov-Hartree-Fock model . In section 2, we describe these results in a language approachable for non-mathematics minded readers. The celebrated phase operator is nothing but the canonical fluctuation operator conjugate to the number operator fluctuation operator. This section is not just for warming up, but it should also shed a different light, than one is used to, on the status of the existence and meaning of the phase operator. In section 3, we study a model of two Bose-Einstein condensates with a Josephson-like coupling. We look at the static and thermodynamic properties of this model of two condensates. As far as we know, the problem of Josephson oscillations between coupled Bose-Einstein condensates has not yet been considered in a mathematically rigorous manner. Here we present a solvable model in a full quantum field theoretical setting and in section 4 and 5 we give a full description and analysis of the dynamical equations of the total and relative number operator fluctuations and the total and relative phase operator fluctuations. About the dynamics we find the exact oscillatory time behaviour of all these fluctuation operators. We also detect the so-called collapse and revival phenomena. Our work is, as far as we know, the first rigorous one on this topic for homogeneous systems, it should also put in a new perspective much of the discussions which are going on in the large activity dealing with trapped Bosons and their interference patterns (see e.g. ). ## 2 Number and phase fluctuation operators In order to fix our ideas and in view of the model we describe in section 3 for the study of the phase interference between two condensates, we present the number and phase fluctuation operators for the imperfect Bose gas (or mean field Bose gas) . We follow the lines of but it should be clear that its validity is much larger . The leading idea of this section is to make clear that the up to now rather ‘mysterious’ phase operator, which everybody uses in the field, but about which there are doubts on its very existence, has a firm mathematical definition in an equilibrium condensed state of a Bose gas. It should be realised that such a condensed state is neither a coherent, nor a Fock state in the technical narrow sense. It is defined as the fluctuation operator canonically adjoint to the number fluctuation operator. This definition is not just a formal thing, but it gives a physical interpretation of the phase operator, different from the existing ones. We are not repeating here all the mathematics of the definition of the phase fluctuation operator, which can be found in various papers (see e.g. ). We content ourself here in making these definitions plausible for the imperfect Bose gas. Let $`\mathrm{\Lambda }R^3`$ be the centered cubic box of length $`L`$, with periodic boundary conditions. The Boson creation and annihilation operators in the one-particle state $`\psi _{L,k}(x)=V^{1/2}e^{ik.x}`$, $`x\mathrm{\Lambda }`$, $`k\mathrm{\Lambda }^{}=\frac{2\pi }{L}Z^\nu `$ are given by $$a_{L,k}^{}=\frac{1}{\sqrt{V}}_\mathrm{\Lambda }𝑑xa^{}(x)e^{ik.x}\text{and}a_{L,k}=\frac{1}{\sqrt{V}}_\mathrm{\Lambda }𝑑xa(x)e^{ik.x},$$ with $$[a(x),a^{}(y)]=\delta (xy).$$ The imperfect Bose gas is specified by the local Hamiltonian $`H_L`$ : $$H_L=T_L\mu _LN_L+\frac{\lambda }{2V}N_L^2,$$ (1) where $`T_L`$ $`={\displaystyle \underset{k\mathrm{\Lambda }^{}}{}}ϵ_ka_{L,k}^{}a_{L,k},ϵ_k={\displaystyle \frac{\|k\|^2}{2m}}`$ $`N_L`$ $`={\displaystyle \underset{k\mathrm{\Lambda }^{}}{}}a_{L,k}^{}a_{L,k};`$ $`\lambda >0`$ measures the strength of the mean field inter particle repulsion. This model is exactly soluble in the thermodynamic limit $`L\mathrm{}`$, keeping the particle density in the Gibbs state $`\omega _L()`$ for (1) constant, equal to $`\rho `$, i.e. for all L: $$\frac{1}{V}\omega _L(N_L)=\rho .$$ It is proved rigorously that for $`T<T_c`$ or $`\rho `$ large enough, the limit state $`\omega _\beta ()=lim_L\mathrm{}\omega _L()`$ exists as an integral over ergodic states ($`\omega _\beta ^\alpha `$, $`\alpha [0,2\pi ]`$): $$\omega _\beta ()=\frac{1}{2\pi }_0^{2\pi }𝑑\alpha \omega _\beta ^\alpha (),$$ with $$\omega _\beta ^\alpha \left(e^{i[a^{}(f)+a(f)]}\right)=e^{\frac{1}{2}(f,Kf)+2i\sqrt{\rho _0}\|\widehat{f}(0)\|\mathrm{cos}\alpha },$$ where $$\widehat{Kf}(k)=\frac{1}{2}\widehat{f}(k)\mathrm{coth}\frac{\beta ϵ_k}{2},fL^2(R^\nu ).$$ The spontaneous gauge symmetry breaking accompanying the phase transition is visible in the states $`\omega _\beta ^\alpha `$ having the property: $$\underset{L\mathrm{}}{lim}\omega _\beta ^\alpha \left(\frac{a_{L,0}^{}}{\sqrt{V}}\right)=\sqrt{\rho _0}e^{i\alpha }.$$ Clearly, $`\rho _0`$ is the condensate and $`\alpha `$ is the phase of the order parameter. One also proves that in the state $`\omega _\beta ^\alpha `$, one has the operator limit: $$\underset{L\mathrm{}}{lim}\frac{a_{L,0}^{}}{\sqrt{V}}=\sqrt{\rho _0}e^{i\alpha }.$$ From now on we limit our attention to one of the ergodic states $`\omega _\beta ^\alpha `$ for some fixed $`\alpha `$, and without restriction of generality we take $`\alpha =0`$, and with a condensate density $`\rho _00`$. For simplicity, denote the state $`\omega _\beta ^0`$ by $`\omega `$. The state $`\omega `$ does not have the gauge symmetry. The generator of the gauge symmetry is the number operator $$N_L=_\mathrm{\Lambda }𝑑xa^{}(x)a(x)$$ with local number density operator $`n(x)=a^{}(x)a(x)`$. The common choice of order parameter operator is $`V^{1/2}a_{L,0}^{\mathrm{}}`$, or taking a self-adjoint combination: $$O_L=\frac{i}{\sqrt{V}}(a_{L,0}^{}a_{L,0})=\frac{i}{V}_\mathrm{\Lambda }𝑑x\left(a^{}(x)a(x)\right),$$ with local order parameter density operator $`o(x)=i\left(a^{}(x)a(x)\right)`$. We concentrate now on the $`k`$-mode fluctuations, with $`k0`$, of the local number and order parameter density operators, i.e. on $`F_{L,k}(n)`$ $`={\displaystyle \frac{1}{\sqrt{V}}}{\displaystyle _\mathrm{\Lambda }}𝑑x\left(n(x)\omega (n(x))\right)\mathrm{cos}k.x`$ $`F_{L,k}(o)`$ $`={\displaystyle \frac{1}{\sqrt{V}}}{\displaystyle _\mathrm{\Lambda }}𝑑x\left(o(x)\omega (o(x))\right)\mathrm{cos}k.x.`$ Remark that for all finite $`L`$, the quantities $`F_{L,k}(n)`$ and $`F_{L,k}(o)`$ are operators and do represent the fluctuations of the number density and of the order parameter density. The first tedious question that is posed, is to characterize the limit operators: $`F_k(n)`$ $`=\underset{L\mathrm{}}{lim}F_{L,k}(n)`$ $`F_k(o)`$ $`=\underset{L\mathrm{}}{lim}F_{L,k}(o).`$ The details of the proof of these limits can be found in . Here we just mention that the limits are taken in the sense of a central limit theorem. The main result is that the limits $`F_k(n)`$ and $`F_k(o)`$ are operators on a well specified Hilbert space, $`\stackrel{~}{}_k`$, generated by a normalised vector $`\stackrel{~}{\mathrm{\Omega }}_k`$ and vectors $`F_k(A_1)\mathrm{}F_k(A_n)\stackrel{~}{\mathrm{\Omega }}_k`$, with the $`A_i`$ local operators, like e.g. $`n(x)`$ and $`o(x)`$, and with arbitrary $`n`$; the scalar product of $`\stackrel{~}{}_k`$ is given by $$(F_k(A_1)\mathrm{}F_k(A_n)\stackrel{~}{\mathrm{\Omega }}_k,F_k(B_1)\mathrm{}F_k(B_n)\stackrel{~}{\mathrm{\Omega }}_k)=\delta _{n,m}\text{Perm}\left((\stackrel{~}{\mathrm{\Omega }}_k,F_k(A_i)F_k(B_j)\stackrel{~}{\mathrm{\Omega }}_k)\right)_{i,j}$$ (2) with two-point function given by $$(\stackrel{~}{\mathrm{\Omega }}_k,F_k(A_i)F_k(B_j)\stackrel{~}{\mathrm{\Omega }}_k)=\underset{L\mathrm{}}{lim}\omega \left(F_{L,k}(A_i)F_{L,k}(B_j)\right)$$ (3) i.e. essentially determined by the two-point functions of the given state $`\omega `$. Therefore as is clear from (2), all $`(n+m)`$-point functions are given by the two-point function (3). The definition (2) defines completely the fluctuation operators $`F_k(A)`$ on the Hilbert space $`\stackrel{~}{}_k`$. On the other hand, the non-commutative law of large numbers, here of large operators, leads straightforwardly to the canonical commutation relation $$\underset{L\mathrm{}}{lim}[F_{L,k}(n),F_{L,k}(o)]=\underset{L\mathrm{}}{lim}\frac{1}{2V}_\mathrm{\Lambda }𝑑x[n(x),o(x)]=\omega ([n(x),o(x)])=i\sqrt{\rho _0},$$ or $$[F_k(n),F_k(o)]=i\sqrt{\rho _0}.$$ (4) This is the basic result for the definition of the phase operator. Equation (4) means that the number fluctuation operator $`F_k(n)`$ and the order parameter fluctuation operator $`F_k(o)`$ are canonically conjugate (compare with $`[q,p]=i\mathrm{}`$). Hence for the physics of BEC we found on the level of fluctuations the canonical pair $`(F_k(n),F_k(o))`$. Clearly the operator $`F_k(o)`$ satisfies all basic physical requirements for playing the role of what is usually called the phase operator of the condensate. The reader will have recognised from (2) and (4) that the fluctuation operators $$F_k(A),F_k(B),\mathrm{}$$ form an algebra of Boson field operators and in particular that $`F_k(n)`$ and $`F_k(o)`$ form quantum canonical variables with a quantisation parameter $`\sqrt{\rho _0}`$ (compare with $`\mathrm{}`$). On the other hand, from (2) it is clear that the vector $`\stackrel{~}{\mathrm{\Omega }}_k`$ defines a generalised or quasi free (gaussian) state $$\stackrel{~}{\omega }_k()=(\stackrel{~}{\mathrm{\Omega }}_k,\stackrel{~}{\mathrm{\Omega }}_k)$$ on the Boson field algebra (see ). This means that the central limit theorem for the $`k`$-mode fluctuations in the state $`\omega `$ defines an equilibrium state $`\stackrel{~}{\omega }_k`$ on the fluctuation operators. The mentioned quasi free character means that all correlation functions of limit fluctuation operators are polynomial functions only of the one- and two-point functions. For more details we refer once more to . The reader might ask for the unicity of this phase operator, if being defined only as the canonically adjoint operator to the number fluctuation operator. The interested reader is referred to section 5 and for that discussion. ## 3 The model and equilibrium states We consider two coupled Bose-Einstein condensates, each of them modelled by an imperfect or mean field Bose gas. Denote $`a_i^{\mathrm{}}(x)`$, $`i=1,2`$ the creation and annihilation operators for the two Bose gases, i.e. $$[a_i(x),a_j^{}(y)]=\delta _{i,j}\delta (xy).$$ We assume that the two gases have the same particle density $`\frac{\rho }{2}`$ (hence $`\rho `$ is the total particle density $`\rho =\frac{1}{V}\omega _L(N_L)`$ with $`N_L=N_{1,L}+N_{2,L}`$), and also that they are of the same type of particles (i.e. there is only one mean field constant $`\lambda `$). We also assume a phase difference $`\phi `$ between the gases, and model the Josephson coupling between the gases by a term $$C_L^{1,2}=\gamma \underset{k\mathrm{\Lambda }^{}}{}a_{1,k}^{}a_{2,k}e^{i\phi }+a_{2,k}^{}a_{1,k}e^{i\phi },$$ (5) with $`\gamma >0`$ the coupling constant. Hence the local Hamiltonian of the system we study is given by: $`H_L`$ $`=T_{1,L}+T_{2,L}\mu _LN_L+{\displaystyle \frac{\lambda }{2V}}N_L^2+C_L^{1,2}`$ $`={\displaystyle \underset{k\mathrm{\Lambda }^{}}{}}(ϵ_k\mu _L)(a_{1,k}^{}a_{1,k}+a_{2,k}^{}a_{2,k})+{\displaystyle \frac{\lambda }{2V}}(N_{1,L}+N_{2,L})^2`$ $`\gamma {\displaystyle \underset{k\mathrm{\Lambda }^{}}{}}a_{1,k}^{}a_{2,k}e^{i\phi }+a_{2,k}^{}a_{1,k}e^{i\phi }.`$ (6) In this section we find the limiting Gibbs states $`\omega =lim_L\mathrm{}\omega _L`$ at inverse temperature $`\beta `$ of this model. A rigorous study along the lines of is perfectly possible, but we permit ourselves here a more intuitive approach. As in any mean field model, we replace the Hamiltonian (3) by a state dependent effective Hamiltonian $$H_L^\omega =\underset{k\mathrm{\Lambda }^{}}{}(ϵ_k\mu +\lambda \rho )(a_{1,k}^{}a_{1,k}+a_{2,k}^{}a_{2,k})\gamma \underset{k\mathrm{\Lambda }^{}}{}a_{1,k}^{}a_{2,k}e^{i\phi }+a_{2,k}^{}a_{1,k}e^{i\phi }$$ (7) with $`\mu =lim_L\mathrm{}\mu _L`$ in correspondence with the constraint $`\rho =\frac{1}{V}\omega (N_L)`$. This effective Hamiltonian is bilinear in the creation and annihilation operators and therefore it can be diagonalized by a Bogoliubov transformation. Let $`\delta _L^\omega ()=[H_L^\omega ,]`$ be the generator of this dynamics and $`f_k=ϵ_k\mu +\lambda \rho `$. Then $$\delta _L^\omega \left(\begin{array}{c}a_{1,k}^{}\\ a_{2,k}^{}\end{array}\right)=\left(\begin{array}{cc}f_k& \gamma e^{i\phi }\\ \gamma e^{i\phi }& f_k\end{array}\right)\left(\begin{array}{c}a_{1,k}^{}\\ a_{2,k}^{}\end{array}\right).$$ The matrix $$\left(\begin{array}{cc}f_k& \gamma e^{i\phi }\\ \gamma e^{i\phi }& f_k\end{array}\right)$$ has eigenvalues $`E_k^\pm `$, $$E_k^\pm =f_k\pm \gamma ,$$ with corresponding eigenoperators $`b_{\pm ,k}^{}`$, $$b_{\pm ,k}^{}=\frac{1}{\sqrt{2}}\left(a_{1,k}^{}e^{\frac{i}{2}\phi }a_{2,k}^{}e^{\frac{i}{2}\phi }\right),$$ (8) i.e. $$\delta _L^\omega (b_{\pm ,k}^{})=E_k^\pm b_{\pm ,k}^{}.$$ (9) The $`b_\pm ^{\mathrm{}}(x)`$ still satisfy Boson commutation rules. The energy spectrum of the quasi-particles $`b_\pm ^{\mathrm{}}(x)`$ has two branches, $`\{E_k^\pm \|kR^\nu \}`$. By a standard argument, using the Boson commutation rules and the correlation inequalities , characterizing the limit equilibrium states, $$\underset{L\mathrm{}}{lim}\beta \omega \left(X^{}\delta _L^\omega (X)\right)\omega (X^{}X)\mathrm{ln}\frac{\omega (X^{}X)}{\omega (XX^{})},$$ (10) for $`X`$ any polynomial in the creation and annihilation operators, one finds in a straightforward manner: $$\omega (b_{\pm ,k}^{}b_{\pm ,k})=\frac{1}{e^{\beta E_k^\pm }1}.$$ (11) Along the usual lines of the derivation of Bose-Einstein condensation, we find a critical density (or inverse temperature) above which one derives the following value of the chemical potential: $$\mu =\lambda \rho \gamma ,$$ i.e. $`E_k^{}=ϵ_k`$ and $`E_k^+>2\gamma `$ for all $`k`$. Hence there is a macroscopic occupation of the 0-momentum state of the ‘$``$’-mode possible; the condensate density is given by: $$\underset{V\mathrm{}}{lim}\frac{1}{V}\omega (b_{,0}^{}b_{,0})=\rho _0>0.$$ (12) There is no condensation for the ‘+’-mode since $`lim_{k0}E_k^+=2\gamma >0`$ for $`\mu =\lambda \rho \gamma `$. From $`\rho =lim_L\mathrm{}V^1\omega (N_L)`$ and (8), one finds $`\rho `$ $`=\underset{L\mathrm{}}{lim}{\displaystyle \frac{1}{V}}{\displaystyle \underset{k}{}}\omega \left(b_{,k}^{}b_{,k}\right)+\omega \left(b_{+,k}^{}b_{+,k}\right)`$ $`=\rho _0+{\displaystyle _{R^\nu }}{\displaystyle \frac{dk}{(2\pi )^3}}{\displaystyle \frac{1}{e^{\beta ϵ_k}1}}+{\displaystyle _{R^\nu }}{\displaystyle \frac{dk}{(2\pi )^3}}{\displaystyle \frac{1}{e^{\beta (ϵ_k+2\gamma )}1}}.`$ For $`\alpha 0`$, denote $$\rho (\alpha )=_{R^\nu }\frac{dk}{(2\pi )^3}\frac{1}{e^{\beta (ϵ_k+\alpha )}1},$$ (13) then (12) becomes $$\rho _0=\rho \rho (0)\rho (2\gamma ).$$ It is clear that $`\rho (0)+\rho (2\gamma )`$ is the critical density. The Bose-Einstein condensation (12) implies a spontaneous breaking of the gauge symmetry, i.e. the limiting Gibbs state $`\omega `$ decomposes with respect to the $`U(1)`$ gauge group into distinct extremal equilibrium states: $$\omega =\frac{1}{2\pi }_0^{2\pi }𝑑\theta \omega _\theta ,$$ where each of the states $`\omega _\theta `$ is determined by the two-point function (11) and the one-point function $$\omega _\theta \left(b_{}^{}(x)\right)=\sqrt{\rho _0}e^{i\theta }.$$ (14) This important point, whose non-triviality is often overlooked, is proved in . Of course, the other way round, namely that gauge symmetry breaking implies condensation, is well known and follows trivially from the Schwartz inequality. From now on we choose a particular extremal equilibrium state $`\omega _\theta `$, and without loss of generality we take $`\theta =0`$. For notational simplicity, this state is again denoted by $`\omega `$. Using once more the correlation inequalities (10), it is not difficult to show that the higher order correlations decompose into sums and products of one- and two-point correlations, given by (14) respectively (11), i.e. the state $`\omega `$ is quasi-free. Therefore we have completely characterized the equilibrium states. It is clear that the gauge symmetry breaking state under discussion indeed corresponds to a state of two different condensates interacting through a Josephson coupling (5). Since $`\omega (b_+^{}(x))=0`$ and $`\omega (b_{}^{}(x))=\sqrt{\rho _0}`$, we find using (8) $$\omega (a_1^{}(x))e^{\frac{i}{2}\phi }=\omega (a_2^{}(x))e^{\frac{i}{2}\phi }=\sqrt{\frac{\rho _0}{2}}.$$ (15) Notice that $`\phi `$ is indeed the phase difference between the condensates. Our arbitrary choice $`\theta =0`$ then actually determines both phases to be $`\pm \frac{\phi }{2}`$ and the choice of equal particle densities $`\frac{\rho }{2}`$ yields equal condensate densities $`\frac{\rho _0}{2}`$. Moreover it is obvious that the gapless mode $`E_k^{}`$ is related to the broken gauge symmetry, i.e. to the Bose-Einstein condensation, and that the mode $`E_k^+`$ with energy gap $`2\gamma `$ arises due to the presence of the Josephson coupling. The detailed study of the fluctuation operators corresponding to these two excitation branches and modes is the subject of the subsequent sections. ## 4 Total number and phase fluctuation operators Motivated by the discussion in section 2 and , define the total number and phase fluctuations in the box $`\mathrm{\Lambda }`$ by ($`k0`$) $`F_{L,k}(n_{tot})`$ $`={\displaystyle \frac{1}{\sqrt{V}}}{\displaystyle _\mathrm{\Lambda }}𝑑x\left(a_1^{}(x)a_1(x)+a_2^{}(x)a_2(x)\rho \right)\mathrm{cos}k.x`$ (16) $`F_{L,k}(\varphi _{tot})`$ $`={\displaystyle \frac{i}{\sqrt{V}}}{\displaystyle _\mathrm{\Lambda }}𝑑x\left(b_{}^{}(x)b_{}(x)\right)\mathrm{cos}k.x,`$ (17) where $`b_{}(x)`$ is defined in (8). Again on the basis of the law of large numbers: $$\underset{L\mathrm{}}{lim}[F_{L,k}(n_{tot}),F_{L,k}(\varphi _{tot})]=\underset{L\mathrm{}}{lim}\frac{i}{2\sqrt{V}}(b_{,0}^{}+b_{,0})=i\sqrt{\rho _0}.$$ (18) Using (8) one writes also $$a_1^{}(x)a_1(x)+a_2^{}(x)a_2(x)=b_{}^{}(x)b_{}(x)+b_+^{}(x)b_+(x),$$ or $$F_{L,k}(n_{tot})=F_{L,k}(n_{})+F_{L,k}(n_+),$$ with $`n_\pm `$ defined in the obvious sense. Hence the total number operator fluctuation is the sum of two number operator fluctuations of two imperfect Bose gases. For a single imperfect Bose gas (see section 2), one finds this analysis as the subject of , making the present analysis straightforward. As already discussed in section 2, and apparent from (18), the limiting fluctuation operators $`F_k()`$ satisfy Bosonic commutation rules (although the local fluctuation operators do not). Equation (18) learns also that the fluctuation operators $`F_k(n_{tot})`$ and $`F_k(\varphi _{tot})`$ constitute a canonical pair, generating an algebra of canonical commutation relations (CCR) of fluctuation observables of the system. Furthermore, it is shown in and briefly discussed in section 2 that the central limit theorem also fixes an equilibrium state $`\stackrel{~}{\omega }_k`$ on this algebra of limiting fluctuation operators, which is a CCR-algebra of Bosonic field operators. This state is shown to be quasi-free and gauge invariant, and hence completely determined by its two-point function, given by : $$\stackrel{~}{\omega }_k\left(F_k(A)F_k(B)\right)=\underset{L\mathrm{}}{lim}_\mathrm{\Lambda }𝑑z\omega \left(A(z)B(0)\right)\mathrm{cos}k.z,$$ where $`A,B`$ are (in the present case) polynomials in the microscopic canonical Bosonic field operators. Remark that there are no technical problems related to the central limit theorem for $`k0`$, off-diagonal long-range order correlations do appear only at $`k=0`$. The first step in our study of the total number- and phase fluctuation operators is to determine their variances. ###### Proposition 1. For $`k0`$ we have $`\begin{array}{cc}\hfill (i)\stackrel{~}{\omega }_k\left(F_k(n_{tot})^2\right)& =\stackrel{~}{\omega }_k\left(F_k(n_{})^2\right)+\stackrel{~}{\omega }_k\left(F_k(n_+)^2\right)\hfill \\ & ={\displaystyle \frac{\rho _0}{2}}\mathrm{coth}{\displaystyle \frac{\beta ϵ_k}{2}}+{\displaystyle \frac{1}{2}}{\displaystyle _{R^\nu }}{\displaystyle \frac{dp}{(2\pi )^3}}{\displaystyle \frac{e^{\beta ϵ_{p+k}}+e^{\beta ϵ_p}}{\left(e^{\beta ϵ_{p+k}}1\right)\left(e^{\beta ϵ_p}1\right)}}\hfill \\ & +{\displaystyle \frac{1}{2}}{\displaystyle _{R^\nu }}{\displaystyle \frac{dp}{(2\pi )^3}}{\displaystyle \frac{e^{\beta E_{p+k}^+}+e^{\beta E_p^+}}{\left(e^{\beta E_{p+k}^+}1\right)\left(e^{\beta E_p^+}1\right)}}\hfill \end{array}`$ (19) $`(ii)\stackrel{~}{\omega }_k\left(F_k(\varphi _{tot})^2\right)`$ $`={\displaystyle \frac{1}{2}}\mathrm{coth}{\displaystyle \frac{\beta ϵ_k}{2}}.`$ (20) ###### Proof. $`(i)`$ The fact that $`\stackrel{~}{\omega }_k\left(F_k(n_{tot})^2\right)=\stackrel{~}{\omega }_k\left(F_k(n_{})^2\right)+\stackrel{~}{\omega }_k\left(F_k(n_+)^2\right)`$ follows from the fact that the ‘$`+`$’- and ‘$``$’-mode are independent of each other. For a calculation of the explicit expression for $`\stackrel{~}{\omega }_k\left(F_k(n_{tot})^2\right)`$, see . $`(ii)`$ This is simply the two-point function of the state $`\omega `$, see . ∎ Remark that the result of this proposition, supplemented with (18) is sufficient in order to characterize completely the limiting fluctuation operators $`F_k(n_{tot})`$ and $`F_k(\varphi _{tot})`$ on a well defined Hilbert space (see section 2). We do not enter into these technical details. Of course we are interested particularly in the limit $`k0`$ of these operators, in order to see how the long-range order due to the Bose-Einstein condensation manifests itself on the level of the fluctuations. In one can find a discussion demonstrating that quantum effects will only be present in the limit $`k0`$ if one works in the groundstate $`\omega ^g`$, defined as the zero-temperature limit of the equilibrium state $`\omega `$: $$\omega ^g=\underset{\beta \mathrm{}}{lim}\omega .$$ At $`T=0`$, the variances (19) and (20) simplify to $`\stackrel{~}{\omega }_k^g\left(F_k(n_{tot})^2\right)`$ $`={\displaystyle \frac{\rho _0}{2}}`$ $`\stackrel{~}{\omega }_k^g\left(F_k(\varphi _{tot})^2\right)`$ $`={\displaystyle \frac{1}{2}},`$ and the limit $`k0`$ is trivial: $$F_0(n_{tot})=\underset{k0}{lim}F_k(n_{tot})\text{and}F_0(\varphi _{tot})=\underset{k0}{lim}F_k(\varphi _{tot}).$$ These are well defined fluctuation operators, satisfying $`[F_0(n_{tot}),F_0(\varphi _{tot})]`$ $`=i\sqrt{\rho _0}`$ $`\stackrel{~}{\omega }_0^g\left(F_0(n_{tot})^2\right)`$ $`={\displaystyle \frac{\rho _0}{2}}`$ $`\stackrel{~}{\omega }_0^g\left(F_0(\varphi _{tot})^2\right)`$ $`={\displaystyle \frac{1}{2}}.`$ We derived the exact uncertainty relation between the number operator and phase operator, given by $$\stackrel{~}{\omega }_0^g\left(F_0(n_{tot})^2\right)\stackrel{~}{\omega }_0^g\left(F_0(\varphi _{tot})^2\right)=\frac{\rho _0}{4}.$$ Remark that the condensate density $`\rho _0`$, in fact $`\sqrt{\rho _0}`$, acts in this equation, as well as in equation (18), as a quantisation parameter (compare with $`\mathrm{}`$). Consequently, the whole content of our results, as well as all physical interpretations, do disappear in the absence of condensation, i.e. if $`\rho _0=0`$. Finally remark that we have omitted in our notation the state dependence of the fluctuation operators throughout, although this dependence is important. Fluctuation operators corresponding to different states (e.g. corresponding to different temperature or different phase) in fact are not comparable as they act on completely different Hilbert spaces. We do not enter into these mathematical subtleties. ## 5 Relative number and phase fluctuation operators The relative number operator in a finite volume is: $$N_{L,rel}=N_{1,L}N_{2,L},$$ and its $`k`$-mode fluctuation $$F_{L,k}(n_{rel})=\frac{1}{V^{1/2}}_\mathrm{\Lambda }𝑑x\left(a_1^{}(x)a_1(x)a_2^{}(x)a_2(x)\right)\mathrm{cos}k.x.$$ As before, we are primarily interested in the limit $`L\mathrm{}`$, followed by the limit $`k0`$. The relative number operator $`N_{L,rel}`$ is not the generator of a symmetry of the local Hamiltonian $`H_L`$ (3) because of the Josephson coupling term. This means that there is no question of spontaneous symmetry breaking for the relative number operator, as it is not a symmetry. Furthermore a straightforward computation of the dynamics of $`N_{L,rel}`$ learns that its spectrum belongs to the excitation branch $`E^+E^{}`$. Since $`E_k^+E_k^{}=2\gamma >0`$, the spectrum of $`N_{L,rel}`$ shows an energy gap. Hence the $`0`$-mode fluctuations of the relative number operator and its adjoint will be normal. Therefore the limits $`L\mathrm{}`$ and $`k0`$ may be interchanged and hence the starting point of our investigation can be the operator (i.e. the case $`k=0`$): $$F_L(n_{rel})=\frac{1}{V^{1/2}}_\mathrm{\Lambda }𝑑x\left(a_1^{}(x)a_1(x)a_2^{}(x)a_2(x)\right)=\frac{1}{V^{1/2}}\underset{k}{}a_{1,k}^{}a_{1,k}a_{2,k}^{}a_{2,k}.$$ (21) Next define the fluctuation operator of the relative current: $`F_L(j_{rel})`$ $`={\displaystyle \frac{i}{2\gamma V^{1/2}}}{\displaystyle _\mathrm{\Lambda }}𝑑x\left(a_1^{}(x)a_2(x)e^{i\phi }a_2^{}(x)a_1(x)e^{i\phi }\right)`$ $`={\displaystyle \frac{i}{2\gamma V^{1/2}}}{\displaystyle \underset{k}{}}a_{1,k}^{}a_{2,k}e^{i\phi }a_{2,k}^{}a_{1,k}e^{i\phi }.`$ (22) This operator clearly corresponds to the 0-mode fluctuations of the relative current from one gas into the other. Below, we show in a series of steps that the operators (21) and (5) are each others adjoint in the limit $`L\mathrm{}`$. Afterwards we show the relations between the relative current fluctuation operator on the one hand and the relative phase fluctuation operator on the other hand. A central limit theorem and reconstruction theorem can again be proved for these operators (see ), proving the existence of the Bosonic field operators $`F(n_{rel})`$ $`=\underset{L\mathrm{}}{lim}F_L(n_{rel})`$ $`F(j_{rel})`$ $`=\underset{L\mathrm{}}{lim}F_L(j_{rel})`$ in a rigorous mathematical sense. Let $`\delta _\omega ()=lim_L\mathrm{}[H_L^\omega ,]`$ be the infinitesimal generator of the dynamics ($`H_L^\omega `$ is the effective Hamiltonian defined in equation (7), with $`\mu =\lambda \rho \gamma `$). In order to prove the properties below, it is convenient to write (21) and (5) in terms of the quasi-particle operators (8): $`F_L(n_{rel})`$ $`={\displaystyle \frac{1}{V^{1/2}}}{\displaystyle \underset{k}{}}b_{+,k}^{}b_{,k}+b_{,k}^{}b_{+,k}`$ (23) $`F_L(j_{rel})`$ $`={\displaystyle \frac{i}{2\gamma V^{1/2}}}{\displaystyle \underset{k}{}}b_{+,k}^{}b_{,k}b_{,k}^{}b_{+,k}.`$ (24) ###### Proposition 2. The operators $`F(n_{rel})`$ and $`F(j_{rel})`$ form a canonical pair and satisfy $$[F(n_{rel}),F(j_{rel})]=ic_{rel},$$ where $$c_{rel}=\beta \underset{L\mathrm{}}{lim}(F_L(n_{rel}),F_L(n_{rel}))_{}=\frac{1}{\gamma }\left(\rho _0+\rho (0)\rho (2\gamma )\right)>0;$$ $`\rho (\alpha )`$ is defined in (13), and $`(,)_{}`$ is the Duhamel two-point function defined below. ###### Proof. The first statement follows again from the general theory on normal fluctuation operators . Also, it is an easy calculation to show that $$[H_L^\omega ,F_L(j_{rel})]=iF_L(n_{rel}).$$ Because of the presence of the energy gap $`2\gamma >0`$, this can be written as: $$F_L(j_{rel})=i\delta _\omega ^1\left(F_L(n_{rel})\right).$$ Therefore, using the fact that $`\omega `$ is an equilibrium (KMS) state, one gets $`[F(n_{rel}),F(j_{rel})]`$ $`=\underset{L\mathrm{}}{lim}\omega \left([F_L(n_{rel}),F_L(j_{rel})]\right)`$ $`=i\underset{L\mathrm{}}{lim}\omega \left(F_L(n_{rel})[1e^{\beta \delta _\omega }]\delta _\omega ^1F_L(n_{rel})\right)`$ $`=i\beta \underset{L\mathrm{}}{lim}(F_L(n_{rel}),F_L(n_{rel}))_{}.`$ The explicit expression for the Duhamel two-point function $`(F_L(n_{rel}),F_L(n_{rel}))_{}`$ then follows from $$\omega \left([F_L(n_{rel}),F_L(j_{rel})]\right)=\frac{i}{\gamma V}\underset{k}{}\omega \left(b_{,k}^{}b_{,k}b_{+,k}^{}b_{+,k}\right)\stackrel{L\mathrm{}}{}\frac{i}{\gamma }\left(\rho _0+\rho (0)\rho (2\gamma )\right).$$ The infinitesimal generator $`\delta _\omega `$ of the microdynamics induces a natural infinitesimal generator $`\stackrel{~}{\delta }_\omega `$ of a dynamics on the macroscopic fluctuation operators by the formula : $$\stackrel{~}{\delta }_\omega F(A)=F(\delta _\omega (A)).$$ ###### Proposition 3. The infinitesimal generator $`\stackrel{~}{\delta }_\omega `$ on the macroscopic fluctuations is given by: $`\stackrel{~}{\delta }_\omega F(n_{rel})`$ $`=i(2\gamma )^2F(j_{rel})`$ (25) $`\stackrel{~}{\delta }_\omega F(j_{rel})`$ $`=iF(n_{rel}).`$ (26) Hence $`F(n_{rel})`$ and $`F(j_{rel})`$ are eigenvectors of $`\stackrel{~}{\delta }_\omega ^2`$: $$\stackrel{~}{\delta }_\omega ^2F(n_{rel})=(2\gamma )^2F(n_{rel}),\stackrel{~}{\delta }_\omega ^2F(j_{rel})=(2\gamma )^2F(j_{rel}),$$ yielding the macrodynamics $`\stackrel{~}{\alpha }_t`$ on the fluctuation operators: $`\stackrel{~}{\alpha }_tF(n_{rel})`$ $`=e^{it\stackrel{~}{\delta }_\omega }F(n_{rel})=F(n_{rel})\mathrm{cos}(2\gamma t)+(2\gamma )F(j_{rel})\mathrm{sin}(2\gamma t)`$ $`\stackrel{~}{\alpha }_tF(j_{rel})`$ $`=e^{it\stackrel{~}{\delta }_\omega }F(j_{rel})={\displaystyle \frac{1}{2\gamma }}F(n_{rel})\mathrm{sin}(2\gamma t)+F(j_{rel})\mathrm{cos}(2\gamma t).`$ ###### Proof. This follows immediately from the relations $`[H_L^\omega ,F_L(j_{rel})]`$ $`=iF_L(n_{rel})`$ $`[H_L^\omega ,F_L(n_{rel})]`$ $`=i(2\gamma )^2F_L(j_{rel}).`$ Remark that we proved that the pair of variables $`(F(n_{rel}),F(j_{rel}))`$ is dynamically independent from the other variables of the system. The pair behaves dynamically as a pair of quantum oscillator variables with a frequency equal to $`2\gamma `$. ###### Proposition 4 (Virial theorem). The mean square fluctuation of the relative number operator is proportional to the mean square fluctuation of the relative current operator, in particular: $$\stackrel{~}{\omega }\left(F(n_{rel})^2\right)=(2\gamma )^2\stackrel{~}{\omega }\left(F(j_{rel})^2\right).$$ ###### Proof. This follows from the time invariance of $`\stackrel{~}{\omega }`$, i.e. $`\stackrel{~}{\omega }\stackrel{~}{\delta }_\omega =0`$: $$0=\stackrel{~}{\omega }\left(\stackrel{~}{\delta }_\omega [F(n_{rel})F(j_{rel})]\right)=\stackrel{~}{\omega }\left(\stackrel{~}{\delta }_\omega [F(n_{rel})]F(j_{rel})\right)+\stackrel{~}{\omega }\left(F(n_{rel})\stackrel{~}{\delta }_\omega [F(j_{rel})]\right),$$ and the equations of motion (25) and (26). ∎ ###### Proposition 5. The mean square fluctuation of the relative number operator is given by $$\stackrel{~}{\omega }\left(F(n_{rel})^2\right)=c_{rel}\gamma \mathrm{coth}\beta \gamma .$$ ###### Proof. We compute this quantity using the correlation inequalities (10), rewritten in the form $$\frac{\beta \omega \left(X\delta _\omega (X^{})\right)}{\omega (XX^{})}\mathrm{ln}\frac{\omega (X^{}X)}{\omega (XX^{})}\frac{\beta \omega \left(X^{}\delta _\omega (X)\right)}{\omega (X^{}X)}.$$ (27) We take for $`X`$ the operator $`A_L=F_L(n_{rel})+i(2\gamma )F_L(j_{rel})`$ and then let $`L\mathrm{}`$, and use proposition 2. One gets $$\underset{L\mathrm{}}{lim}\omega (A_LA_L^{})=\stackrel{~}{\omega }\left(F(n_{rel})^2\right)+(2\gamma )^2\stackrel{~}{\omega }\left(F(j_{rel})^2\right)+(2\gamma )c_{rel},$$ and by the virial theorem (proposition 4): $$\underset{L\mathrm{}}{lim}\omega (A_LA_L^{})=2\stackrel{~}{\omega }\left(F(n_{rel})^2\right)+(2\gamma )c_{rel}.$$ Analogously $$\underset{L\mathrm{}}{lim}\omega (A_L^{}A_L)=2\stackrel{~}{\omega }\left(F(n_{rel})^2\right)(2\gamma )c_{rel}.$$ On the other hand, $$\delta _\omega (A_L)=i(2\gamma )^2F_L(j_{rel})(2\gamma )F_L(n_{rel})=(2\gamma )A_L,$$ and hence $`\underset{L\mathrm{}}{lim}\omega \left(A_L^{}\delta _\omega (A_L)\right)`$ $`=4\gamma \stackrel{~}{\omega }\left(F(n_{rel})^2\right)+(2\gamma )^2c_{rel}`$ $`\underset{L\mathrm{}}{lim}\omega \left(A_L\delta _\omega (A_L^{})\right)`$ $`=4\gamma \stackrel{~}{\omega }\left(F(n_{rel})^2\right)+(2\gamma )^2c_{rel}.`$ After substitution in (27) one gets $$\mathrm{ln}\frac{2\stackrel{~}{\omega }\left(F(n_{rel})^2\right)(2\gamma )c_{rel}}{2\stackrel{~}{\omega }\left(F(n_{rel})^2\right)+(2\gamma )c_{rel}}=2\beta \gamma ,$$ or alternatively $$\stackrel{~}{\omega }\left(F(n_{rel})^2\right)=c_{rel}\gamma \mathrm{coth}\beta \gamma .$$ This finishes the complete study of the static and dynamic properties of the canonical pair $`(F(n_{rel}),F(j_{rel}))`$ of the relative density and current fluctuations. We proved rigorously that for all temperatures below the condensation temperature and with non-zero condensate, this pair behaves like a pair of quantum harmonic oscillator variables, describing oscillations of the fluid from type 1 into type 2 and vice versa, yielding the typical interference pattern. The plasmon frequency is given by $`2\gamma `$. All this is physically clear. Our next and final problem is to find out what the position of the phase is in all this. We turn our attention now to look for a relation between the relative current fluctuation operator $`F(j_{rel})`$ and the relative phase fluctuation operator, which we define in an analogous form as the total phase fluctuation operator, as follows: $$F_L(\varphi _{rel})=\frac{i}{2V^{1/2}}_\mathrm{\Lambda }𝑑x\left(b_+^{}(x)b_+(x)\right)=\frac{i}{2}\left(b_{+,0}^{}b_{+,0}\right).$$ Denote its central limit by $`F(\varphi _{rel})`$. First, observe that one can distinguish two terms in the relative current fluctuation (24), namely the $`k=0`$ part and the rest: $$F_L(j_{rel})=\frac{i}{2\gamma V^{1/2}}\left(b_{+,0}^{}b_{,0}b_{,0}^{}b_{+,0}\right)+\frac{i}{2\gamma V^{1/2}}\underset{k0}{}b_{+,k}^{}b_{,k}b_{,k}^{}b_{+,k}.$$ Denote the first term by $$F_L(j_{rel}^0)=\frac{i}{2\gamma V^{1/2}}\left(b_{+,0}^{}b_{,0}b_{,0}^{}b_{+,0}\right),$$ and its central limit by $`F(j_{rel}^0)`$. Also denote by $`\omega ^g`$ the ground state, obtained as the zero-temperature limit of the equilibrium state $`\omega `$: $$\omega ^g()=\underset{\beta \mathrm{}}{lim}\omega (),$$ and $`\stackrel{~}{\omega }^g`$ the corresponding ground state for the limiting fluctuation operators observables: $$\stackrel{~}{\omega }^g()=\underset{\beta \mathrm{}}{lim}\stackrel{~}{\omega }().$$ ###### Proposition 6. We have the following relationships between the limiting fluctuation operators: 1. $`\beta >0`$, $`\beta =\mathrm{}`$ included, $$F(j_{rel}^0)=\frac{\sqrt{\rho _0}}{\gamma }F(\varphi _{rel});$$ 2. for $`\beta =\mathrm{}`$, $$F(j_{rel})=F(j_{rel}^0)=\frac{\sqrt{\rho _0}}{\gamma }F(\varphi _{rel}).$$ ###### Proof. As shown in , two fluctuation operators $`F(A),F(B)`$ are equal in the algebra of fluctuation operators whenever $$\stackrel{~}{\omega }\left(F(AB)^2\right)=0,$$ (28) i.e. whenever the variance of the difference $`AB`$ of the operators vanishes. This is expressing in a mathematical rigorous setting, the phenomenon of coarse graining on the level of fluctuations. Therefore we calculate $$\begin{array}{cc}\hfill \stackrel{~}{\omega }\left(\left[F(j_{rel}^0)\frac{\sqrt{\rho _0}}{\gamma }F(\varphi _{rel})\right]^2\right)& =\stackrel{~}{\omega }\left(F(j_{rel}^0)^2\right)+\frac{\rho _0}{\gamma ^2}\stackrel{~}{\omega }\left(F(\varphi _{rel})^2\right)\hfill \\ & \frac{\sqrt{\rho _0}}{\gamma }\stackrel{~}{\omega }\left(F(j_{rel}^0)F(\varphi _{rel})\right)\frac{\sqrt{\rho _0}}{\gamma }\stackrel{~}{\omega }\left(F(\varphi _{rel})F(j_{rel}^0)\right).\hfill \end{array}$$ One finds, using the explicit knowledge of the state $`\omega `$ (section 3): $`\stackrel{~}{\omega }\left(F(j_{rel}^0)^2\right)`$ $`=\stackrel{~}{\omega }\left(F(\varphi _{rel})^2\right)={\displaystyle \frac{\rho _0}{4\gamma ^2}}\mathrm{coth}\beta \gamma `$ $`\stackrel{~}{\omega }\left(F(j_{rel}^0)F(\varphi _{rel})\right)`$ $`=\stackrel{~}{\omega }\left(F(\varphi _{rel})F(j_{rel}^0)\right)={\displaystyle \frac{\sqrt{\rho _0}}{4\gamma }}\mathrm{coth}\beta \gamma ,`$ hence leading to the following equality, as operators: $$F(j_{rel}^0)=\frac{\sqrt{\rho _0}}{\gamma }F(\varphi _{rel}).$$ From proposition 4 and 5, it follows that $$\stackrel{~}{\omega }\left(F(j_{rel})^2\right)=\frac{c_{rel}}{4\gamma }\mathrm{coth}\beta \gamma ,$$ and from proposition 2, $$\underset{\beta \mathrm{}}{lim}c_{rel}=\frac{\rho _0}{\gamma }.$$ Therefore $$\stackrel{~}{\omega }^g\left(F(j_{rel})^2\right)=\frac{\rho _0}{4\gamma ^2}=\stackrel{~}{\omega }^g\left(F(j_{rel}^0)^2\right).$$ This implies necessarily $$\stackrel{~}{\omega }^g\left(\left[F(j_{rel})F(j_{rel}^0)\right]^2\right)=0,$$ and again the equality of the operators $$F(j_{rel})=F(j_{rel}^0)$$ in the ground state, as a result of coarse graining. ∎ The physical interpretation of this proposition is the following. For non-zero temperatures, the relative current consists of two terms. One of them is $`j_{rel}^0`$, which has a non-trivial contribution to the fluctuation of the relative current only if $`\rho _0>0`$, i.e. whenever the gauge symmetry is spontaneously broken. The other term contains no more reference to the zero mode, in other words to the condensate. Therefore it is clear that the fluctuation operator $`F(j_{rel}^0)`$ contains all the information of the fluctuations of what one could call the condensate current, or the current between the condensates interacting through the Josephson junction. The important equality $$F(j_{rel}^0)=\frac{\sqrt{\rho _0}}{\gamma }F(\varphi _{rel})$$ is nothing but a rigorous translation, on the level of the fluctuations, of the popular statement: “the (superfluid, condensate) current is the gradient of the phase”. The second statement of the proposition shows that the quantum effects on the level of the fluctuations, originating from the spontaneous symmetry breaking, are only present in the ground state. This of course is popular wisdom, already experienced in many models , but expressed here in a mathematically rigorous fashion for our model. Finally, it may come as a surprise that our results show no dependence on $`\phi `$, the expectation value for the phase difference between the condensates (see equation (15)). This however is a simple consequence of the description of the system in its mathematically simplest form, using the operators $`b_{\pm ,k}^{\mathrm{}}`$ (8), which yield a $`\phi `$-independent description of the system. Indeed, from a mathematical point of view, the system can not be expected to behave different for different $`\phi `$, since e.g. the eigenvalues of the Hamiltonian $`E_k^\pm `$ are $`\phi `$-independent. If one is interested in the physics following from a non-zero $`\phi `$, i.e. if one wants to derive typical Josephson currents proportional to $`\mathrm{sin}\phi `$, one needs to work with the bare operators $`a_{(1,2),k}^{\mathrm{}}`$. In particular, consider the relative current fluctuation operator defined by $$F_L(j_{rel}^\phi )=\frac{1}{2\gamma V^{1/2}}\underset{k}{}a_{1,k}^{}a_{2,k}+a_{2,k}^{}a_{1,k}\omega (a_{1,k}^{}a_{2,k}+a_{2,k}^{}a_{1,k}),$$ in stead of (5). It can easily be calculated that this operator satisfies $`\underset{L\mathrm{}}{lim}[F_L(j_{rel}),F_L(j_{rel}^\phi )]`$ $`=0`$ $`\underset{L\mathrm{}}{lim}[F_L(n_{rel}),F_L(j_{rel}^\phi )]`$ $`=ic_{rel}\mathrm{sin}\phi `$ $`\delta ^\omega \left(F_L(j_{rel}^\phi )\right)`$ $`=i\mathrm{sin}\phi F_L(n_{rel}).`$ Together with the results above, this establishes that the limiting fluctuation operator $`F(j_{rel}^\phi )`$ is given by $$F(j_{rel}^\phi )=F(j_{rel})\mathrm{sin}\phi ,$$ (29) where (29) is to be understood in terms of the equivalence between fluctuation operators (28), i.e. (29) follows from $$\underset{L\mathrm{}}{lim}\omega \left(\left(F_L(j_{rel}^\phi )\mathrm{sin}\phi F_L(j_{rel})\right)^2\right)=0.$$ The variances of $`F(j_{rel}^\varphi )`$ and its dynamics, computed from proposition 3 show the explicit $`\phi `$-dependence, and its typical ‘collapse and revival’ properties, found in the experiments .
warning/0004/astro-ph0004055.html	ar5iv	text	# On the Two-Phase Structure of Protogalactic Clouds1 ## 1 Introduction According to CDM models of galaxy formation (Blumenthal et al. 1984; Navarro, Frenk, & White 1997; Klypin, Nolthenius, & Primack 1997), large galaxies emerge through the coalescence of smaller systems, normally identified as dwarf galaxies, which contain mostly dark matter with a fraction of gas. The dark matter components merge with each other in a dissipationless manner to form the halos of larger galaxies. The crossing of gas streamlines leads to shocks, heating the gas to the virial temperature of the galactic halo (Binney 1977; Rees & Ostriker 1977; White & Rees 1978). In order for the heated gas in this extended halo (hereafter, the protogalactic cloud or PGC) to collapse, its cooling timescale $`\tau _c`$ must be shorter than its dynamical time scale, $`\tau _d`$ during each growth stage of the galaxy. When large galaxies acquire a mass comparable to the Galaxy, this condition is satisfied once the characteristic length scale of the PGC is $`<100`$ kpc (Blumenthal et al. 1984). For typical values ($`10^6`$ K) of the virial temperature, the cooling timescale increases with temperature, and the PGC’s are thermally unstable (Field 1965). Thermal instability leads to the rapid growth of perturbations and fragmentation of PGCs (Murray & Lin 1990). The result is that a two-phase medium develops during the initial cooling of the PGC, in which a population of warm fragmentary clouds (WFC’s) are confined by the pressure of hot, residual halo gas (RHG) (Burkert & Lin 2000). The RHG is primarily heated by the release of the gravitational energy of the collapsing PGC and cooled by radiative emission and conductive transport. The WFC’s settle into the central region of the halo potential. They are unable to cool below 10<sup>4</sup> K until their density reaches a sufficiently high value that the WFC’s become self-shielded from external photodissociating UV radiation (Couchman & Rees 1986; Haiman, Rees, & Loeb 1997; Haiman, Abel, & Rees in press). Thereafter, molecular hydrogen can form within them as a consequence of non-equilibrium cooling. Efficient cooling reduces the temperature to $`T<100`$ K even in the absence of metals. Consequently, WFC’s evolve into cold molecular clouds (CMCs). The formation of massive stars within the CMC’s provides ultraviolet heating, limiting the formation rate of CMC’s from WFC’s, and leading to the formation of an equilibrium three-phase structure in the galaxy. In this paper, we discuss the evolution of the hot phase that is left after the initial thermal instability. Detailed investigations of the warm and cold phases will be made in upcoming papers. For the current work, we assume a power-law size distribution for the WFC’s. The CMC’s are not included explicitly, but star formation within them is explicitly assumed to occur at a rate which would provide an adequate flux of UV photons to heat and ionize the WFC’s. In §2, we briefly recapitulate the physical processes associated with this scenario. In § 3, we determine the energy balance between the warm and hot phases, while in § 4, we discuss the resulting star formation rate in the WFC’s. In § 5, we show model results on the evolution of the warm and hot phases. Finally, we summarize our results and discuss their implications in § 6. ## 2 The Emergence of a Two-phase Medium The temperature of a large, virialized PGC, $`TV_c^2\mu /k10^6(V_c/100\mathrm{km}\mathrm{sec}^1)^2`$ K, where $`V_c`$ is the circular velocity of the halo potential, $`\mu `$ is the mean atomic weight in grams, and $`k`$ is the Boltzmann constant. At such temperatures, the dominant radiative cooling mechanisms are bremsstrahlung and recombination processes (Dalgarno & McCray 1972). The radiative cooling rates for metal poor clouds (with \[Fe/H\] < 2 < absent2\mathrel{\vbox{ \offinterlineskip\hbox{$<$} \kern 1.29167pt\hbox{$\sim$} }}-2) are $$\mathrm{\Lambda }(T)10^{22}\left(\frac{T}{10^5\mathrm{K}}\right)^{1.2}\mathrm{ergs}\mathrm{cm}^3\mathrm{s}^1$$ (1) for $`2\times 10^4`$ K $`<T<10^6`$ K, and $$\mathrm{\Lambda }(T)6\times 10^{24}\left(\frac{T}{10^6\mathrm{K}}\right)^{1/2}\mathrm{ergs}\mathrm{cm}^3\mathrm{s}^1$$ (2) for T > 106 > 𝑇superscript106T\mathrel{\vbox{ \offinterlineskip\hbox{$>$} \kern 1.29167pt\hbox{$\sim$} }}10^{6} K. Equation (1) is a highly simplified average representation of the behavior of the cooling rate in the given temperature range. In actuality, there is a peak due to H emission at $`2\times 10^4`$ K, another due to emission by He<sup>+</sup> at 10<sup>5</sup> K, and a minimum at $`10^6`$ K, before the cooling efficiency rises again, due to bremsstrahlung emission (Gould & Thakur 1970). The cooling timescale is $$\tau _c=\frac{\frac{3}{2}\rho kT}{\mu n^2\mathrm{\Lambda }}$$ (3) where $`\rho `$ and $`n`$ are the mass and number density of the gas. The necessary collapse condition, $`\tau _c\tau _d`$ where the dynamical timescale $`\tau _d=R/V_c`$ (White & Rees 1978), for a virialized PGC is satisfied when $`n`$ exceeds a critical value $$n_{crit}0.2R_{kpc}^1\left(\frac{T}{10^6\mathrm{K}}\right)^{\frac{1}{2}}\mathrm{cm}^3,$$ (4) where $`R_{kpc}=R/1`$ kpc, and R is the galactocentric radius. Above $`3\times 10^4`$ K, $`\tau _c`$ is an increasing function of $`T`$, such that the contrast between slightly cooler regions and the background grows during the cooling process. The relatively rapid loss of entropy also leads to a deficit of pressure, $`P`$, in the clouds compared to their background. The higher pressure of the background compresses the clouds in an attempt to maintain pressure balance. The rise in $`n`$ leads to a reduction in $`\tau _c`$ for the clouds, further enhancing the dichotomy in cooling timescale between the clouds and the background. The interface separating cool clouds from the background retreats at an accelerating pace. Although $`\tau _c\tau _d`$ at the onset of thermal runaway, it decreases rapidly to $`\tau _d`$ as the perturbed region cools significantly below the virial temperature. Clouds, with radius $`S`$ and sound speed $`c_s`$, establish hydrostatic equilibrium with the RHG and undergo isobaric evolution on a timescale $`S/c_s`$. If this timescale is $`>\tau _c`$, $`n`$ cannot adjust such that cooling is isochoric. Relatively large clouds undergo a transition from isobaric to isochoric cooling before their density becomes significantly larger than that of the background. This transition in small clouds occurs at a later cooling stage. Their density growth becomes nonlinear before the cooling becomes isochoric. Consequently, they emerge as dense, cool fragments (Burkert & Lin 2000). When $`T`$ in the clouds decreases to $`10^4`$K, recombination becomes the dominant cooling process. However, these clouds are also exposed to UV radiation due to emission from background active galactic nuclei. At redshift $`z=2.5`$, the flux of this emission at the Lyman limit is estimated to be $`F_{UV}=0.5\times 10^{21}`$ ergs Hz<sup>-1</sup> ster <sup>-1</sup> cm<sup>-2</sup> s<sup>-1</sup>, with a dropoff towards higher energies as steep as $`\nu ^3`$ (Haardt & Madau 1996). The above value varies approximately as $`(1+z)^3`$. In the outer regions of PGC where the pressure of the RHG is low, the internal density of these clouds is also relatively small, such that they may be mostly photoionized, with $`T10^4`$K. But, in the inner regions of the PGC where the background pressure is higher, WFC’s have larger internal density, such that they become self shielded against the extragalactic UV flux and their temperature cools below $`100`$ K due to emission by H<sub>2</sub> and metal ions (Dalgarno & McCray 1972; Hollenbach & McKee 1979). When these clouds become gravitationally unstable, stars rapidly form within them. The UV flux provided by these massive stars is able to heat gas within their Strömgren radii to $`10^4`$ K, quenching star formation in those regions. At all galactocentric radii, then, the PGC fragments into a population of WFC’s which are maintained at $`10^4`$ K by the UV flux. In most regions, the ionizing photons are emitted by a population of massive stars whose formation rate is self-regulated. The WFC’s are embedded within RHG, whose density is sufficiently low that $`\tau _c>\tau _d`$, ie. it is thermally stable (Field 1965), and it can remain near the virial temperature of the galaxy. The pressure of the RHG confines all but the largest WFC’s, for which self gravity is important. Consequently, the PGC becomes a two-phase medium, with the density contrast between the WFC’s and the RHG inversely proportional to their temperature ratio ($`100`$). ## 3 Energy Budget of the Two-Phase Medium ### 3.1 Dynamical Interaction Between the Warm and Hot Phases Following its collapse into the potential of the galactic halo, the RHG is shock-heated to the virial temperature of the potential, and rapidly attains a quasistatic equilibrium. It then adjusts to have density and temperature appropriate for the maintenance of thermal equilibrium constrained by the processes discussed below. The inverse buoyancy of the pressure-confined WFC’s causes them to settle towards the center of the potential provided by the dark matter halo. During their descent, the WFC’s experience a drag by the RHG, reaching a terminal speed relative to the RHG given by $$V_t\left(D_\rho \frac{S}{R}\right)^{1/2}V_c,$$ (5) where $`D_\rho \rho _w/\rho _h`$ is the ratio of the density of the WFC’s to that of the RHGs. (In all relations below, the subscripts “h” and “w” refer to quantities of the RHG and of the WFC, respectively.) For WFC’s with sizes $`S>R/D_\rho `$, $`V_tV_c`$. Here, we neglect the effect of the clouds’ size evolution on their kinematics. The value of $`V_t`$ derived in equation (5) assumes a constant density background. In reality, the pressure gradient of the background causes the density of the RHG to rise towards the center of the galaxy. The radii of the WFC’s $`SR`$, and so the pressure gradient across an individual cloud is negligible. If the RHG is approximately isothermal, and the temperature of the WFC’s are determined by atomic cooling to be $`10^4`$ K, then $`D_\rho `$ remains constant during the motion of the clouds. As a result, the primary effect of the density gradient of the RHG upon $`V_t`$ is in altering the ratio $`S/R`$. If $`\rho _hR^a`$, then, under the above assumptions, $`\rho _wR^a`$, $`SR^{a/3}`$, and $`V_tR^{(a6)/6}`$. The motion of the WFC’s through the RHG also leads to mass loss from the WFC’s due to Kelvin-Helmholtz instability (Murray et al. 1993). In the limit $`D_\rho 1`$, the growth timescale of KH instability is given by $$\tau _{KH}=\lambda \frac{(\rho _h+\rho _w)}{(\rho _h\rho _w)^{1/2}V_t}\frac{\lambda D_\rho ^{1/2}}{V_t}.$$ (6) Short wavelength perturbations rapidly saturate. Perturbations with wavelengths $`\lambda S`$ provide the dominant contribution to mass loss, which occurs at a rate $$\dot{M}\frac{4\pi \rho _wS^2\lambda }{\tau _{KH}}\frac{4\pi \rho _wS^2V_t}{D_\rho ^{\frac{1}{2}}}.$$ (7) For a WFC with a mass $`M=4\pi \rho _wS^3/3`$, the stripping timescale $`\tau _s=M/\dot{M}(R/3V_c)(S/R)^{1/2}<\tau _d`$, i.e. WFC’s are disintegrated before they settle to the center of the galactic potential. Thus, both thermal instability associated with the cooling, and Kelvin-Helmholtz instability arising during the infall of the WFC’s leads to the formation of WFC’s with arbitrarily small sizes. The break down process increases the collective area filling factor of the WFC’s, leading to an increase in their collision frequency. Collisions and the resulting mergers increase the masses of the WFC’s. When the rates of coagulation and disruption (including other processes discussed below) balance with each other, an equilibrium size distribution is established (Dong et al. in preparation). ### 3.2 Conductive Heat Transfer Between Warm Fragments and Hot Medium The existence of WFC’s in the RHG also leads to heat transfer from the RHG to the WFC’s through conduction (McKee & Cowie 1977). This process is particularly important for small WFC’s. The conductive heat flux into a WFC is given by $$F_c=\kappa T\mathrm{ergs}\mathrm{cm}^2\mathrm{s}^1,$$ (8) where $`\kappa 6\times 10^7T^{5/2}`$ (Spitzer 1956). To order of magnitude, $$F_c\kappa \frac{T}{l}\left(\frac{1}{3}\kappa _hT_hn_w^2\mathrm{\Lambda }_w\right)^{\frac{1}{2}},$$ (9) where $`l`$ is estimated as the depth into the WFC’s over which radiative emission balances conduction into the WFC’s, and $`\mathrm{\Lambda }_w`$ is the cooling efficiency within the WFC. In terms of the density contrast, $$F_c\frac{\mu _h}{\mu _w}D_\rho \left(\frac{1}{3}\kappa _hT_h\mathrm{\Lambda }_w\right)^{\frac{1}{2}}n_h,$$ (10) where $`\mu _h`$ and $`\mu _w`$ are the mean molecular weights in the RHG and WFC’s, respectively, and $`n_h=\rho _h/\mu _h`$. In small clouds where the density scale height is small, the magnitude of $`F_c`$ is limited by the saturated flux, $$F_cF_{sat}=\rho _hc_h^3,$$ (11) where $`c_h`$ is the isothermal sound speed in the RHG ($``$ 220 km s <sup>-1</sup> for the Galaxy). The presence of magnetic field reduces the mean free path of the electrons and $`F_c`$ (Rosner & Tucker 1989). In our calculations, we neglect the effect of magnetic field and adopt $`F_c=F_{sat}`$ based on our model parameters. Large WFC’s, which attain large $`V_t`$, may also advect thermal energy from the RHG to the WFCs at a rate $$F_a\frac{1}{2}\rho _hV^3.$$ (12) As discussed above, the sizes of the WFC’s are modified by Kelvin-Helmholtz instability and coagulation before they are accelerated to $`V_t`$. The WFC’s may therefore spend most of their time with smaller velocities relative to the RHG (Murray et al. 1993). The advective heating is therefore expected to be weaker than that due to conduction. ### 3.3 The Warm Fragmentary Clouds Energy equilibrium is established within a cloud when $$\frac{4}{3}\pi S^3\mathrm{\Gamma }_r+4\pi S^2F_c=\frac{4}{3}\pi S^3n_w^2\mathrm{\Lambda }_w.$$ The photoionization heating $$\mathrm{\Gamma }_r=(1x_w)n_wh_{\nu _0}^{\mathrm{}}\frac{4\pi J_\nu }{h\nu }a_\nu (\nu \nu _0)𝑑\nu =(1x_w)n_wf_G4\pi J_{\nu _0}a_{\nu _0}\nu _0,$$ (13) where $`x_w`$ is the ionization fraction in the WFC’s, $`h`$ is Planck’s constant, $`J_\nu `$ is the mean intensity in ergs ster<sup>-1</sup> Hz<sup>-1</sup> s<sup>-1</sup>, $`a_\nu `$ is the photoionization cross section as a function of frequency, $`h\nu _0=13.6`$ eV, and $`f_G`$ depends upon the shape of the spectrum ($`f_G<0.1`$ for realistic background spectra). This equilibrium determines the temperature of the WFC’s. The above assumes a top-hat density distribution within the clouds. This is appropriate for clouds which are at most weakly self-gravitating, given the sharp temperature transition between the WFC’s and the RHG resulting from thermal instability. In a mostly ionized WFC, the heat input due to conduction is larger than that due to photoionization when $$x_w>x_c=1\frac{3}{S}\frac{F_c}{n_wf_G4\pi J_{\nu _0}a_{\nu _0}\nu _0}.$$ (14) In this limit, the rate of conductive heat input exceeds that of the radiative loss due to recombination when $$S<S_{es}=\frac{3\mu _hn_hc_h^3}{\mathrm{\Lambda }_wn_w^2}$$ (15) in the limit of saturated conduction and $$S<S_{eu}=\left(\frac{3\kappa _hT_h}{n_w^2\mathrm{\Lambda }_w}\right)^{1/2}$$ (16) in the limit of unsaturated conduction. The above results also apply to the limit in which the background UV flux is negligibly small, so that the gas in WFC’s is mostly neutral. For the intermediate case in which photoionization contributes significantly to the heating, the additional energy input would tend to increase the value of $`S_{eu}`$ and $`S_{es}`$. WFC’s with $`S<S_{es}`$ (or $`S_{eu}`$) are unstable to conductive heat transport. Their inability to radiate away all the conductive heat input leads to an increase in their temperatures. In order to maintain pressure equilibrium with the RHG, these WFC’s expand, decreasing their internal density. The increase in $`S`$ enhances conductive heat transport into the clouds, whereas their cooling efficiency decreases with the decrease in $`n_w`$. Unless $`\mathrm{\Lambda }_w`$ increases rapidly with $`T_w`$, conduction then leads to the total evaporation of these small WFC’s. The stability criterion can be obtained by analyzing the dependence of $`S_{es}`$, $`S_{eu}`$, and $`S`$ upon $`T_w`$. If we approximate $`\mathrm{\Lambda }_wT_w^s`$, then from equations 15 and 16 we find that both $`S_{es}`$ and $`S_{eu}`$ vary as $`T_w^{(2s)/2}`$. As the pressure-confined WFC’s are heated and expand, however, their radii vary as $`T_w^{1/3}`$. The clouds are therefore unstable if $`s<4/3`$. In that case, we can picture a cloud with $`S=S_{es}`$ initially. As it is heated and expands, the steep dependence of $`S_{es}`$ upon $`T_w`$ leads to $`S<S_{es}`$. Radiative cooling can then no longer offset conduction into the cloud, leading to further temperature increase and expansion. The same stability criterion of $`s>4/3`$ holds for clouds where $`S_{eu}`$ is the appropriate limiting radius. This stability condition is met over the temperature range dominated by hydrogen emission, $`9000\mathrm{K}T_w20,000`$ K, such that clouds with $`S<S_{es}`$ or $`S<S_{eu}`$ remain stable until their temperatures reach 20,000 K. The minimum possible stable cloud size is therefore set by the value of $`S_{es}`$ or $`S_{eu}`$ at 20,000 K, where $`\mathrm{\Lambda }(T)`$ is a maximum. In the above analysis, the WFC’s are assumed to be confined by the ambient pressure of the RHG. If the WFC’s are able to cool to T < 102 < 𝑇superscript102T\mathrel{\vbox{ \offinterlineskip\hbox{$<$} \kern 1.29167pt\hbox{$\sim$} }}10^{2} K, then self-gravity contributes to the confinement of even relatively small clouds. In this limit, thermal equilibrium between photoionization heating and radiative losses may still be established, but it is unstable because of the negative specific heat of self-gravitating clouds (Murray & Lin 1992). The result is similar to that for non self-gravitating clouds, except that, for a self-gravitating cloud, the radius varies as $`T_w^1`$, and stability requires $`s>4`$, which is, again, only satisfied over a very narrow range of temperatures. While the lower limit of cloud radii is set by their ability to radiate away the energy flux due to conduction in equations (15) and (16), the upper limit of cloud radii is set by gravitational instability. Low mass clouds are confined by the pressure of the RHG. At higher masses, self-gravity becomes important. Clouds become unstable above a critical mass, which depends upon the external pressure $$M_{crit}=\left[3.15\left(\frac{kT_w}{\mu _w}\right)^4\frac{1}{G^3P_h}\right]^{\frac{1}{2}},$$ (17) where $`P_h`$ is the pressure of the RHG, and $`T_w`$ is the internal temperature of the clouds (Ebert 1955; Bonner 1956). This condition corresponds to an approximate maximum cloud radius of $$S_{BE}=1.2\left(\frac{k}{G}\right)^{\frac{1}{2}}\left(\frac{3}{4\pi }\right)^{\frac{1}{3}}\left(\frac{T_w}{\mu _w}\right)\left(n_hT_h\right)^{\frac{1}{2}}.$$ (18) We have adopted the Bonner-Ebert criterion for stability, rather than the Jeans criterion, due to the fact that the WFC’s are pressure-confined. While there might be an apparent inconsistency, in that we assume the WFC’s to have nearly constant density, whereas the Bonner-Ebert criterion assumes a density gradient. Had we used the Jeans criterion, however, our value for $`S_{BE}`$ would be very close to that used above, and so the choice of stability criterion does not affect our results below. The above results assume that the WFC’s are subject only to heating by photoionization and by conduction from the hot gas. In addition, the drag of the clouds through the RHG, and the growth of Kelvin-Helmholtz instability will both lead to additional mechanical energy input into the WFC’s which could, in principle, provide additional support for the clouds. These mechanisms are least important, however, in the largest clouds, which experience the least drag, and whose self-gravity stabilizes the growth of deeply-penetrating, long wavelength perturbations by Kelvin-Helmholtz instability (Murray et al. 1993). Stability of the most massive clouds should therefore still be governed by the Bonner-Ebert criterion. Detailed models of clouds moving through background gas will be used in future work to confirm these conclusions. ### 3.4 Total Rate of Conductive Energy Transport The physical processes involved in determining the mass spectrum of the WFC’s will be examined in a future paper. In the results below, we assume that the fragmentation and coagulation of the WFC’s leads to a power-law size distribution of the WFC’s, in which the number density of clouds within the PGC, $`n_c`$, per unit cloud radius $`S`$ is given by $$\frac{dn_c}{dS}=AS^\gamma .$$ (19) The normalization constant, $`A`$, is related to the spatially averaged density of the warm phase by $$\rho _w=\frac{4}{3}\pi S^3\rho _w\frac{\mathrm{d}n_c}{\mathrm{d}S}dS=\frac{4}{3}\pi \rho _wA\left(\frac{1}{\gamma +4}\right)\left(S_{max}^{\gamma +4}S_{min}^{\gamma +4}\right)$$ (20) where $`S_{max}`$ and $`S_{min}`$ are the maximum and minimum cloud size, respectively. As discussed above, $`S_{min}`$ is determined by the smallest size of clouds which can radiate away the energy conducted into them from the RHG, while $`S_{max}`$ is determined by the Bonner-Ebert mass. The value of $`A`$ is a function of galactocentric radius, whereas we assume $`\gamma `$ to be the same everywhere within the PGC. If the clouds follow a distribution close to that observed within Galactic cloud complexes, for which $`\gamma =2.5`$ (Scalo 1985; Zinnecker, McCaughrean, & Wilking 1993 ), then the results above and below are dominated by the massive end of the distribution. Alternatively, if the efficiency of cloud disruption is high relative to the rate of coagulation, then the cloud distribution will be dominated by the low-mass end. The determination of which limit is appropriate shall be made in an upcoming paper. In the current work, we treat the size distribution as an uncertainty, and consider two extreme limits. In the first, we assume $`\gamma >3`$, and all terms derived both above and below are dominated by the high mass end of the cloud distribution. In the second, we assume $`\gamma <4.5`$, and all terms are dominated by the low mass end of the distribution. The total rate of conductive energy transport per unit volume from the RHG into the WFC’s is given by $$L_{cond}=\frac{dn_c}{dS}F_c4\pi S^2𝑑S.$$ (21) Using the results above, this becomes $$L_{cond}=\frac{4\pi AF_c}{\gamma +3}\left(S_{max}^{\gamma +3}S_{min}^{\gamma +3}\right)$$ (22) For $`\gamma >3`$, the above result becomes $$L_{cond}=3\left(\frac{\gamma +4}{\gamma +3}\right)\frac{\rho _w}{\rho _w}F_cS_{BE}^1,$$ (23) while for $`\gamma <4`$ it is $$L_{cond}=3\left(\frac{\gamma +4}{\gamma +3}\right)\frac{\rho _w}{\rho _w}F_cS_{es}^1.$$ (24) The above results assume that $`S_{max}S_{min},`$ where $`S_{max}=S_{BE}`$ and $`S_{min}=S_{es}`$. As can be seen by a comparison of equations (15) and (18), $`S_{min}/S_{max}0.01(\mathrm{\Lambda }_w/10^{24})^1n_h^{1/2},`$ and so this is a reasonable approximation. Also, because $`S_{BE}S_{es}`$, the greater surface-to-volume ratio of the case $`\gamma <4`$ leads to a significantly greater cooling rate due to conduction into the WFC’s relative to the case with $`\gamma >3`$. ### 3.5 The Residual Halo Gas Equations (23) and (24) give the conductive energy transfer rate from the RHG to the WFC’s. The RHG also experiences radiative losses, mainly through bremsstrahlung emission, which occurs at a rate given by equation (2). Using the conductive flux given in equations (23) and (24), we find that conductive losses into the WFC’s exceed bremsstrahlung cooling in the limit that $$n_h^2<3\left(\frac{\gamma +4}{\gamma +3}\right)\left(\frac{S_{max}^{\gamma +3}S_{min}^{\gamma +3}}{S_{max}^{\gamma +4}S_{min}^{\gamma +4}}\right)\frac{\mu _h}{mu_w}\frac{T_w}{T_h}\frac{c_h^3}{\mathrm{\Lambda }_h}\rho _w.$$ (25) when conduction is saturated, and $$n_h^2<\left(\frac{\gamma +4}{\gamma +3}\right)\left(\frac{S_{max}^{\gamma +3}S_{min}^{\gamma +3}}{S_{max}^{\gamma +4}S_{min}^{\gamma +4}}\right)\frac{\left(3\kappa _hT_h\mathrm{\Lambda }_w\right)^{\frac{1}{2}}}{\mathrm{\Lambda }_h}\rho _w$$ (26) when conduction is unsaturated. In the above relations, we assume pressure balance between the warm and hot phases, ie. $`n_hT_h=n_wT_w.`$ Heat is supplied to the RHG through its own quasi-static contraction. More importantly, heat is also deposited into the RHG as a result of the frictional drag on collapsing WFC’s, by which the WFC’s are able to transfer a fraction of their gravitational potential energy into thermal energy of the RHG. (In the absence of any drag, the total energy of individual WFC’s is approximately conserved as they fall into the galactic halo potential.) The heating per unit volume of a flux of WFC’s moving at their terminal speeds is given by $$\mathrm{\Gamma }=\frac{V_c^2}{R}\frac{\mathrm{d}n_c}{\mathrm{d}S}\frac{4}{3}\pi \rho _wS^3V_tdS.$$ (27) Using $`V_t`$ from eq (6), the above result becomes $$\mathrm{\Gamma }=\left(\frac{\gamma +4}{\gamma +\frac{9}{2}}\right)\left(\frac{S_{max}^{\gamma +\frac{9}{2}}S_{min}^{\gamma +\frac{9}{2}}}{S_{max}^{\gamma +4}S_{min}^{\gamma +4}}\right)\left(\frac{\mu _wT_h}{\mu _hT_w}\right)\frac{V_c^3}{R^{\frac{3}{2}}}\rho _w.$$ (28) For $`\gamma >4`$, the above result is dominated by the contribution of the largest WFC’s, for which $`V_t`$ may become comparable to $`V_c`$ at small $`R`$. Whether or not the clouds are able to attain $`V_t`$ depends upon the ratio of their acceleration timescales to the timescales upon which they are affected by conduction, collisions, or Kelvin-Helmholtz instability. For a cloud to be accelerated to a speed $`V_c`$, it must travel a distance $`R/2`$ in the halo, whereas, Kelvin-Helmholtz instability leads to significant mass loss by the time a cloud has moved through the RHG over a distance of several times its diameter (Murray et al. 1993). Unless the RHG is also contracting with a velocity $`V_HV_c`$, the massive WFC’s may not be able to attain their terminal velocities (Dong et al. in preparation). Instead of $`V_t`$ above, we parametrize the terminal differential velocity as $`fV_c`$, with $`f`$ being a free parameter. Using this prescription, the heating rate is given by $$\mathrm{\Gamma }=\frac{fV_c^3}{R}\rho _w.$$ (29) In the limit $`\gamma <4.5`$, the heating is dominated by drag from low mass clouds, for which $`V_t<V_c`$. The above result then becomes $$\mathrm{\Gamma }=\left(\frac{\gamma +4}{\gamma +\frac{9}{2}}\right)\frac{V_c^3}{R^{\frac{3}{2}}}\left(\frac{\mu _wT_h}{\mu _hT_w}\right)^{\frac{1}{2}}\rho _wS_{es}^{\frac{1}{2}}.$$ (30) If the radiative and conductive cooling of the RHG exceed the thermal energy input, some of the RHG will precipitate to form additional WFC’s. As the gas density of the RHG is depleted, both the conductive heat flux into, and drag on the infalling WFC’s decrease. The density $`n_h`$ therefore adjusts so as to achieve thermal equilibrium in which $`GL_{cond}+L_{brem}`$. ## 4 Ionization of the clouds ### 4.1 The Nature of the External Source In the above analysis, the WFC’s are generally taken to be highly ionized. To maintain this ionization state, the UV flux to which the WFC’s are exposed must be sufficiently large to offset recombination. In Figure 1, we show the depths into which a cloud may be ionized by various sources of UV flux, plotted as a function of $`n`$. These results were obtained with CLOUDY (Ferland 1991). The solid, dashed, and dot-dashed curves represent results obtained with the standard AGN spectral energy distribution, for three values of $`F_{UV}`$. The solid curve was calculated using a value of $`0.5\times 10^{21}`$ ergs Hz<sup>-1</sup> ster<sup>-1</sup> cm<sup>-2</sup> s<sup>-1</sup> at 13.6 eV, representing the UV background at z=2.5 (Haardt & Madau 1996). The dashed and dot-dashed curves were computed using fluxes that were, respectively, smaller and larger by an order of magnitude, approximately representing redshifts $`<0.5`$ and 5 respectively. The dotted curve was computed using the Kurucz spectrum in CLOUDY, with the flux set to that of an O5 star at a distance of 10 pc. From Figure 1, we can infer the integrated depth into the PGC outside of which all of the warm fragments are ionized. In the early evolutionary epochs of the Galaxy before most of the ordinary matter was converted into stars, only the regions outside $`30`$ kpc are photoionized by the UV background (see below). At smaller galactocentric radii, local sources of UV radiation are required. In diffuse and low mass dwarf galaxies, however, the entire galaxy may be photoionized by the diffuse background radiation (Kepner, Babul, & Spergel 1997). In addition to the extragalactic background, the RHG radiates soft X-ray photons following a bremsstrahlung spectrum. Previous work (Fall & Rees 1985), which used analytic approximations to estimate the structure of the RHG, found the resulting x-ray flux to be insufficient to ionize large clouds beyond a few kpc from the Galactic center. In the models presented below, we find the bremsstrahlung luminosities to be less than estimated by Fall & Rees (1985). Ionization by the RHG is therefore not a significant contributor to determining the state of the WFC’s. Hydrodynamic drag on WFC’s also leads to energy input into them, as described in § 3.4. For relatively large WFC’s $`V_t>10`$ km s<sup>-1</sup>. If these fragments are not ionized initially, then shocks will cause internal heating. From above, though, the ratio of heating by drag to that by conduction is found to be $$\frac{\frac{\pi }{2}\rho _1S^2V_t^3}{4\pi S^2F_c}=3\times 10^4\left(\frac{V_t}{10\mathrm{km}\mathrm{s}^1}\right)^3.$$ (31) Thus, even for velocities as high as 100 km s<sup>-1</sup>, the heating rate due to drag is negligible. For survivable clouds with $`S>S_{eu}`$ and $`S_{es}`$, conduction does not provide adequate heat transfer to offset radiative loses. Since the dominant radiative process at $`10^4`$K is recombination, conductive heat flow is also insufficient to ionize the WFC’s. ### 4.2 Self-Regulated Star Formation In the absence of an adequate ionizing photon flux, the cooling efficiency of bremsstrahlung, recombination, and atomic hydrogen emission decrease rapidly below $`10^4`$ K. In a metal-free protogalactic cloud, however, non-equilibrium recombination leads to the formation of a small amount of $`H^{}`$ ions which recombine with neutral $`H`$ to form $`H_2`$. Radiative losses due to rotational and vibrational transitions of $`H_2`$ reduce the gas temperature to $`10^2`$ K (Murray & Lin 1990). If \[Fe/H\]$`>3`$, even lower temperatures ($`10`$ K) are attainable due to cooling by heavy elements (primarily by CII, CO, and Silicate grains) (Hellsten & Lin 2000), leading to the formation of cold molecular clouds (CMC’s). When the temperature of the CMC’s is reduced by over two orders of magnitude from that of the WFC’s, their Bonner-Ebert masses are reduced by over four orders of magnitude (cf. equation 17). Self-gravity is therefore much more important for the evolution of the CMC’s than for that of the WFC’s. Cloudlets close to the Bonner-Ebert mass are centrally condensed. When their mass exceeds the Bonner-Ebert mass, thermal pressure can no longer support the weight of the envelope, and the cloudlets undergo inside-out collapse (Shu 1977). During the collapse the Jean’s mass decreases with the increasing density. Due to the centrally-concentrated cloud structure, however, the Jeans mass is always larger than the mass contained inside any given radius. The collapse is therefore stable, and does not lead to fragmentation in the absence of further unstable cooling, such that the initial value of $`M_{crit}`$ represents the minimum mass for isothermal collapsing clouds (Tsai 1992). Numerical simulations show that even rotating centrally-concentrated clouds have difficulties in breaking up into fragments with a small fraction of $`M_{crit}`$ (Burkert & Bodenheimer 1996). This conjecture is in contrast to the conventional opacity-limited fragmentation scenario (Hoyle 1953; Low & Lynden-Bell 1976) which may be more appropriate for the collapse of unstable homogeneous clouds (Burkert & Bodenheimer 1993). The initial conditions required by the opacity-limited fragmentation scenario may be difficult to accomplish in nature. In a metal-free environment, the temperature of the CFC’s, $`T10^2`$ K and $`M_{crit}1010^2M_{}`$. Stars formed in such a metal-poor environment are massive and short-lived, consistent with their rarity today (Hellsten & Lin 2000). The early chemical enrichment is dominated by the output of Type II supernovae, consistent with the abundance distribution observed among stars with \[Fe/H\]$`<1`$ (Wheeler, Sneden, & Truran 1989). The massive stars are also copious sources of UV radiation. The effect is twofold. Photons in the range 11-13.6 eV are able to penetrate clouds and photodissociate H<sub>2</sub> within them, unless their column densities are so high the clouds are self-shielded. Higher energy photons ionize and heat the gas to 10<sup>4</sup> K. The loss of coolants, and increased heating means that a population of $`10^5`$ O5 stars is adequate to photoionize all the CMC’s, re-heat them to $`T10^4`$ K, and prevent the formation of any additional CMC’s out to a few tens of kpc. Photoionization heating also increases $`M_{crit}`$ to $`10^6M_{}`$, stablizing small cloudlets, and quenching star formation. As the massive stars evolve off the main sequence, the UV flux decreases, cooling again leads to $`T10^2`$ K in sheltered regions, and spontaneous star formation is resumed. Thus, the formation rate of massive stars may be estimated from the assumption that it is self-regulated at a level necessary to sustain marginal ionization of all WFC’s. This self-regulated star formation rate naturally yields values of \[Fe/H\] comparable to that observed in halo stars (Lin & Murray 1992). The rate at which massive stars must be formed is determined by the need to balance ionization with recombination in the WFC’s. The mean free path for Lyman continuum photons is large in the RHG, but not in the WFC’s. For sufficiently high covering factor, we can assume that a large fraction of the UV photons produced by hot stars are absorbed locally. In equilibrium then, the total production rate of UV photons per unit volume must balance the total recombination rate per unit volume. This radiative equilibrium results in high ionization in the limit $$n_{}Q_0=n_w^2\alpha _B\frac{4}{3}\pi S^3\frac{\mathrm{d}n_c}{\mathrm{d}S}dS=\left(\frac{T_h}{T_w}\right)\frac{\alpha _B}{\mu _w}\rho _w,$$ (32) where $`n_{}`$ is the number density of stars, $`Q_0`$ is the average UV output in photons s<sup>-1</sup> per star, and $`\alpha _B`$ is the case B recombination coefficient (Osterbrock 1989). The production of UV photons is dominated by O stars, with lifetimes $`\tau _{}3\times 10^6`$ yr. The average mass loss rate per unit volume from the warm phase due to star formation is then given by $$\dot{\rho }_{SF}=\frac{n_{}m_{sf}}{\tau _{}}=\frac{m_{sf}}{\tau _{}}\frac{\alpha _B}{Q_0}\frac{T_h}{T_w}\frac{\rho _w}{\mu _w}n_h,$$ (33) where $`m_{sf}`$ is the mass of stars that form to produce $`Q_0`$ photons per unit time. If the stars follow a Salpeter mass function, with a minimum mass of 0.1 M and maximum mass of 100 M, then $$\frac{Q_0}{m_{sf}}4\times 10^{47}\mathrm{photons}\mathrm{M}_{}^1\mathrm{s}^1.$$ (34) The rate of mass depletion by star formation in the warm phase derived above is independent of the size distribution of the WFC’s, and depends only upon their filling factor and internal density. ## 5 Protogalactic Models ### 5.1 The Numerical Scheme We have performed one-dimensional Lagrangian hydrodynamic models of the evolution of the RHG subject to the physical processes discussed above. The models treat the warm and hot phases as separate fluids, evolving within a fixed external gravitational potential set by the dark matter. Angular momentum of the gas is not included. The models are applicable to the very early stages of star formation in galaxies or after mergers, or to cooling flows, before the gas has contracted sufficiently for its motion is dominated by angular momentum. During the later disk formation phase of a galaxy, or in the innermost regions of mergers or cooling flows, angular momentum would play an important role, whereas for the situations modelled here it would not strongly affect the physical processes under consideration. The momentum equation for the RHC $$\frac{\mathrm{d}u_h}{\mathrm{d}t}=4\pi R^2\frac{P_h}{m}gF_{drag},$$ (35) includes contributions from the gas pressure, $`P_h`$, the dark matter potential, $`g`$, and drag from the infalling WFC’s. In equation (35), $`u_h`$ is the radial velocity of the RHG, and $`4\pi \rho _hR^2\mathrm{d}R=\mathrm{d}m`$. For computational simplicity, the gravitational potential is taken to be isothermal out to a radius of 50 kpc, such that $`g=V_c^2/R`$. The magnitude of $`g`$ is held constant within a core radius $`R_{core}=1`$ kpc and dark matter is assumed to vanish outside of 50 kpc, such that $`g`$ becomes Keplerian. Although this prescription of $`g`$ is highly simplified, our results do not depend sensitively upon the detailed form of $`g`$. The form of the drag term is $$F_{drag}=\pi S^2\left(V_tu_h\right)^2\frac{\mathrm{d}n_c}{\mathrm{d}S}dS.$$ (36) As noted above, in models where high mass clouds dominate, we set $`V_t=V_c`$, whereas equation (5) is used in models where low mass clouds dominate. The equation of continuity for the RHG, $$\frac{\mathrm{d}\rho _h}{\mathrm{d}t}=\left[\frac{\rho _h}{R^2}\frac{}{R}\left(R^2u_h\right)+\dot{\rho }_c(R,t)\right],$$ (37) includes a sink term, due to condensation of gas out of the hot phase. The amount of mass lost from the hot phase depends upon the relative magnitudes of heating and cooling. If heating either exceeds or equals cooling in the RHG, then $`\dot{\rho }_c=0`$. If cooling exceeds heating in the RHG, however, mass is transferred from the RHG to the WFC’s at a rate given by $$\dot{\rho }_c=\frac{\rho _h\rho _{bal}}{\tau _c},$$ (38) where $`\rho _{bal}`$ is the density which would be required for heating to balance cooling in the hot phase. The energy equation for the RHG, $$\frac{\mathrm{d}e_h}{\mathrm{d}t}=\frac{\mathrm{\Gamma }\mathrm{\Lambda }_{brem}\mathrm{\Lambda }_{cond}}{\rho _h},$$ (39) includes heating by drag from the warm phase, bremsstrahlung emission, and conduction losses into the warm phase. The warm phase is not pressure-supported in the galactic halo, and so the mean velocity of the WFC evolves as $$\frac{\mathrm{d}u_w}{\mathrm{d}t}=g+F_{drag},$$ (40) where the drag term is similar to that in equation (35). A maximum infall velocity of $`fV_c`$ is imposed upon the warm phase. As discussed above, the parameter $`f`$ depends upon the details of the evolution of the warm phase. In our models, we consider values of $`f=1`$ and 0.1, thus examining the differences in our results for widely varying values of $`f`$. The mass of the warm phase increases by mass loss from the hot phase, and decreases by mass loss due to star formation, following the prescription given in equation (33), such that $$\frac{\mathrm{d}\rho _w}{\mathrm{d}t}=\dot{\rho }_c(R)\left[\frac{\rho _w}{R^2}\frac{}{R}\left(R^2\rho _w\right)+\dot{\rho }_{SF}\right].$$ (41) In principle, the smallest WFC’s (with $`S<S_{eu}`$ or $`S<S_{es}`$) are continually being evaporated into the RHG. But, if this phase transition increases $`\rho _h`$ such that cooling exceeds heating in the RHG, then an appropriate amount of hot gas would precipitate to retain thermal equilibrium in the RHG. Thus, the ratio $`<\rho _w>/\rho _h`$ depends indirectly on the size distribution of WFC’s which determines the conductive flux between the two phases. Finally, the WFC’s are assumed to be ionized and heated by the UV flux of the nearby massive stars, such that $`T_w10^4`$ K. They are also assumed to be in a pressure equilibrium with the RGH such that $`n_w=n_hT_h/T_w`$. Equations (35)-(41) completely describe the $`\rho _h`$, $`T_h`$, $`V_h`$, $`V_w`$, and $`<\rho _w>`$ distribution. From these quantities, $`T_w`$, $`\rho _w`$, as well as the heat transfer luminosity due to bremstrahlung and conduction, $`L_{\mathrm{brem}}`$ and $`L_{\mathrm{cond}}`$ can be obtained. We solve these equations numerically with a 1-D Lagrangian method (Richtmyer & Morton 1995). The warm and hot phases are treated as separate fluids, evolving and interacting as described above. The models described below use a grid of 400 cells for both fluids, with the cell spacing increasing outwards so as to give the best resolution in the region of interest, within 100 kpc. The outer boundary is set at 2 Mpc, with a constant pressure. The large outer radius is chosen so that conditions at that boundary cannot affect the solution. Variations by factors of two in the number of cells are found to have no significant effect upon the results. The hydrodynamic equations are evolved explicitly, while the energy equation (38) is solved implicitly, allowing the use of a simple Courant timestep criterion. The models are evolved until they reach a quasi-steady state, the results of which are shown in the following sections. ### 5.2 Model 1 Using the approach outlined above, we now consider several simple models to describe the evolution of gaseous ordinary matter in the halo of large galaxies. In principle, we should concurrently consider the cosmological evolution which led to the formation of the halo. However, such an investigation would require not only a large dynamical range in dissipationless interaction among dark matter structure but also the nonlinear evolution of ordinary matter in evolving and asymmetrical potentials. A number of authors have performed simulations of systems which follow the evolution of both dissipational and dissipationless components, either in the context of galaxy collisions, or galaxy formation (see, eg. Barnes & Hernquist 1991; Gerber, Lamb, & Balsara 1996; Lamb et al. 1997; Abel et al. 1998; Steinmetz & Navarro 1997, 1999; and references therein). Such studies are generally either unable to resolve the full dynamical range present in a young galaxy, or do not include the full range of physical processes discussed above. Fortunately, we are able to examine the importance of the processes discussed in § 3 and 4 by following the standard approach in galaxy formation simulations by breaking the problem into piecemeal tasks. In contrast to the conventional simulations (eg. Yepes et al. 1997; Elizondo et al. 1999), which mainly deal with the emergence of dark matter structure and in which the star formation process is parameterized, we assume the potential is already formed and focus our attention on the microphysics of the gas. In the first model, we take $`V_c=220`$ km s<sup>-1</sup>, and assume that $`\gamma >3`$, so that the appropriate equations to use are those derived for limit in which high mass clouds dominate. For these massive clouds, the gas drag effect is relatively weak and we initially adopt $`f=1`$, i.e. the WFC’s move inward at speeds of $`V_c`$. The initial conditions of the model are uniform density in the hot phase, with $`n_h=0.01`$ cm<sup>-3</sup>, giving a mass of $`10^{11}`$ M within a radius of 50 kpc. The average density of the warm phase is initially taken to be 0.01 times that of the hot phase. The temperature of the hot phase is initially set to 10<sup>6</sup> K. We have varied the initial conditions substantially with little change in the final distribution of the hot phase. The model was evolved for 10<sup>9</sup> yr, by which time it had reached an asymptotic quasi-stationary state. While the mass of gas continues to decrease due to star formation, the density distributions of the RHC and WFC’s attain asymptotic forms. These are shown in Figure 2. Figure 2a shows the final radial distributions of the density of the hot phase (solid curve), the average density of the warm phase (dashed curve), and the density of stars formed from the warm phase during the evolution (dot-dashed curve). The stellar density is calculated using the positions of stars where they are first formed. Because the stars will form on highly eccentric orbits, for which they spend most of their time near apocenter, the calculated stellar distribution should be representative of the observed distribution, after they have undergone violent relaxation (Aarseth, Lin, & Papaloizou 1988). As can be seen from the figure, $`\rho _h`$ has a fairly shallow radial dependence, varying approximately as $`R^{0.9}`$. The average density of the warm phase drops off somewhat more steeply, varying approximately as $`R^{1.2}`$. The stellar mass density has the steepest radial dependence, varying approximately as $`R^3`$. Figure 2b shows the final radial distributions of the pressure (solid curve) and temperature (dashed curve) of the hot phase. At radii near 100 kpc, the temperature remains at $`10^6`$ K. At smaller radii, however, the increasing density of the warm phase, with the corresponding increase in the heating rate, leads to temperatures several times higher. The combination of the density gradient with the shallow temperature gradient leads to a relatively steep pressure gradient. The magnitude of $`nT7\times 10^5`$ at 10 kpc and it varies approximately as $`R^1`$. The change in slope outside of 30 kpc is related to the form of the dark matter potential used in the model, which becomes Keplerian beyond 50 kpc. Figure 2c shows the radial dependences of the mass ratio of the warm to the hot phase (solid curve), and of the volume filling factor, $`f_V`$ of the WFC’s (dashed curve). We define $`f_V`$ as $`f_V\rho _w/\mu _wn_w`$, ie. it is the ratio of the average density of the warm gas to the internal density of the WFC’s. As can be seen, most of the mass of the system is in the WFC’s. It is also seen that $`f_V`$ shows very little variation with radius, and $`f_V0.01`$ for Rc < R < 100 < subscript𝑅𝑐𝑅 < 100R_{c}\mathrel{\vbox{ \offinterlineskip\hbox{$<$} \kern 1.29167pt\hbox{$\sim$} }}R\mathrel{\vbox{ \offinterlineskip\hbox{$<$} \kern 1.29167pt\hbox{$\sim$} }}100 kpc. The lack of variation in $`f_V`$ may seem somewhat counterintuitive, given that $`\rho _w`$ decreases at larger radii. It is a consequence, however, of the fact that the WFC’s are in pressure balance with the RHG, and the pressure of the latter increases rapidly towards smaller galactic radii. The clouds are therefore compressed into much smaller radii at small $`R`$. Note also that the ratio $`\rho _w/\rho _h>1`$ everywhere, such that most of the ordinary matter is in the WFC’s rather than residing in the RHG. This tendency arises because $`\tau _c<\tau _d`$ initially in this model. Figure 2d shows the radial dependences of the ratio of the depletion timescale of the WFC’s by star formation to the dynamical timescale (solid curve), and the ratio of cooling in the RHG by bremsstrahlung emission to that by conduction into the WFC’s (dashed curve). The depletion timescale is defined as $$\tau _{dep}=\frac{\rho _w}{\dot{\rho }_{sf}},$$ (42) while the dynamical timescale $$\tau _d=\frac{R}{V_c}.$$ (43) The ratio $`\tau _{dep}/\tau _d`$ gives a rough estimate of the ability of the system to replenish warm gas that has been lost to star formation (it is also replaced by mass lost from the warm phase, but this can be replenished only by motions at the sound speed of the hot phase, $`c_sV_c`$). As can be seen from the figure, both ratios are $`1`$ between $`R_c`$ and 100 kpc. Only in the inner regions of the halo, $`\tau _{dep}`$ becomes a significant fraction of $`\tau _d`$ such that efficient conversion from gas into stars is possible. For the values of $`T_h`$ and $`\rho _h`$ in the model, most of the flux from bremsstrahlung emission is composed of X-ray photons. In the low density outer regions of the halo, bremsstrahlung is not an efficient radiative process, and a comparable amount of energy is lost to the RHG via conduction. Within 100 kpc, the total luminosity of bremsstrahlung emission is $`L_{brem,Tot}=1.2\times 10^{11}`$ L. By comparison, the total energy lost from the RHG via conduction into the WFC’s (and immediately radiated from the WFC’s) is $`L_{cond,Tot}=1.3\times 10^{10}`$ L. In the models discussed below, a much greater fraction of the energy is lost via conduction. Based on these results, we do not expect these systems to be luminous source of X-rays. ### 5.3 Model 2, Changes with $`V_c`$ We compare the results of the previous model with one having a circular velocity of 100 km s<sup>-1</sup>. Such a model represents a dwarf system either at any epoch, or one which will become a larger system, but is at an early stage of evolution and has not yet undergone significant merging. In CDM models (Blumenthal et al. 1984; Thoul, & Weinberg 1996; Navarro, Frenk, & White 1997; Klypin, Nolthenius, & Primack 1997; Abel & Mo 1998), the first structures which form are significantly smaller than this. However, as can be seen from comparing the results below with those of Model 1, systems with significantly smaller values of V<sub>c</sub> are unlikely to contain a significant mass of RHG, and the assumptions of the models would not be applicable to such systems. The results of Model 2 are shown in Figure 3, which illustrate the same quantities as Figure 2. In Figure 3a, it can be seen that the densities are reduced significantly relative to the model with higher $`V_c`$. The relative sense of the steepness is retained, though the slopes are somewhat different, with $`\rho _h`$ varying approximately as $`R^{2/3}`$, $`\rho _w`$ as $`R^{1.5}`$ and $`\rho _{}`$ as $`R^{2.5}`$. As shown in Figure 3b, the temperature gradient in the RHG is again very shallow. The pressure of the RHG varies approximately as $`R^1`$, but it is lower than in the previous model, with $`n_hT_h3\times 10^4`$ at 10 kpc. The smaller value of $`V_c`$ relative to the previous model reduces the heating efficiency of the WFC’s. In order to balance heating and cooling, the density of the RHG is reduced relative to the previous model. This reduction in the RHG’s density relative to that of the warm phase is accompanied by increases in the mass ratio of the warm to the hot phase, and of the filling factor of the warm phase. Both of these quantities are larger by an order of magnitude relative to the earlier model, as can be seen in Figure 3c. We also note that within 1.5 kpc, where $`f_V>1`$, our model assumptions obviously break down. In this case, the WFC’s should be considered as a coherent entity rather than an assembly of smaller clouds. The relatively small pressure of the RHG leads to low internal density of the WFC’s, and a corresponding reduction in the star formation rate. This is reflected in the large values of $`\tau _{dep}/\tau _d`$ shown in Figure 3d. The reduced density of the RHG and the large $`<\rho _w>/\rho _h`$ ratio also increase the relative importance of conductive cooling of the RHG relative to bremsstrahlung emission. We find total luminosities out to 100 kpc of $`L_{cond,T}=7.2\times 10^8`$ L and $`L_{brem,T}=5.1\times 10^9`$ L. ### 5.4 Model 3, Changes with $`f`$ In Figure 4, we show the results of a model in which $`V_c=220`$ km s<sup>-1</sup>, and in which $`f=0.1`$, ie. the WFC’s move inwards with speeds of at most $`V_t=0.1V_c`$. This prescription would represent a situation in which disruption, coalescence, and coagulation occur on timescales short compared to the acceleration timescale of the clouds. Figure 4 shows the same quantities as Figure 1. The reduction of the infall velocity has a similar affect as a relatively small $`V_c`$, that is, a decrease in the heating efficiency of the RHG by the WFC’s, relative to the $`f=1`$ case. Both density and pressure of the RHG are relatively small in this model, similar to that seen in Model 2, which has a smaller value of $`V_c`$. With a relatively small density, the RHG loses more energy through conduction into the WFC’s than by bremsstrahlung everywhere in the system (Figure 4d). The total luminosities of the two processes are found to be $`L_{cond,Tot}=8.5\times 10^9`$ L and $`L_{brem,Tot}=1.4\times 10^8`$ L. The reduction in the infall velocity also leads to a shallower density profile in the WFC’s. Comparing with the earlier models, this tendency causes star formation to be enhanced at large radii and suppressed at smaller radii. Consequently, the stellar density becomes relatively shallow compared with previous results. ### 5.5 Model 4, Domination by small clouds We have also considered a model in which $`\gamma <4.5`$ so that the total mass and surface area of the cloud population is dominated by the small rather than the large clouds. In this limit, heating of the RHG, and conductive losses from the RHG into the WFC’s are taken from equations (24) and (30). The infall speed of the WFC’s is limited to $$V_t=\left(D_\rho \frac{S_{es}}{R}\right)^{\frac{1}{2}}V_c.$$ (44) Due to their relatively low infall velocities, individual small clouds are less efficient at heating the RHG than are large clouds. For a given $`\rho _w`$, therefore, a distribution dominated by small clouds leads to a significantly smaller value of $`G`$. In addition, the larger surface-to-volume ratio of a distribution of WFC’s dominated by small clouds leads to a much greater conductive cooling rate as compared to that in a distribution dominated by large clouds. The increase in its efficiency relative to the above models means that cooling by conduction into the WFC’s dominates bremsstrahlung everywhere in the RHG. The combined effects of higher cooling efficiency and lower heating efficiency lead to the depletion of the RHG and a much larger mass ratio of the WFC’s to the RHG relative to the models above. The lower density in the RHG relative to earlier models leads directly to smaller pressures, causing a reduction in both $`n_w`$ and the star formation rate. The change was so dramatic that we found $`f_V>1`$, ie. the RHG effectively does not exist in the system. But, in the absence of RHG, the WFC’s cannot fragment through Kelvin-Helmholz instability to offset their growth through cohesive collisions (see §3.3). Thus, the large $`<\rho _w>/\rho _h`$ ratio and the dominance of the small clouds appears to be incompatible. A self-consistent treatment of both phases, including the size distribution of the WFC’s will be examined and discussed in a future paper. The above models assume that $`\mathrm{\Lambda }_w=10^{22}`$ergs cm<sup>3</sup> s<sup>-1</sup>, a high value for warm clouds, near that of the peak from hydrogen emission. Such would be expected, given the stability arguments presented in § 3, if the small cloud population is dominated by those clouds which are just able to radiate away the energy conducted into them from the RHG. If, however, the small cloud population is dominated by slightly larger clouds, which can radiate more efficiently, then a much small value of $`\mathrm{\Lambda }_w`$ may be more appropriate. We have therefore repeated Model 4 using $`\mathrm{\Lambda }_w=10^{24}`$ergs cm<sup>3</sup> s<sup>-1</sup>. The decreased emission efficiency of the small clouds led to results that were significantly different from those found above, and closely resembled those found for Model 3. The difference between these two possible cases remains an uncertainty in this work, and shall be resolved in upcoming models which examine the structure of WFC’s embedded within the RHG. ## 6 Discussion In this paper, we examine the microphysics of gas dynamics in the early epoch of galactic evolution. Our objective is to provide a description of the dominant physical processes in ordinary matter which may be applied to a large class of galaxy formation models. For example, in the canonical hierarchical galaxy formation scenario, small dwarf-galaxy building blocks, containing non-interactive dark matter and gaseous ordinary matter, form first and subsequently merge to form larger entities such as our Galaxy. In the present analysis, we neglect the dynamical evolution of the dark matter halo which is undoubtedly important in determining not only the formation process but also the present kinematic properties of galaxies. This approximation enables us to focus our attention upon the evolution of gas in the early epoch of galactic evolution, which regulates the rate and location of star formation and therefore the light distribution and chemical properties of the emerging galaxies. We discuss below three possible applications of our results. ### 6.1 Galactic Stellar halo The results presented here show that thermal instability results in the formation of the residual halo gas (RHG), with density and pressure appropriate for quasi-hydrostatic and energy equilibria, and warm fragmentary clouds (WFC’s), which are pressure confined by RHG and heated by ionizing UV photons. Such systems have been proposed as being the source of observed Lyman-limit systems (Mo & Miralda-Escude 1996). The mass limits of the WFC’s are set by the same criteria as in the earlier work, but in this work we include many more details of the interactions between the phases. Inside $`10100`$ kpc, WFC’s in galaxies with masses comparable to the Galaxy are self-shielded from the extragalactic UV flux (see Figure 1). Unless these WFC’s are continually heated by the UV flux from nearby massive stars, further cooling reduces their temperature to $`10^2`$ K. Gravitational instability in the large WFC’s leads to spontaneous formation of stars among which the massive stars radiate UV photons, ionize their surroundings, and quench the formation of additional stars during their lifetime (Lin & Murray 1992). Through such a self-regulating feedback process, the maximum rate at which gas may be converted into stars is determined by the maintenance of an adequate UV flux to photoionize all the WFC’s (§4.2). In the regions far from the center of the halo where the density of the RHG is relatively low, the internal density of the WFC’s and their average density $`<\rho _w>`$ are also low. The effect of self-regulation limits the star formation timescale $`\rho _w/\dot{\rho }_{SF}>\tau _d`$. But, the star formation efficiency is much higher at smaller galactic distances. The results in Figures 2-4 show that within a few kpc (depending on the model), the stellar density $`\rho _{}`$ already exceeds both $`<\rho _w>`$ and $`\rho _h`$ after 1 Gyr. Once the stars are formed out of gas, they cannot dissipate their orbital energy such that their orbital radii cannot contract further. Thus, the interaction between RHG, WFC’s, and the newly formed massive stars essentially determines the asymptotic surface brightness distribution in galaxies. In Figure 5, we show the computed surface density profiles computed from the stellar distributions of Models 1-3. As can be seen, the models are fairly well fit by deVaucouleur profiles within the inner 50 kpc, in good agreement with the profiles of spheroidal systems. The onset of rapid and efficient star formation invalidates the instantaneous mixing and incremental gas-to-star conversion assumptions which are essential to the closed-box models for galactic chemical evolution and enrichment (cf Binney & Tremaine 1987), despite the apparent consistency between it and the observed metalicity distribution among population II stars (Chiappini et al. 1999). If these stars are formed in a series of starburst events in which a large fraction of the remaining gas is converted into stars on a timescale shorter than the dynamical timescale of the galaxy or the lifespan of the massive stars, the metallicity distribution would reflect the metallicity inhomogeneity in WFC’s. In this case, the deficiency of extremely metal poor stars (with \[Fe/H\] < < \mathrel{\vbox{ \offinterlineskip\hbox{$<$} \kern 1.29167pt\hbox{$\sim$} }}-3) would be consistent with an evolving initial mass function which gradually becomes less biased towards massive stars as the WFC’s are chemically enriched. A necessary condition for the formation of stars with long-lived, low-mass stars is small $`M_{crit}`$ (see §4.2). In a self-regulated environment, WFC’s with \[Fe/H\] > > \mathrel{\vbox{ \offinterlineskip\hbox{$>$} \kern 1.29167pt\hbox{$\sim$} }}-3 can spontaneously cool from $`10^4`$ K to $`10`$ K (so that Mcrit < 1M < subscript𝑀𝑐𝑟𝑖𝑡1subscript𝑀direct-productM_{crit}\mathrel{\vbox{ \offinterlineskip\hbox{$<$} \kern 1.29167pt\hbox{$\sim$} }}1M_{\odot}) between successive generations of nearby massive stars (with an interval $``$ a few $`10^6`$ yr) provided the initial $`n_w>0.11`$ (Hellsten & Lin 2000). The corresponding external pressure needed to confine such WFC’s is $`nT10^{34}`$. In Figures 2-4, we see that these values of $`n_hT_h`$ are attained outside 10 kpc, similar to the regions where the metal poor population II stars are located. In the above scenario, we have neglected the effect of magnetic fields which stablize cold clouds against gravitational instability through field-ion coupling and ion-neutral collisions (cf Shu 1985; Shu, Adams, & Lizano 1993). In dense cores of molecular clouds around the solar neighborhood, low-level ionization is maintained through cosmic ray heating and the field strength declines through ambipolar diffusion (Spitzer 1978; Shu 1985). If this process is important in WFC’s, the star formation rate would be much reduced from those illustrated in Figures 2-4. Collisions between WFC’s would occur at velocities in excess of 10 km s<sup>-1</sup>. Shock compression and rapid cooling near the collision interface could lead to a rapid expulsion of magnetic field, rendering its support ineffective. In the models discussed above, the covering factor of large WFC’s is of order unity. A large cloud would therefore collide with another large cloud only about once during a galactic crossing time. Collisions between large and small clouds would, however, occur much more frequently. In the limit that the size distribution of the clouds is dominated by relatively small clouds, the area covering factor is large, and small clouds collide with each other much more frequently. These issues are beyond the scope of the present investigation and they need to be thoroughly investigated in the future. ### 6.2 Formation of globular clusters The analysis presented here can also be applied to the formation of Galactic globular clusters (Lin & Murray 1996). Prior to the conversion of ordinary matter from gas into stars, the progenitors of these clusters were protocluster clouds (PCC’s). The chemical homogeneities within individual clusters and the large metallicity variation among different clusters suggest that PCC’s are a distinct entities which must be confined either by their own self-gravity or external pressure. But, if PCC’s are entirely bound by their own self-gravity, external UV heating would not be adequate to suppress thermal instability within them (see §3.3). The magnitude of $`M_{crit}`$ for the WFC’s is comparable to the mass of globular clusters. We identify these warm, marginally self-gravitating and partially pressure-confined WFC’s as PCC’s. At Galactic distances $`D330`$ kpc, PCC’s with mass ($`M`$) $`<M_{crit}(10^6M_{})`$ are confined by the the pressure of the RHG, $`n_hT_h10^{35}`$, depending on the Galactic halo structure during the epoch of cluster formation (Figs. 2-4). If these PCC’s are completely ionized and have a $`T_w10^4`$, then $`n_w0.110`$ cm<sup>-3</sup>. For these structural parameters, the UV flux needed is equivalent to that emitted by a few O5 stars at a distance comparable to or greater than their size (typical a few pc) (see Figure 1). (These stars could also reside within the PCC’s). These clouds could persist for a significant fraction of $`\tau _d`$ if the accretion of smaller clouds or condensation from the RHG is adequate to compensate for their mass loss due to stripping by the RHG. On the observational side, in order to verify that the PCC’s were pressure-confined, we first estimate their $`n_wT_w`$ from the current properties of globular clusters, averaged over their half-mass radii ($`r_h`$) (Murray & Lin 1996). We use these quantities because the stellar density and the velocity dispersion at $`r_h`$ do not change significantly during post-formation evolution. Extrapolation to the stage prior to star formation is, however, highly uncertain. If, after their formation, the young stellar objects undergo collapse and virialization from rest, the clouds’ initial radii ($`r_i`$) would be $`2r_h`$. Larger $`r_i`$ would be expected if star formation requires dissipative collisions and coagulation of substellar fragments (Murray & Lin 1996). But $`r_i`$ is unlikely to be larger than the tidal radii of the PCC’s, which are typically a few times larger than $`r_h`$. Thus, the initial density of the PCC’s may be 1-3 orders of magnitude smaller than the average cluster density at $`r_h`$ today. Based on the present velocity dispersion of the clusters, we infer the initial temperature of the PCC’s to be $`10^4`$ K, comparable to that expected if they were photoionized. From these estimates, we infer $`n_wT_w10^{25}`$ (Murray & Lin 1992). The dependence of the pressure upon galactic radius, $`D`$, is very poorly determined from the observational parameters. Of more significance is that the magnitude of the pressure inferred from the observations is very similar to that found in the RHG in our models (Figs. 2-4). From these results, and the cluster metallicities, we can also estimate the cooling time scale ($`\tau _{cc}`$) and dynamical time scale ($`\tau _{dc}`$), of the PCC’s. The ratio $`\tau _{cc}/\tau _{cd}`$ increases from $`10^4`$ near the Galactic bulge to $`1`$ at $`100`$ kpc. In most PCC’s, $`\tau _{cc}<<\tau _{cd}`$ and thermal equilibrium is only possible in the presence of external UV photons with a flux comparable to that required by self-regulated star formation in the halo. ### 6.3 X-ray Luminosity In the hierarchical galaxy formation process, coalescence of small stellar systems (dwarf galaxies) occurs within a few dynamical timescales. If the ordinary matter contained initially within these building blocks is to be heated through shock dissipation to the virial temperature of the common halo, the mergers would become luminous X-ray sources (Eke, Navarro, & Frenk 1998). But the results of our models indicate that; 1) thermal instability leads to the formation of WFC’s which contain most of the ordinary matter, and 2) even in the RHG, conduction may be a more efficient channel of heat loss than bremsstrahlung emission. As discussed in § 5.2-5.4, the resulting x-ray luminosities are small, ranging from 10<sup>11</sup> L in Model 1 to 10<sup>8</sup> L in Model 3. The spectra of the models are shown in Figure 6. The spectra of the low luminosity systems may be modified somewhat by supernova emission. Using the total star formation rates of Models 1-3, we estimate supernova luminosities of 10<sup>10</sup> L for Model 1, 6$`\times `$10<sup>9</sup> L for Model 2, and 5$`\times `$10<sup>8</sup> L for Model 3. Based upon these results, we expect relatively little X-ray luminosity to be released from regions with ongoing galaxy merging events, consistent with x-ray observations (Fabbiano & Schweizer 1995; Read & Ponman 1998). ### 6.4 Limitations In our first attempt to investigate the complex physics of multi phase gas dynamics during the early epoch of galaxy formation, we have adopted various simplifying assumptions such as 1-D spherical symmetry and power-law size distribution for WFC’s. These treatments can be improved with a self-consistent analysis of the WFC’s evolution which will be presented in a follow up paper. We have also neglected the evolution of dark matter which dominates the potential. The results in Figs 2-4 show that the distribution of pressure, density, filling factor of WFC’s and the star formation rate depend, though not sensitively on the potential. A full study of this problem will require the use of multi-dimensional integrated simulations of both dark and ordinary matter. While such work has been done in a cosmological context, the extreme temperature range spanned by the gas within an individual galaxy (at least four orders of magnitude) places severe demands upon a code to be able to resolve structures over many orders of magnitude in size. Our basic approach and the prescription provided here can be readily used in such investigations as they begin to be made. We thank Drs. A. Burkert, S. Faber, U. Hellsten, C. Frenk and A. Wolfe for valuable conversations. This research has been supported in part by the NSF through grant AST-9618548 and by NASA through an astrophysics theory program grant which supports a joint Center for Star Formation Studies at NASA-Ames Research Center, UC Berkeley, and UC Santa Cruz. This work was performed in part under the auspices of the U.S. Department of Energy by the University of California Lawrence Livermore National Laboratory under contract No. W-7405-Eng-48.
warning/0004/hep-th0004022.html	ar5iv	text	# Brane surgery: energy conditions, traversable wormholes, and voids ## 1 Introduction Branes, ubiquitous elements of any low-energy limit of string theory, have recently attracted much attention as essential ingredients of the semi-phenomenological Randall–Sundrum models . These models have been used to both ameliorate the “hierarchy problem” and to explore the possibility of “exotic” Kaluza–Klein theories with their infinitely large extra dimensions . Essential ingredients in these RS models are the existence of both positive and negative tension branes. Now a brane tension is normally thought of as being completely equivalent to an internal cosmological constant, and from the point of view of physics constrained to the brane this is certainly correct. However, from the higher-dimensional point of view (that is, as seen from the embedding space) this is not correct: For a (p+1)-brane embedded in (n+1) dimensions a brane tension leads to the stress energy $$T^{\mu \nu }=\mathrm{\Lambda }_Dg_{\mathrm{induced}}^{\mu \nu }\delta ^{np}(\eta ^a)=\mathrm{\Lambda }_D\left(g^{\mu \nu }\underset{a=1}{\overset{np}{}}n_a^\mu n_a^\nu \right)\delta ^{np}(\eta ^a),$$ (1.1) where the sum runs over the $`np`$ normals to the brane, and the $`\eta ^a`$ are suitable Gaussian normal coordinates. Contracting with a higher-dimensional null vector, $`k_\mu `$, we see $$T^{\mu \nu }k_\mu k_\nu =\mathrm{\Lambda }_Dg_{\mathrm{induced}}^{\mu \nu }k_\mu k_\nu \delta ^{np}(\eta ^a)=\mathrm{\Lambda }_D\left[\underset{a=1}{\overset{np}{}}(n_a^\mu k_\mu )^2\right]\delta ^{np}(\eta ^a).$$ (1.2) If the brane tension is negative, $`\mathrm{\Lambda }_D<0`$, and the null vector is even slightly orthogonal to the brane, then on the brane $$T^{\mu \nu }k_\mu k_\nu <0.$$ (1.3) That is, the embedding-space null energy condition (NEC) is violated. In fact, integrating across the brane, even the averaged null energy condition (ANEC) is violated. (Ipso facto, all the energy conditions are violated.) This is a classical violation of the energy conditions, which we shall soon see is even more profound than the classical violations due to non-minimally coupled scalar fields . In a recent series of papers we have made a critical assessment of the current status of the energy conditions, finding a variety of both classical and quantum violations of the energy conditions. We now see that uncontrolled violations of the energy conditions are also a fundamental and intrinsic part of any brane-based low-energy approximation to fundamental string theory. Among the possible consequences of these energy condition violations we mention the occurrence of traversable wormholes (violations of topological censorship), possible violations of the singularity theorems (more properly, evasions of the singularity theorems), and even the possibility of negative asymptotic mass. A particular example of this sort of phenomenon occurs in the (finite size) Randall–Sundrum models, where one has two parallel branes (our universe plus a hidden brane) of equal but opposite brane tension. One or the other of these branes (depending on whether one is considering the RS1 or RS2 model) violates the (4+1)-dimensional energy conditions and exhibits the “flare out” behaviour reminiscent of a traversable wormhole . That these branes do not quite represent traversable wormholes in the usual sense follows from the fact that the “throat” is an entire flat (3+1) Minkowski space, instead of the more usual $`R^1\times S^{d1}`$. Furthermore, in the infinite-size version of the Randall–Sundrum (RS2) model, where the hidden sector has been pushed out to hyperspatial infinity, our universe is itself represented by a positive-tension (3+1)-brane, which does not violate any (4+1)-dimensional energy conditions. The energy-condition violating brane has in this particular model been pushed out to hyperspatial infinity and discarded. Be that as it may, the occurrence of negative tension branes in modern semi-phenomenological models is generic, and a feel for the some of the peculiar geometries they can engender is essential to developing any deep understanding of the physics. In this particular paper we shall for illustrative purposes choose a particularly simple model: We work with a (3+1)-dimensional bulk, which contains a (2+1)-dimensional brane (of either positive or negative brane tension). We choose this particular model because it is sufficiently close to reality to make the points we wish to make as forcefully as possible, and because it arises naturally in certain types of fundamental string theory. While it is most often the case that fundamental string theories (or their various offspring: membrane models, M-theory, etc.) are formulated in either (9+1) or (10+1) dimensions,<sup>1</sup><sup>1</sup>1 In many specific cases the actual implementation is directly in terms of a Euclidean-signature $`10`$ or $`11`$ dimensional spacetime; with the underlying Lorentzian-signature reality hidden under several layers of scaffolding. this is not absolutely necessary: There is an entire industry based on formulating string theories directly in (3+1) dimensions, with the price that has to be paid being the inclusion of extra (1+1)-dimensional quantum fields propagating on the world-sheet .<sup>2</sup><sup>2</sup>2 Consider for example the bosonic string, which is most often viewed as a (1+1)-dimensional world sheet propagating in (25+1) dimensions: There is a trivial re-interpretation in which the bosonic string propagates in (3+1) dimensions and there are $`22`$ free scalar fields propagating on the world-sheet. These $`22`$ scalar fields are there just to soak up the conformal anomaly and make the theory manageable. If these scalar fields are now constrained by appropriate identifications the re-interpretation is less trivial — it is an example of the fact that compactifications of some of the dimensions of the higher-dimensional embedding spacetime that the world sheet propagates through can be traded off for a lower-dimensional uncompactified embedding spacetime plus interacting fields on the world-sheet. When this procedure is applied to superstrings the technical details are considerably more complex, but the basic result still holds. Now even in such a (3+1)-dimensional incarnation of string theory, open strings will terminate on D-branes (Dirichlet branes), and an effective theory involving the (3+1)-dimensional bulk plus (2+1), (1+1), and (0+1) dimensional D-branes (“domain walls”, “cosmic strings”, and “soliton-like particles”) can be contemplated as a low-energy approximation.<sup>3</sup><sup>3</sup>3 More traditional string theorists who absolutely insist on working directly in the higher-dimensional embedding space can view the current calculations as a particular toy model in which only selected sub-sectors of the grand total degrees of freedom are excited. Additionally, it should be borne in mind that many of the generic features of the analysis presented in this paper will extend mutatis mutandis to embedding spaces and branes of higher dimensionality. You do not want the bulk to have fewer than (3+1) dimensions since then bulk gravity is either completely or almost trivial. You do not want the bulk to have more than (10+1) dimensions since the model is then difficult to interpret in terms of fundamental string theory. For technical reasons (to be able to use the thin-shell formalism) you want the brane to be of co-dimension $`1`$, so if the bulk is (n+1)-dimensional the brane should be (\[n–1\]+1)-dimensional. Within these dimensional limitations, the qualitative features of this paper are generic. While D-branes are perhaps the most straightforward examples of membrane-like solitons in string theory, they do come with additional technical baggage: the most elementary implementation of D-branes occurs in bosonic string theories , but often D-branes are associated with specific implementations of supersymmetric string theories and carry various types of Ramond–Ramond or Neveu–Schwarz charge. There are in addition other types of brane-like configurations that sometimes arise in fundamental string theory such as non-dynamical “orientifold planes” , which generate gravitational fields corresponding to negative tensions, but which do not themselves exhibit internal dynamics. We will not delve further into this beastiary, but will instead content ourselves with the observation that the low-energy limit of fundamental string theory (of whatever persuasion) generically leads to an effective theory containing brane-like excitations. This overall picture is actually very similar to the notion of extended topological defects arising from symmetry breaking in point-particle field theories: There are many semi-phenomenological GUT-based point particle field theories that naturally contain domain walls, cosmic strings, and/or solitons. The key difference here is that point particle field theories inevitably lead to positive brane tensions, with negative brane tensions being energetically disfavoured (they correspond to an unnatural form of symmetry breaking that forces one to the top of the potential). The key difference in brane-based models is that there is no longer any particular barrier to negative brane tension — in fact negative brane tensions are ubiquitous, now being so commonly used as to almost not require explicit mention . Within the model we have chosen, we demonstrate that negative tension branes lead to traversable wormholes — in some cases to stable traversable wormholes. (Positive tension branes quite naturally lead to closed baby universes; these are not FLRW universes, and are not suitable for cosmology, but are perhaps of interest in their own right.) We also explore the possibility of viewing the brane as an actual physical boundary of spacetime, with the region on the “other side” of the brane being null and void. The basic tools used are the idea of “Schwarzschild surgery” as developed in (see also the more detailed presentation in ), which we first extend to “brane surgery”, specialize to “Reissner–Nordström–de Sitter” surgery, and then use to present an analysis of both static and dynamic spherically-symmetric (2+1)-dimensional branes in a (3+1)-dimensional Reissner–Nordström–de Sitter background geometry.<sup>4</sup><sup>4</sup>4 As we shall soon see, brane surgery is essentially a specific implementation of the Israel–Lanczos–Sen junction conditions of general relativity; as such it has been used implicitly in many brane-related papers (see for example ); the key difference in the present paper is in the details and in the questions we address. We find both stable and unstable traversable wormhole solutions, stable and unstable baby universes, and stable and unstable voids. ## 2 Brane surgery We start by considering a rather general static spherically symmetric geometry (not the most general, but quite sufficient for our purposes) $$\mathrm{d}s^2=F(r)\mathrm{d}t^2+\frac{\mathrm{d}r^2}{F(r)}+r^2\mathrm{d}\mathrm{\Omega }_2^2.$$ (2.4) To build the class of geometries we are interested in, we start by taking two copies of this geometry, truncating them at some time-dependent radius $`a(t)`$, and sewing the resulting geometries together along the boundary $`a(t)`$. The result is a manifold without boundary that has a “kink” in the geometry at $`a(t)`$. If we sew together the two external regions $`r(a(t),\mathrm{})`$, then the result is a wormhole spacetime with two asymptotic regions. On the other hand, if we sew together the two internal regions $`r(0,a(t))`$, then the result is a closed baby universe. At the “kink” $`a(t)`$ the spacetime geometry is continuous, but the radial derivative (and hence the affine connexion) has a step-function discontinuity. The Riemann tensor in this situation has a delta-function contribution at $`a(t)`$, and this geometry can be analyzed using the Israel–Lanczos–Sen “thin shell” formalism of general relativity . The relevant specific implementation of the thin-shell formalism can be developed by extending the formalism of and . Because of its relative simplicity we shall start with the static case $`a=\mathrm{constant}`$. ### 2.1 Brane statics The unit normal vector to the sphere $`a=\mathrm{constant}`$ is (depending on whether one is considering inward or outward normals) $$n^\mu =\pm (0,\sqrt{F(a)},0,0);n_\mu =\pm (0,\frac{1}{\sqrt{F(a)}},0,0).$$ (2.5) The extrinsic curvature (second fundamental form) can be written in terms of the normal derivative $$K_{\mu \nu }=\frac{1}{2}\frac{g_{\mu \nu }}{\eta }=\frac{1}{2}n^\sigma \frac{g_{\mu \nu }}{x^\sigma }=\pm \frac{1}{2}\sqrt{F(a)}\frac{g_{\mu \nu }}{r}.$$ (2.6) If we go to an orthonormal basis, the relevant components are<sup>5</sup><sup>5</sup>5 The use of an orthonormal basis makes it particularly easy to phase the calculation in terms of the physical density and physical pressure. $$K_{\widehat{t}\widehat{t}}=\frac{1}{2}\sqrt{F(r)}\frac{g_{tt}}{r}g^{tt}=\frac{1}{2}\sqrt{F(r)}\frac{F(r)}{r}\frac{1}{F(r)}=\frac{1}{2}F(r)^{1/2}\frac{F(r)}{r}\|_{r=a}.$$ (2.7) $$K_{\widehat{\theta }\widehat{\theta }}=\pm \frac{1}{2}\sqrt{F(r)}\frac{g_{\theta \theta }}{r}g^{\theta \theta }=\pm \frac{1}{2}\sqrt{F(r)}\frac{r^2}{r}\frac{1}{r^2}=\pm \frac{\sqrt{F(r)}}{r}\|_{r=a}.$$ (2.8) The discontinuity in the extrinsic curvature is related to the jump in the normal derivative of the metric as one crosses the brane $$\kappa _{\mu \nu }=K_{\mu \nu }^+K_{\mu \nu }^{}.$$ (2.9) In general, one could take the geometry on the two sides of the brane to be different $`[F^+(r)F^{}(r)]`$, but in the interests of clarity the present models will all be taken to have a $`Z_2`$ symmetry under interchange of the two bulk regions.<sup>6</sup><sup>6</sup>6 Remember that we have already decided to take the range of the $`r`$ coordinate to be either two copies of $`(a(t),\mathrm{})`$, corresponding to a wormhole; or two copies of $`(0,a(t))`$, corresponding to a baby universe. Then $`Z_2`$ symmetry corresponds to $`F^+(r)=F^{}(r)`$, with a kink in the geometry at $`r=a(t)`$. Our normal vectors do not flip sign as we cross the brane. Under these conditions $$\kappa _{\widehat{t}\widehat{t}}=F(r)^{1/2}\frac{F(r)}{r}\|_{r=a}.$$ (2.10) $$\kappa _{\widehat{\theta }\widehat{\theta }}=\pm 2\frac{\sqrt{F(r)}}{r}\|_{r=a}.$$ (2.11) The upper sign refers to a wormhole geometry where the two exterior regions have been sewn together (discarding the two interior regions), while the lower sign is relevant if one has kept the two interior bulk regions. The thin-shell formalism of general relativity relates the discontinuity in extrinsic curvature to the energy density and tension localized on the junction:<sup>7</sup><sup>7</sup>7 The numerical coefficients appearing herein are dimension-dependent (because of the implicit trace over the Ricci tensor and extrinsic curvature hidden in the Einstein equations). $$\sigma =\frac{1}{4\pi }\kappa _{\widehat{\theta }\widehat{\theta }}=\frac{1}{2\pi r}\sqrt{F(r)}\|_{r=a}.$$ (2.12) $$\theta =\frac{1}{8\pi }\left[\kappa _{\widehat{\theta }\widehat{\theta }}\kappa _{\widehat{t}\widehat{t}}\right]=\frac{1}{4\pi r}\frac{}{r}\left(r\sqrt{F(r)}\right)\|_{r=a}.$$ (2.13) If the material located in the junction is a “clean” brane (a brane in its ground state, without extra trapped matter in the form of stringy excitations), then its equation of state is $`\sigma =\theta `$ and the condition for a static brane configuration (either a wormhole or baby universe geometry) is simply $$\sigma =\theta 2\sqrt{F(r)}\|_{r=a}=\frac{}{r}\left(r\sqrt{F(r)}\right)\|_{r=a}\frac{}{r}\left(\frac{F(r)}{r^2}\right)\|_{r=a}=0.$$ (2.14) Thus we have a very simple result: static wormholes (baby universes) correspond to extrema of the function $`F(r)/r^2`$, though at this stage we have not yet made any assertions about stability or dynamics. The only difference between wormholes and baby universes is that for wormholes the brane tension must be negative, whereas for baby universes it is positive. It is instructive to note that the locations of these static brane solutions correspond to circular photon orbits in the original spacetime (and this is true for arbitrary $`F(r)`$). That is: at these static brane solutions any “particle” that is emitted form the brane, which then follows null geodesics (of the bulk spacetime), and which initially has no radial momentum, will just skim along the brane; never moving off into the bulk. (Note that this is a purely kinematic effect that occurs over and above any “trapping” due to stringy interactions between the brane and excited string states.) This may easily be verified by considering the photon orbits for arbitrary $`F(r)`$. The time-translation and rotational Killing vectors lead to conserved quantities $$(\frac{}{t},k)=ϵg_{tt}\frac{\mathrm{d}t}{\mathrm{d}\lambda }=ϵF\frac{\mathrm{d}t}{\mathrm{d}\lambda }=ϵ.$$ (2.15) $$(\frac{}{\varphi },k)=\mathrm{}g_{\varphi \varphi }\frac{\mathrm{d}\varphi }{\mathrm{d}\lambda }=\mathrm{}r^2\frac{\mathrm{d}\varphi }{\mathrm{d}\lambda }=\mathrm{}.$$ (2.16) Inserting this back into the condition that the photon momentum be a null vector, $`(k,k)=0`$, we see $$\left(\frac{\mathrm{d}r}{\mathrm{d}\lambda }\right)^2+\frac{F(r)\mathrm{}^2}{r^2}=ϵ^2.$$ (2.17) Now $`\lambda `$ is an arbitrary affine parameter, so we can reparameterize $`\lambda \mathrm{}\lambda `$ and define $`b=ϵ/\mathrm{}`$ to see that photon orbits are described by the equation $$\left(\frac{\mathrm{d}r}{\mathrm{d}\lambda }\right)^2+\frac{F(r)}{r^2}=b^2.$$ (2.18) The circular photon orbits (and at this stage we make no claims about stable versus unstable circular photon orbits) are, as claimed, at extrema of the function $`F(r)/r^2`$ (which coincide with the location of the stable brane configurations). ### 2.2 Brane dynamics If now the brane is allowed to move radially $`aa(t)`$, we start the analysis by first parameterizing the motion in terms of proper time along a curve of fixed $`\theta `$ and $`\varphi `$. That is: the brane sweeps out a world-volume $$X^\mu (\tau ,\theta ,\varphi )=(t(\tau ),a(\tau ),\theta ,\varphi ).$$ (2.19) The 4-velocity of the $`(\theta ,\varphi )`$ element of the brane can then be defined as $$V^\mu =(\frac{\mathrm{d}t}{\mathrm{d}\tau },\frac{\mathrm{d}a}{\mathrm{d}\tau },0,0).$$ (2.20) Using the normalization condition and the assumed form of the metric, and defining $`\dot{a}=\mathrm{d}a/\mathrm{d}\tau `$, $$V^\mu =(\frac{\sqrt{F(a)+\dot{a}^2}}{F(a)},\dot{a},0,0);V_\mu =(\sqrt{F(a)+\dot{a}^2},\frac{\dot{a}}{F(a)},0,0).$$ (2.21) The unit normal vector to the sphere $`a(\tau )`$ is $$n^\mu =\pm (\frac{\dot{a}}{F(a)},\sqrt{F(a)+\dot{a}^2},0,0);n_\mu =\pm (\dot{a},\frac{\sqrt{F(a)+\dot{a}^2}}{F(a)},0,0).$$ (2.22) The extrinsic curvature can still be written in terms of the normal derivative $$K_{\mu \nu }=\frac{1}{2}n^\sigma \frac{g_{\mu \nu }}{x^\sigma }.$$ (2.23) If we go to an orthonormal basis, the $`\widehat{\theta }\widehat{\theta }`$ component is easily evaluated $$K_{\widehat{\theta }\widehat{\theta }}=\pm \frac{1}{2}\sqrt{F(a)+\dot{a}^2}\frac{g_{\theta \theta }}{r}g^{\theta \theta }=\pm \frac{\sqrt{F(a)+\dot{a}^2}}{a}$$ (2.24) The $`\tau \tau `$ component is a little messier, but generalizing the calculation of or (which amounts to calculating the four-acceleration of the brane) quickly leads to<sup>8</sup><sup>8</sup>8 We do not repeat the details here since this calculation is now standard textbook fare , pp 182–183. If one wishes to avoid the need for this particular calculation one can instead work backwards from the conservation of stress-energy, together with the already-calculated expression for $`K_{\widehat{\theta }\widehat{\theta }}`$, to deduce an expression for $`K_{\widehat{\tau }\widehat{\tau }}`$. But if you choose this route you lose the opportunity to make a consistency check. $$K_{\widehat{\tau }\widehat{\tau }}=\frac{1}{2}\frac{1}{\sqrt{F(a)+\dot{a}^2}}\left(\frac{\mathrm{d}F(r)}{\mathrm{d}a}+2\ddot{a}\right)=\frac{\mathrm{d}}{\mathrm{d}a}\sqrt{F(a)+\dot{a}^2}.$$ (2.25) Applying the thin-shell formalism now gives: $$\sigma =\frac{1}{2\pi a}\sqrt{F(a)+\dot{a}^2}.$$ (2.26) $$\theta =\frac{1}{4\pi a}\frac{\mathrm{d}}{\mathrm{d}a}\left(a\sqrt{F(r)+\dot{a}^2}\right).$$ (2.27) These equations can easily be seen to be compatible with the conservation of the stress energy localized on the brane $$\frac{\mathrm{d}}{\mathrm{d}\tau }(\sigma a^2)=\theta \frac{\mathrm{d}}{\mathrm{d}\tau }(a^2).$$ (2.28) So as usual, two of these three equations are independent, and the third is redundant. From the above we see that traversable wormhole solutions, corresponding to the minus sign above, require negative brane tension (and so positive internal pressure and negative internal energy density). This is in complete agreement with where it was demonstrated that even for dynamical wormholes there must be violations of the null energy condition at (or near) the throat. If the material located in the junction is again assumed to be a “clean” brane ($`\sigma =\theta `$) then all the dynamics can be reduced to a single equation <sup>9</sup><sup>9</sup>9 And if the brane is not “clean” in this sense one only needs to keep track of one additional piece of information — the on-brane conservation equation (2.28). $$\dot{a}^2+F(a)=(2\pi \sigma )^2a^2.$$ (2.29) This single dynamical equation applies equally well to both wormholes and baby universes, the $``$ that shows up in the brane Einstein equations quietly goes away upon squaring — thus for questions of dynamics and stability these surgically constructed baby universes and wormholes can be dealt with simultaneously — the only difference lies in question of whether the brane tension is positive or negative. Note that we could re-write this dynamical equation as $$\left(\frac{\mathrm{d}\mathrm{ln}(a)}{\mathrm{d}\tau }\right)^2+\frac{F(a)}{a^2}=(2\pi \sigma )^2.$$ (2.30) From this it is clear that static solutions must be located at extrema of the function $`F(a)/a^2`$, in agreement with the static analysis. In the next section we shall make use of this general formalism by specializing $`F(r)`$ to the Reissner–Nordström–de Sitter form. We shall then exhibit some explicit solutions to these brane equations of motion, and perform the relevant stability analysis. ## 3 Reissner–Nordström–de Sitter surgery For the Reissner–Nordström–de Sitter geometry $$F(r)=1\frac{2M}{r}+\frac{Q^2}{r^2}\frac{\mathrm{\Lambda }}{3}r^2.$$ (3.31) It is then most instructive to write the dynamical equation in the form $$\left(\frac{\mathrm{d}\mathrm{ln}(a)}{\mathrm{d}\tau }\right)^2+V(a)=E,$$ (3.32) with a “potential” $$V(a)=\frac{F(a)}{a^2}=\frac{1}{a^2}\frac{2M}{a^3}+\frac{Q^2}{a^4}\frac{\mathrm{\Lambda }}{3},$$ (3.33) and an “energy” $$E=+(2\pi \sigma )^2.$$ (3.34) The extrema of this potential are easily located, their positions are independent of $`\mathrm{\Lambda }`$ and occur at $$r_\pm =\frac{3M}{2}\pm \sqrt{\left(\frac{3M}{2}\right)^22Q^2}.$$ (3.35) (As promised, these are indeed the locations of the circular photon orbits of the Reissner–Nordström–de Sitter geometry; note that the cosmological constant does not affect the location of these circular photon orbits.) The value of this potential at these extrema is somewhat tedious to calculate, we find $`V_\pm (M,Q,\mathrm{\Lambda })`$ $``$ $`V(r_\pm (M,Q))`$ $`=`$ $`{\displaystyle \frac{1}{4Q^2}}\left(1{\displaystyle \frac{9}{2}}{\displaystyle \frac{M^2}{Q^2}}+{\displaystyle \frac{27}{8}}{\displaystyle \frac{M^4}{Q^4}}\right)\pm {\displaystyle \frac{M}{4Q^6}}\left[\left({\displaystyle \frac{3M}{2}}\right)^22Q^2\right]^{3/2}{\displaystyle \frac{\mathrm{\Lambda }}{3}}.`$ Though it may not be obvious, the $`Q0`$ limit formally exists and is given by $$V_{}(M,Q0,\mathrm{\Lambda })\mathrm{};V_+(M,Q0,\mathrm{\Lambda })\frac{1}{27M^2}\frac{\mathrm{\Lambda }}{3}.$$ (3.37) The behaviour of the potential $`V(a)`$ is qualitatively: * $`V(a)+\mathrm{}`$ as $`a0`$, ($`Q0`$). * $`V(a)\frac{\mathrm{\Lambda }}{3}`$ as $`a+\mathrm{}`$. * There is at most one local minimum ($`V_{}`$ located at $`r_{}`$) and one local maximum ($`V_+`$ located at $`r_+`$). Two figures, where we have plotted $`V(a)`$ for two special cases, are provided in the discussion below. When looking for a stable brane solution we want to satisfy the following: 1. We want the local minimum to exist, and the brane to be located in its basin of attraction. 2. The energy must be at least equal to $`V_{}`$ (to even get a solution), and should not exceed $`V_+`$ (to avoid having the solution escape from the local well located around $`r_{}`$). 3. We also do not want (at least for now) the brane to fall inside (or even touch) any horizon the original Reissner–Nordström–de Sitter geometry might have — for two reasons: 1. Because if it does fall inside (or even touch) an event horizon the wormhole geometry is operationally indistinguishable from a Reissner–Nordström–de Sitter black hole and therefore not particularly interesting (but see the discussion regarding singularity avoidance later in this paper) whereas the baby universe geometry is for $`Q=0`$ doomed to a brief and unhappy life, and for $`Q0`$ is just plain weird. 2. For technical reasons ($`r`$ is now timelike and $`t`$ spacelike) a few key minus signs flip at intermediate steps of the calculation, more on this later. These physical constraints now imply: 1. To get a local minimum we need $`M>\sqrt{8/9}Q`$. 2. To then trap the solution, to make it one of bounded excursion, we need $$V_{}(M,Q,\mathrm{\Lambda })+(2\pi \sigma )^2V_+(M,Q,\mathrm{\Lambda }).$$ (3.38) 3. Horizon avoidance requires $`F(a)>0`$ over the entire range of motion; this implies $$V(a)=\frac{F(a)}{a^2}>0;V_{}(M,Q,\mathrm{\Lambda })>0.$$ (3.39) In view of this the horizon avoidance condition might more properly be called horizon elimination — horizons can be avoided if and only if the inner and outer horizons are actually eliminated. (We could however still have a cosmological horizon at very large distances, this cosmological horizon is never reached if the bounded excursion constraint is satisfied.) We can also explicitly separate out the cosmological constant to write the horizon elimination condition as $$\mathrm{\Lambda }<3V_{}(M,Q,\mathrm{\Lambda }0),$$ (3.40) which makes it clear that a powerful enough negative (bulk) cosmological constant is guaranteed to eliminate all the event horizons from the geometry. That these constraints can simultaneously be satisfied (at least in certain parameter regimes) can now be verified by inspection. The best way to proceed is to sub-divide the discussion into several special cases. ### 3.1 $`M>\|Q\|=0`$: There is one maxima (at $`a=3M`$) and no minimum. There are no stable solutions, though the “arbitrarily advanced civilization” posited by Morris and Thorne might like to try to artificially maintain the unstable static solution at $`a=3M`$. (This solution is unstable to both collapse and explosion.) Adding $`Q0`$ provides a hard core to the potential so that collapse is avoided (modulo the horizon crossing issue which must be dealt with separately). ### 3.2 $`M>\|Q\|0`$: There are now both a local maximum (at $`r_+<3M`$) and a global minimum (at $`r_{}>0`$). The potential is plotted in figure 1. Stable solutions exist, (both static stable solutions and stable solutions of bounded excursion), but since $`V_{}<0`$ ($`\mathrm{\Lambda }=0`$) at the global minimum horizon avoidance requires $$\mathrm{\Lambda }<3V_{}(M,Q,\mathrm{\Lambda }0)<0.$$ (3.41) That is, stable traversable wormhole or baby universe solutions exist only if the bulk is anti-de Sitter space (adS<sub>(3+1)</sub>) with a strong enough negative cosmological constant. Indeed if you consider the original geometry prior to brane surgery and extend it down to $`r=0`$ then for this choice of parameters (because of the large negative cosmological constant) you encounter a naked singularity. For the stable wormhole geometries based on this brane prescription this is not a problem since the naked singularity was in the part of the spacetime that you threw away in setting up the brane construction. (The baby universe models on the other hand, while stable, explicitly do contain naked singularities.)<sup>10</sup><sup>10</sup>10 This is part of a general pattern: The stable (or even merely static) brane configurations that do not possess naked singularities in the bulk region are the wormhole configurations with negative brane tension. This observation also applies to the other sub-cases discussed below. This is compatible with the discussion of Chamblin, Perry, and Reall who discovered qualitatively similar behaviour for (8+1)-dimensional branes in a (9+1)-dimensional bulk. Specifically, they found that static (8+1)-dimensional brane configurations with positive brane tension led to naked singularities in the bulk, and that eliminating the naked singularities forced the adoption of negative brane tension (and implicitly a wormhole configuration). This observation also serves to buttress our previous comments to the effect that the qualitative features of the calculations presented in this paper are generic, and are not just limited to (2+1) branes in (3+1) dimensions. A particularly simple sub-class of these solutions occurs when the bulk cosmological constant is tuned to a special value in terms of the brane tension. This is the analog of the Randall–Sundrum fine tuning and corresponds to a zero “effective” cosmological constant, in the sense that the brane equation of motion can be rearranged and reinterpreted as being governed by $`E=0`$ and $$\mathrm{\Lambda }_{\mathrm{effective}}=\mathrm{\Lambda }+3(2\pi \sigma )^2.$$ (3.42) If this $`\mathrm{\Lambda }_{\mathrm{effective}}`$ is now tuned to zero $$\mathrm{\Lambda }=3(2\pi \sigma )^2<3V_{}(M,Q,\mathrm{\Lambda }0)<0.$$ (3.43) ### 3.3 $`M=\|Q\|`$: There are still both a local maximum (at $`r_+=2M`$) and a global minimum (at $`r_{}=M`$). Stable solutions exist. Since now $`V_{}(\mathrm{\Lambda }0)=0`$at the global minimum horizon avoidance requires anti de Sitter space with an arbitrarily weak cosmological constant. (And again this is an example of horizon elimination.) ### 3.4 $`M(\sqrt{8/9}\|Q\|,\|Q\|)`$: There are still both a local maximum (at $`r_+<2M`$) and a global minimum (at $`r_{}>M`$). The potential is plotted in figure 2. Stable solutions exist. Since now $`V_{}(\mathrm{\Lambda }0)>0`$ at the global minimum, horizon avoidance can be achieved with zero cosmological constant in the bulk. For instance, picking $$\mathrm{\Lambda }=0;(2\pi \sigma )^2=V_{}(M,Q,\mathrm{\Lambda }0),$$ (3.44) yields the stable static solution at $`r_{}`$. This is perhaps the most “physical” of these traversable wormholes in that it resides in an asymptotically flat spacetime. ### 3.5 $`M=\sqrt{8/9}\|Q\|`$: The maximum and minimum merge into a single point of inflection (at $`r_\pm =3M/2`$). There are no stable solutions. All the solutions exhibit runaway to large radius. ### 3.6 $`M(\sqrt{8/9}\|Q\|,0)`$: There is not even a point of inflection: the potential is monotonic decreasing. There are no stable solutions. ### 3.7 $`M=0`$; $`Q0`$: There is not even a point of inflection: the potential is monotonic decreasing. There are no stable solutions. ### 3.8 $`M<0`$: Letting the central mass $`M`$ go negative is not helpful — $`M<0`$ helps stabilize against collapse, but actually destroys the possibility of stable solutions because the location of “extrema” $`r_\pm `$ is pushed to unphysical nominally negative values of the radius. ### 3.9 Baby bangs? The fact that so many of these baby universe models are unstable to explosion is intriguing, and potentially of phenomenological interest. While these particular baby-universe models are not suitable cosmologies for our own universe, we believe that more realistic scenarios can be developed. ### 3.10 Singularity avoidance? We have so far sought to implement horizon avoidance in our models: we have sought conditions that would prevent the brane from falling through or even touching any horizon that might be present in the underlying pre-surgery spacetime. Suppose we now relax that constraint. The best way to analyze the situation is to note that inside the horizon (more precisely between the outer horizon and the inner horizon) the pre-surgery metric can be written in the form $$\mathrm{d}s^2=+\|F(r)\|\mathrm{d}t^2\frac{\mathrm{d}r^2}{\|F(r)\|}+r^2\mathrm{d}\mathrm{\Omega }_2^2.$$ (3.45) The calculation of the four-velocity, normal, extrinsic curvatures, and their discontinuities can be repeated, with the result that in this region \[$`F(r)<0`$\] $$V^\mu =(\frac{\sqrt{\dot{a}^2\|F(a)\|}}{\|F(a)\|},\dot{a},0,0);n^\mu =\pm (\frac{\dot{a}}{\|F(a)\|},\sqrt{\dot{a}^2\|F(a)\|},0,0);$$ (3.46) and $$\sigma =\frac{1}{2\pi a}\sqrt{\dot{a}^2\|F(a)\|}.$$ (3.47) After rearrangement this leads to the same dynamical equation as before $$\left(\frac{\mathrm{d}\mathrm{ln}(a)}{\mathrm{d}\tau }\right)^2+\frac{F(a)}{a^2}=(2\pi \sigma )^2.$$ (3.48) So that all of our previous arguments can be extended inside the event horizon. A few key observations: * The two turning points occur at $`F(a)/a^2=(2\pi \sigma )^2>0`$. Thus $`F(a)>0`$ at the turning points. So if there are horizons present (that is, if $`F(a)=0`$ has solutions $`r_{\mathrm{horizon}}^\pm r_\pm `$), and one is in the potential well near $`r_{}`$, then one turning point will be outside the outer horizon, and the second turning point will be inside the inner horizon. * Even though the brane oscillation will take finite proper time $`\tau `$ this corresponds to infinite $`t`$-parameter time — when the brane re-emerges from the outer horizon it will emerge from a past outer horizon of a “future incarnation” of the universe; the brane will not re-emerge into our own universe. (For simplicity you may wish to set $`\mathrm{\Lambda }=0`$ and consider the Penrose diagram of the maximally extended Reissner–Nordström geometry as presented, for instance, on page 158 of Hawking and Ellis . A partial Penrose diagram for Reissner–Nordström–de Sitter may be found in . See also figure 3.) * Operationally, from “our” asymptotically flat region, once the brane passes the horizon the geometry will be indistinguishable from an ordinary Reissner–Nordström–de Sitter black hole. * The original pre-surgery spacetime has two asymptotic regions, two outer horizons, and two inner horizons, which are then repeated an infinite number of times in the maximal analytic extension. If the brane starts out in the rightmost asymptotic region and falls through the right (future) outer horizon, then you can quickly convince yourself that it must pass through the left inner horizon (twice, once on the way in, and once again on the rebound) before moving back out through the right (past) outer horizon back into the (next incarnation of) the right asymptotic region. (See figure 3.) * The wormhole geometry based on this brane surgery is an explicit example of partial evasion of the usual singularity theorems . (We say evasion, not violation, because the presence of the negative tension brane vitiates the usual hypotheses used in proving the singularity theorems.) The wormhole geometry certainly has trapped surfaces once the brane falls inside the horizon, but by construction there is no left curvature singularity. (The right curvature singularity is still there, and the right inner horizon is still a Cauchy horizon.)<sup>11</sup><sup>11</sup>11 If you think of the Reissner–Nordström–de Sitter geometry as arising from gravitational collapse of an electrically charged star, then it is the left curvature singularity (which is eliminated by the present construction) that would arise from the central density of the star growing to infinity. The right curvature singularity (which is unaffected by the present construction) has a totally different genesis as it arises in a matter-free region due to gravitational focussing of the electromagnetic field. Note that this is a idealized statement appropriate to “clean” wormhole universes containing only a few test particles of matter: in any more realistic model where the universe contains a finite amount of radiation, inner event horizons are typically unstable to a violent blue shift instability, and are typically converted by back-reaction effects to some sort of curvature singularity . This process however, lies far beyond the scope of the usual singularity theorems. If you wish to eliminate both left and right singularities a more drastic fix is called for: You will need to use a (3+0)-dimensional brane, something you might call an instanton-brane because it represents a spacelike hypersurface through the spacetime — at early times there’s nothing there, the brane “switches on” for an instant, and then it’s gone again. The simplest example of such a instanton-brane is to place one at $`r_{}`$, the static minimum of the potential $`V(a)`$.<sup>12</sup><sup>12</sup>12 Although this is a static minimum of the usual $`V(a)`$ it is in the present context not stable. This arises because for a spacelike shell the overall sign of the potential flips. If there are event horizons then this minima will be inside the event horizon (between inner and outer horizons) and a hypersurface placed at $`r_{}`$ will be spacelike. Placing the instanton-brane at this location will eliminate both singularities and both inner horizons — you are left with two asymptotic regions and two (outer) event horizons, infinitely repeated. More generally one could think of an instanton-brane described by a location $`a(\mathrm{})`$ where $`\mathrm{}`$ is now proper length along the brane (and the notion of dynamics is somewhat obscure). The spacelike tangent and timelike normal are now (outside the horizon) $$V^\mu =(\frac{\sqrt{(\mathrm{d}a/\mathrm{d}\mathrm{})^2F(a)}}{F(a)},\frac{\mathrm{d}a}{\mathrm{d}\mathrm{}},0,0);n^\mu =\pm (\frac{1}{F(a)}\frac{\mathrm{d}a}{\mathrm{d}\mathrm{}},\sqrt{(\mathrm{d}a/\mathrm{d}\mathrm{})^2F(a)},0,0);$$ (3.49) and a brief computation yields $$\sigma =\frac{1}{2\pi a}\sqrt{(\mathrm{d}a/\mathrm{d}\mathrm{})^2F(a)}.$$ (3.50) This can be rearranged to give $$\left(\frac{\mathrm{d}\mathrm{ln}(a)}{\mathrm{d}\mathrm{}}\right)^2\frac{F(a)}{a^2}=(2\pi \sigma )^2.$$ (3.51) So the net result is that for an instanton-brane the sign of the potential has flipped, but that of the brane contribution to the energy has not. (And exactly the same result continues to hold inside the horizon, a few intermediate signs flip, but that’s all.) In summary: certain varieties of brane wormhole provide explicit evasions (either partial or complete) of the usual singularity theorems. ## 4 Voids: the brane as a spacetime boundary A somewhat unusual feature of brane physics is that the brane could also be viewed as an actual physical boundary to spacetime, with the “other side” of the brane being null and void. In general relativity as it is normally formulated the notion of an actual physical boundary to spacetime (that is, an accessible boundary reachable at finite distance) is anathema. The reason that spacetime boundaries are so thoroughly deprecated in general relativity is that they become highly artificial special places in the manifold where some sort of boundary condition has to be placed on the physics by an act of black magic. Without such a postulated boundary condition all predictability is lost, and the theory is not physically acceptable. Since there is no physically justifiable reason for picking any one particular type of boundary condition (Dirichlet, Neumann, Robin, or something more complicated), the attitude in standard general relativity has been to exclude boundaries, by appealing to the cosmic censor whenever possible and by hand if necessary. The key difference when a brane is used as a boundary is that now there is a specific and well-defined boundary condition for the physics: D-branes (D for Dirichlet) are defined as the loci on which the fundamental open strings end (and satisfy Dirichlet-type boundary conditions). D-branes are therefore capable (at least in principle) of providing both a physical boundary and a plausible boundary condition for spacetime. For Neveu–Schwarz branes the boundary conditions imposed on the fundamental string states are more complicated, but they still (at least in principle) provide physical boundary conditions on the spacetime. When it comes to specific calculations, this may however not be the best mental picture to have in mind — after all, how would you try to calculate the Riemann tensor for the edge of spacetime? And what would happen to the Einstein equations at the edge? There is a specific trick that clarifies the situation: Take the manifold with brane boundary and make a second copy, then sew the two manifolds together along their respective brane boundaries, creating a single manifold without boundary that contains a brane, and exhibits a $`Z_2`$ symmetry on reflection around the brane. Because this new manifold is a perfectly reasonable no-boundary manifold containing a brane, the gravitational field can be analyzed using the usual thin-shell formalism of general relativity : The metric is continuous, the connection exhibits a step-function discontinuity, and the Riemann curvature a delta-function at the brane. The dynamics of the brane can then be investigated in this $`Z_2`$-doubled manifold, and once the dynamical equations and their solutions have been investigated the second surplus copy of spacetime can quietly be forgotten. In particular, all the calculations we have performed for the spherically symmetric wormholes of this paper apply equally well to spherically symmetric holes in spacetime (not black holes, actual voids in the manifold), with the edge of the hole being a brane — we deduce the existence of a large class of stable void solutions, and an equally large class of unstable voids that either collapse to form black holes, or explode to engulf the entire universe. Equally well, the baby universes of the preceding section can, under this new physical interpretation of the relevant mathematics, be used to investigate finite volume universes with boundary. The bulk of the physical universe now lies in the range $`r(0,a)`$, and the edge of the universe is located at $`a`$. Again, we deduce the existence of a large class of stable baby universes with boundary, and an equally large class of unstable baby universes that either collapse to singularity, or explode to provide arbitrarily large universes. Note that these particular exploding universes are not FLRW universes, and are not suitable cosmologies for our own universe. Nevertheless, this notion of using a brane as an actual physical boundary of spacetime is an issue of general applicability, and we hope to return to this topic in future publications. ## 5 Discussion The main point of this paper is that in the brane picture there is nothing wrong with the notion of a negative brane tension, and that once branes of this type are allowed to contribute to the stress-energy, the class of solutions is greatly enhanced, now including many quite peculiar beasts not normally considered to be part of standard general relativity. As specific examples, the energy condition violations caused by negative tension branes allow one to construct classical traversable wormholes, at least some of which (as we have seen) are actually dynamically stable. Now for spherically symmetric wormholes of the type considered in this paper, attempting to cross from one universe to the other requires the traveller to cross the brane, a process which is likely to prove disruptive of the traveller’s internal structure, well being, and overall health. This problem, or rather the no-brane analog of this problem, was already considered by Morris and Thorne in their pioneering work on traversable wormholes . A possible resolution comes from the fact that spherical symmetry is a considerable idealization: One of the present authors has demonstrated that if one uses negative tension cosmic strings instead of negative tension domain walls, then it is possible to construct traversable wormhole spacetimes that do not possess spherical symmetry, and contain perfectly reasonable paths from one asymptotic region to the other that do not involve personal encounters with any form of “exotic matter” . (See also the extensive discussion in .) In a brane context this means we should consider the possibility of a negative tension (1+1)-dimensional brane in (3+1)-dimensional spacetime. Now the peculiarities attendant on widespread violations of the energy conditions are not limited to violations of topological censorship; as we have seen there is also the possibility of violating (evading) the singularity theorem. If this is not enough, then it should be borne in mind that without some form of energy condition we do not have a positive mass theorem. (Looking out into our own universe, we do have a positive mass observation, but it would be nice to be able to deduce this from general principles.) A discussion of some of the peculiarities attendant on negative asymptotic mass can be found in the early work of Bondi , and a possible observational signal (particular types of caustics in the light curves due to gravitational lensing) has been pointed out by Cramer et al. . Finally, energy condition violations are also the sine qua non for the Alcubierre “warp drive” . In summary, all of these somewhat peculiar geometries, which were investigated within the general relativity community more with a view to understanding the limitations of general relativity (and more specifically, of semiclassical general relativity) than in the expectation that they actually exist in reality, are now seen to automatically be part and parcel of the brane models currently being considered as semi-phenomenological models of empirical reality. ## Acknowledgments The research of CB was supported by the Spanish Ministry of Education and Culture (MEC). MV was supported by the US Department of Energy. We wish to thank Harvey Reall and Sumit Das for their comments and interest.
warning/0004/hep-ph0004189.html	ar5iv	text	# DESY 00–034 TTP00–05 hep-ph/0004189 March 2000 RunDec: a Mathematica package for running and decoupling of the strong coupling and quark masses ## Program Summary * Title of program: RunDec * Available from: http://www-ttp.physik.uni-karlsruhe.de/Progdata/ttp00/ttp00-05/ * Computer for which the program is designed and others on which it is operable: Any work-station or PC where Mathematica is running. * Operating system or monitor under which the program has been tested: UNIX, Mathematica 4.0 * No. of bytes in distributed program including test data etc.: $`65000`$ * Distribution format: ASCII * Keywords: Quantum Chromodynamics, running coupling constant, running quark mass, on-shell mass, $`\overline{\mathrm{MS}}`$ mass, decoupling of heavy particles * Nature of physical problem: The values for the coupling constant of Quantum Chromodynamics, $`\alpha _s^{(n_f)}(\mu )`$, actually depends on the considered energy scale, $`\mu `$, and the number of active quark flavours, $`n_f`$. The same applies to light quark masses, $`m_q^{(n_f)}(\mu )`$, if they are, e.g., evaluated in the $`\overline{\mathrm{MS}}`$ scheme. In the program RunDec all relevant formulae are collected and various procedures are provided which allow for a convenient evaluation of $`\alpha _s^{(n_f)}(\mu )`$ and $`m_q^{(n_f)}(\mu )`$ using the state-of-the-art correction terms. * Method of solution: RunDec uses Mathematica functions to perform the different mathematical operations. * Restrictions on the complexity of the problem: It could be that for an unphysical choice of the input parameters the results are nonsensical. * Typical running time: For all operations the running time does not exceed a few seconds. ## 1 Introduction Quantum Chromodynamics (QCD) is nowadays well established as the theory of strong interaction within the Standard Model of elementary particle physics. In recent years there has been a wealth of theoretical results (for a review see ). At the same time perturbative QCD has been extremely successful in describing the experimental data with high precision. The fundamental quantity of QCD is the so-called beta function which connects the value of the strong coupling constant, $`\alpha _s(\mu )`$, at different energy scales $`\mu `$. It is thus particularly important to know the beta function as precise as possible. In the four-loop corrections were evaluated allowing for a consistent running at order $`\alpha _s^4`$. In the majority of all computations performed in QCD the $`\overline{\mathrm{MS}}`$ renormalization scheme is adopted. In this scheme the Appelquist-Carazzone decoupling theorem is not directly applicable. When crossing flavour thresholds, it is thus important to perform the decoupling “by hand”. In order to be consistent, four-loop running must go along with the three-loop decoupling relation which was evaluated in . Similar considerations are also valid for quark masses. Also here the renormalization group function is available up to the four-loop level and the corresponding decoupling relation up to order $`\alpha _s^3`$ (see also ). In this paper all relevant formulae are collected which are necessary for the running and decoupling of $`\alpha _s`$ and for quark masses. Their proper use is discussed and easy-to-use Mathematica procedures collected in the package RunDec are provided. Their handling is described and examples are given. The outline of the paper is as follows. In the next Section the formulae are presented which are needed for the running of the strong coupling constant up to the four-loop level. The corresponding equations for the quark masses are presented in Section 3. In addition the conversion formulae between the $`\overline{\mathrm{MS}}`$ and on-shell scheme are discussed in some detail. Section 4 is concerned with the decoupling of the strong coupling and quark masses. Finally, in Section 5, the most important procedures of the package RunDec are described in an easy-to-use way. For most practical applications they should be sufficient. In the Appendix the complete collection of procedures is given. ## 2 Running strong coupling constant The beta function governing the running of the coupling constant of QCD is defined through $`\mu ^2{\displaystyle \frac{\mathrm{d}}{\mathrm{d}\mu ^2}}{\displaystyle \frac{\alpha _s^{(n_f)}(\mu )}{\pi }}`$ $`=`$ $`\beta ^{(n_f)}\left(\alpha _s^{(n_f)}\right)={\displaystyle \underset{i0}{}}\beta _i^{(n_f)}\left({\displaystyle \frac{\alpha _s^{(n_f)}(\mu )}{\pi }}\right)^{i+2},`$ (1) where $`n_f`$ is the number of active flavours. The coefficients are given by $`\beta _0^{(n_f)}`$ $`=`$ $`{\displaystyle \frac{1}{4}}\left[11{\displaystyle \frac{2}{3}}n_f\right],`$ $`\beta _1^{(n_f)}`$ $`=`$ $`{\displaystyle \frac{1}{16}}\left[102{\displaystyle \frac{38}{3}}n_f\right],`$ $`\beta _2^{(n_f)}`$ $`=`$ $`{\displaystyle \frac{1}{64}}\left[{\displaystyle \frac{2857}{2}}{\displaystyle \frac{5033}{18}}n_f+{\displaystyle \frac{325}{54}}n_f^2\right],`$ $`\beta _3^{(n_f)}`$ $`=`$ $`{\displaystyle \frac{1}{256}}[{\displaystyle \frac{149753}{6}}+3564\zeta _3+({\displaystyle \frac{1078361}{162}}{\displaystyle \frac{6508}{27}}\zeta _3)n_f`$ (2) $`+({\displaystyle \frac{50065}{162}}+{\displaystyle \frac{6472}{81}}\zeta _3)n_f^2+{\displaystyle \frac{1093}{729}}n_f^3].`$ $`\zeta `$ is Riemann’s zeta function, with values $`\zeta _2=\pi ^2/6`$ and $`\zeta _3\mathrm{1.202\hspace{0.17em}057}`$. It is convenient to introduce the following notation: $`b_i^{(n_f)}`$ $`=`$ $`{\displaystyle \frac{\beta _i^{(n_f)}}{\beta _0^{(n_f)}}},`$ $`a^{(n_f)}(\mu )`$ $`=`$ $`{\displaystyle \frac{\alpha ^{(n_f)}(\mu )}{\pi }}.`$ (3) In the following the labels $`\mu `$ and $`n_f`$ are omitted if confusion is impossible. Integrating Eq. (1) leads to $`\mathrm{ln}{\displaystyle \frac{\mu ^2}{\mathrm{\Lambda }^2}}`$ $`=`$ $`{\displaystyle \frac{da}{\beta (a)}}`$ (4) $`=`$ $`{\displaystyle \frac{1}{\beta _0}}\left[{\displaystyle \frac{1}{a}}+b_1\mathrm{ln}a+(b_2b_1^2)a+\left({\displaystyle \frac{b_3}{2}}b_1b_2+{\displaystyle \frac{b_1^3}{2}}\right)a^2\right]+C,`$ where an expansion in $`a`$ has been performed. The integration constant is conveniently split into $`\mathrm{\Lambda }`$, the so-called asymptotic scale parameter, and $`C`$. The conventional $`\overline{\mathrm{MS}}`$ definition of $`\mathrm{\Lambda }`$, which we shall adopt in the following, corresponds to choosing $`C=(b_1/\beta _0)\mathrm{ln}\beta _0`$ . Iteratively solving Eq. (4) yields $`a`$ $`=`$ $`{\displaystyle \frac{1}{\beta _0L}}{\displaystyle \frac{b_1\mathrm{ln}L}{(\beta _0L)^2}}+{\displaystyle \frac{1}{(\beta _0L)^3}}\left[b_1^2(\mathrm{ln}^2L\mathrm{ln}L1)+b_2\right]`$ (5) $`+{\displaystyle \frac{1}{(\beta _0L)^4}}\left[b_1^3\left(\mathrm{ln}^3L+{\displaystyle \frac{5}{2}}\mathrm{ln}^2L+2\mathrm{ln}L{\displaystyle \frac{1}{2}}\right)3b_1b_2\mathrm{ln}L+{\displaystyle \frac{b_3}{2}}\right],`$ where $`L=\mathrm{ln}(\mu ^2/\mathrm{\Lambda }^2)`$ and terms of $`𝒪(1/L^5)`$ have been neglected. $`\mathrm{\Lambda }`$ is defined in such a way that Eq. (5) does not contain a term proportional to $`(\mathrm{const}./L^2)`$ . The canonical way to compute $`a(\mu _2)`$ when $`a(\mu _1)`$ is given for a fixed number of flavours is as follows: 1. Determine $`\mathrm{\Lambda }`$. There are several possibilities to do this. One could, e.g., use the explicit solution given in Eq. (4). Another possibility is the use of (5) and solve the equation iteratively for $`\mathrm{\Lambda }`$. Furthermore the first line of (4) could be used and the integral could be solved numerically without performing any expansion in $`\alpha _s`$. We will see in the examples below that the numerical differences are small. 2. $`a(\mu _2)`$ is computed with the help of Eq. (5) where the value of $`\mathrm{\Lambda }`$ is inserted and $`\mu `$ is set to $`\mu _2`$. It is also possible to avoid the introduction of $`\mathrm{\Lambda }`$ in intermediate steps and to solve the differential equation (1) numerically using $`a(\mu )\|_{\mu =\mu _1}=a(\mu _1)`$ as initial condition. This convention requires the knowledge of both $`\alpha _s`$ and the scale $`\mu `$ in order to determine $`\alpha _s`$ at the new scale. Frequently, $`\mu =M_Z`$ is used as reference scale. On the other hand $`\mathrm{\Lambda }`$ plays the role of an universal parameter which at the same time sets the characteristic scale of QCD. At this point it is instructive to consider an example. Let us assume that $`\alpha _s`$ is given at the $`Z`$-boson scale: $`\alpha _s^{(5)}(M_Z)=0.118`$. Let us further assume that it is determined from the experiment with three-loop accuracy, which means that in the $`\beta `$ function (2) only the coefficients up to $`\beta _2`$ are considered and $`\beta _3`$ is neglected. Let us now evaluate the strong coupling at the scale $`\mu =M_b`$ and compare the results obtained with the different strategies outlined above. In the following a possible Mathematica session is shown. NumDef is a set of Mathematica rules which assigns typical values to the physical parameters used in our procedures. The numbers used in this paper can be found in Eq. (39) and the procedures are described in the Appendix. ``` In[1]:= <<RunDec.m; ``` Comment: evaluation of $`\mathrm{\Lambda }`$ from $`\alpha _s^{(5)}(M_Z)`$ based on the explicit solution, Eq. (4), and subsequent evaluation of $`\alpha _s^{(5)}(M_b)`$ from $`\mathrm{\Lambda }`$ based on Eq. (5). ``` In[2]:= lamex = LamExpl[asMz/.NumDef,Mz/.NumDef,5,3] Out[2]= 0.208905 In[3]:= AlphasLam[lamex,Mb/.NumDef,5,3] Out[3]= 0.216610 ``` Comment: evaluation of $`\mathrm{\Lambda }`$ from $`\alpha _s^{(5)}(M_Z)`$ based on Eq. (5), and subsequent evaluation of $`\alpha _s^{(5)}(M_b)`$ from $`\mathrm{\Lambda }`$ based on Eq. (5). ``` In[4]:= lamim = LamImpl[asMz/.NumDef,Mz/.NumDef,5,3] Out[4]= 0.208348 In[5]:= AlphasLam[lamim,Mb/.NumDef,5,3] Out[5]= 0.216444 ``` Comment: evaluation of $`\alpha _s^{(5)}(M_b)`$ from $`\alpha _s^{(5)}(M_Z)`$ based on Eq. (1). ``` In[6]:= AlphasExact[asMz/.NumDef,Mz/.NumDef,Mb/.NumDef,5,3] Out[6]= 0.216712 ``` Rounding to three significant digits leads to a difference of $`\pm 1`$ in the last digit. Considering the direct integration of (1) as the most precise one we can conclude $`\alpha _s^{(5)}(M_b)=0.217`$ assuming three-loop accuracy. In Tab. 1 the influence of the number of loops is studied in the evaluation of $`\alpha _s^{(5)}`$ at the scale $`M_b`$ and the (hypothetical) scale 1 GeV using $`\alpha _s^{(5)}(M_Z)`$ as input. (For the latter only the function AlphasExact\[\] is used for the computation.) It can be seen that the inclusion of $`\beta _1`$ leads to a significant jump in $`\alpha _s^{(5)}(M_b)`$ whereas the effect of the three- and four-loop coefficients, i.e. $`\beta _2`$ and $`\beta _3`$, is only marginal. Their influence is more pronounced for $`\mu =1`$ GeV. ## 3 Quark masses in the $`\overline{\mathrm{MS}}`$ and on-shell scheme In the $`\overline{\mathrm{MS}}`$ scheme the running of the quark masses is governed by the function $`\gamma _m(\alpha _s)`$ $`\mu ^2{\displaystyle \frac{d}{d\mu ^2}}m^{(n_f)}(\mu )`$ $`=`$ $`m^{(n_f)}(\mu )\gamma _m^{(n_f)}\left(\alpha _s^{(n_f)}\right)=m^{(n_f)}(\mu ){\displaystyle \underset{i0}{}}\gamma _{m,i}^{(n_f)}\left({\displaystyle \frac{\alpha _s^{(n_f)}(\mu )}{\pi }}\right)^{i+1},`$ (6) where the coefficients $`\gamma _{m,i}`$ are known up to the four-loop order $`\gamma _{m,0}^{(n_f)}`$ $`=`$ $`1,`$ $`\gamma _{m,1}^{(n_f)}`$ $`=`$ $`{\displaystyle \frac{1}{16}}\left[{\displaystyle \frac{202}{3}}{\displaystyle \frac{20}{9}}n_f\right],`$ $`\gamma _{m,2}^{(n_f)}`$ $`=`$ $`{\displaystyle \frac{1}{64}}\left[1249+\left({\displaystyle \frac{2216}{27}}{\displaystyle \frac{160}{3}}\zeta _3\right)n_f{\displaystyle \frac{140}{81}}n_f^2\right],`$ $`\gamma _{m,3}^{(n_f)}`$ $`=`$ $`{\displaystyle \frac{1}{256}}[{\displaystyle \frac{4603055}{162}}+{\displaystyle \frac{135680}{27}}\zeta _38800\zeta _5+({\displaystyle \frac{91723}{27}}{\displaystyle \frac{34192}{9}}\zeta _3+880\zeta _4`$ (7) $`+{\displaystyle \frac{18400}{9}}\zeta _5)n_f+({\displaystyle \frac{5242}{243}}+{\displaystyle \frac{800}{9}}\zeta _3{\displaystyle \frac{160}{3}}\zeta _4)n_f^2`$ $`+({\displaystyle \frac{332}{243}}+{\displaystyle \frac{64}{27}}\zeta _3)n_f^3],`$ with $`\zeta _3\mathrm{1.202\hspace{0.17em}057}`$, $`\zeta _4=\pi ^4/90`$ and $`\zeta _5\mathrm{1.036\hspace{0.17em}928}`$. In analogy to (3) we define $`c_i^{(n_f)}`$ $`=`$ $`{\displaystyle \frac{\gamma _{m,i}^{(n_f)}}{\beta _0^{(n_f)}}}.`$ (8) Combining Eqs. (1) and (6) leads to a differential equation for $`m(\mu )`$ as a function of $`\alpha _s(\mu )`$. It has the solution $$\frac{m(\mu )}{m(\mu _0)}=\frac{c(\alpha _s(\mu )/\pi )}{c(\alpha _s(\mu _0)/\pi )},$$ (9) with $`c(x)`$ $`=`$ $`x^{c_0}\{1+(c_1b_1c_0)x+{\displaystyle \frac{1}{2}}[(c_1b_1c_0)^2+c_2b_1c_1+b_1^2c_0b_2c_0]x^2`$ (10) $`+[{\displaystyle \frac{1}{6}}(c_1b_1c_0)^3+{\displaystyle \frac{1}{2}}(c_1b_1c_0)(c_2b_1c_1+b_1^2c_0b_2c_0)`$ $`+{\displaystyle \frac{1}{3}}(c_3b_1c_2+b_1^2c_1b_2c_1b_1^3c_0+2b_1b_2c_0b_3c_0)]x^3\},`$ where terms of $`𝒪(x^4)`$ have been neglected. For a given mass, $`m`$, at scale $`\mu _0`$ and $`\alpha _s(\mu _0)`$ the scale invariant mass $`\mu _m=m(\mu _m)`$ can be obtained from Eq. (9) by iteration. Note the appearance of $`\alpha _s(\mu )`$ on the r.h.s. of (9). Thus for the computation of $`\mu _m`$ it is convenient to use in a first step $`\alpha _s(\mu _0)`$ in combination with Eq. (5) to determine $`\mathrm{\Lambda }`$. Afterwards Eq. (5) is used again for the calculation of $`\alpha _s(\mu )`$ which is inserted in (9) before the iteration. From Eq. (9) it appears natural to define the mass $`\widehat{m}`$ $``$ $`{\displaystyle \frac{m(\mu )}{c(\alpha _s(\mu )/\pi )}},`$ (11) which is often used in the context of lattice calculations. By construction the mass $`\widehat{m}`$ is scale independent. It is furthermore scheme independent (as far as mass-independent schemes are concerned). This can be seen by considering the r.h.s. of (11) in the limit $`\mu \mathrm{}`$ $`\widehat{m}`$ $`=`$ $`\underset{\mu \mathrm{}}{lim}m(\mu )\left({\displaystyle \frac{\alpha _s(\mu )}{\pi }}\right)^{\frac{\gamma _{m,0}}{\beta _0}},`$ (12) and by recalling the fact that the coefficients $`\beta _0`$ and $`\gamma _{m,0}`$ are scheme independent. In the following we will refer to $`\widehat{m}`$ as renormalization group invariant mass. In the following we want to provide the relations between the $`\overline{\mathrm{MS}}`$ and the on-shell mass. Whereas the coefficient of order $`\alpha _s^2`$ has been available since quite some time only recently the three-loop result could be obtained . In an asymptotic expansion in combination with conformal mapping and Padé approximation has been used in order to obtain a numerical result for the $`\overline{\mathrm{MS}}`$–on-shell conversion formula. The numerical results of are in perfect agreement with the subsequent analytical calculation of (cf. Tab. 2). For a given on-shell mass the $`\overline{\mathrm{MS}}`$ quantity can be computed with the help of $`{\displaystyle \frac{m(\mu )}{M}}`$ $`=`$ $`1+{\displaystyle \frac{\alpha _s^{(n_f)}(\mu )}{\pi }}[{\displaystyle \frac{4}{3}}l_{\mu M}]+\left({\displaystyle \frac{\alpha _s^{(n_f)}(\mu )}{\pi }}\right)^2[{\displaystyle \frac{3019}{288}}2\zeta _2{\displaystyle \frac{2}{3}}\zeta _2\mathrm{ln}2+{\displaystyle \frac{1}{6}}\zeta _3`$ (13) $`{\displaystyle \frac{445}{72}}l_{\mu M}{\displaystyle \frac{19}{24}}l_{\mu M}^2+({\displaystyle \frac{71}{144}}+{\displaystyle \frac{1}{3}}\zeta _2+{\displaystyle \frac{13}{36}}l_{\mu M}+{\displaystyle \frac{1}{12}}l_{\mu M}^2)n_l{\displaystyle \frac{4}{3}}{\displaystyle \underset{1in_l}{}}\mathrm{\Delta }\left({\displaystyle \frac{\text{M}_i}{M}}\right)]`$ $`+\left({\displaystyle \frac{\alpha _s^{(n_f)}(\mu )}{\pi }}\right)^3[z_m^{(3)}(M)+({\displaystyle \frac{165635}{2592}}{\displaystyle \frac{25}{3}}\zeta _2{\displaystyle \frac{25}{9}}\zeta _2\mathrm{ln}2+{\displaystyle \frac{55}{36}}\zeta _3)l_{\mu M}`$ $`{\displaystyle \frac{11779}{864}}l_{\mu M}^2{\displaystyle \frac{475}{432}}l_{\mu M}^3+n_l(({\displaystyle \frac{10051}{1296}}+{\displaystyle \frac{37}{18}}\zeta _2+{\displaystyle \frac{2}{9}}\zeta _2\mathrm{ln}2+{\displaystyle \frac{7}{9}}\zeta _3)l_{\mu M}+{\displaystyle \frac{911}{432}}l_{\mu M}^2`$ $`+{\displaystyle \frac{11}{54}}l_{\mu M}^3)+n_l^2(({\displaystyle \frac{89}{648}}{\displaystyle \frac{1}{9}}\zeta _2)l_{\mu M}^2{\displaystyle \frac{13}{216}}l_{\mu M}^2{\displaystyle \frac{1}{108}}l_{\mu M}^3)],`$ where $`\zeta _2=\pi ^2/6`$ and $`l_{\mu M}=\mathrm{ln}\mu ^2/M^2`$. $`n_l`$ is the number of light quarks. The function $`\mathrm{\Delta }(x)`$ arises from the two-loop diagram with a second fermion-loop . For $`0x1`$ it is approximated within an accuracy of 1% by $`\mathrm{\Delta }(x)`$ $`=`$ $`{\displaystyle \frac{\pi ^2}{8}}x0.597x^2+0.230x^3.`$ (14) The corresponding mass effects at order $`\alpha _s^3`$ are not yet known. In the argument of $`\mathrm{\Delta }(x)`$ the ratio of the on-shell mass of the light quarks, $`\text{M}_i`$, and the heavy one, $`M`$, appears. The coefficients $`z_m^{(3)}(M)`$ can be found in Tab. 2 for different values of $`n_l`$ where both the results of and are listed. For completeness also the corresponding two-loop coefficients (without the contribution from $`\mathrm{\Delta }(x)`$) are given. The analytical result for $`z_m^{(3)}(M)`$ reads $`z_m^{(3)}(M)`$ $`=`$ $`{\displaystyle \frac{9478333}{93312}}+{\displaystyle \frac{55}{162}}\mathrm{ln}^42+\left({\displaystyle \frac{644201}{6480}}+{\displaystyle \frac{587}{27}}\mathrm{ln}2+{\displaystyle \frac{44}{27}}\mathrm{ln}^22\right)\zeta _2{\displaystyle \frac{61}{27}}\zeta _3`$ (15) $`+{\displaystyle \frac{3475}{432}}\zeta _4+{\displaystyle \frac{1439}{72}}\zeta _2\zeta _3{\displaystyle \frac{1975}{216}}\zeta _5+{\displaystyle \frac{220}{27}}a_4+n_l[{\displaystyle \frac{246643}{23328}}{\displaystyle \frac{1}{81}}\mathrm{ln}^42`$ $`+({\displaystyle \frac{967}{108}}+{\displaystyle \frac{22}{27}}\mathrm{ln}2{\displaystyle \frac{4}{27}}\mathrm{ln}^22)\zeta _2+{\displaystyle \frac{241}{72}}\zeta _3{\displaystyle \frac{305}{108}}\zeta _4{\displaystyle \frac{8}{27}}a_4]`$ $`+n_l^2\left[{\displaystyle \frac{2353}{23328}}{\displaystyle \frac{13}{54}}\zeta _2{\displaystyle \frac{7}{54}}\zeta _3\right],`$ where $`a_4=\text{Li}_4(1/2)\mathrm{0.517\hspace{0.17em}479}`$. Iterating (13) leads to a relation between the scale-invariant mass, $`\mu _m=m(\mu _m)`$, and the on-shell mass $`{\displaystyle \frac{\mu _m}{M}}`$ $`=`$ $`1+{\displaystyle \frac{\alpha _s^{(n_f)}(M)}{\pi }}[{\displaystyle \frac{4}{3}}]+\left({\displaystyle \frac{\alpha _s^{(n_f)}(M)}{\pi }}\right)^2[{\displaystyle \frac{2251}{288}}2\zeta _2{\displaystyle \frac{2}{3}}\zeta _2\mathrm{ln}2+{\displaystyle \frac{1}{6}}\zeta _3`$ (16) $`+n_l({\displaystyle \frac{71}{144}}+{\displaystyle \frac{1}{3}}\zeta _2){\displaystyle \frac{4}{3}}{\displaystyle \underset{1in_l}{}}\mathrm{\Delta }\left({\displaystyle \frac{\text{M}_i}{M}}\right)]+({\displaystyle \frac{\alpha _s^{(n_f)}(M)}{\pi }})^3z_m^{SI,(3)}(M).`$ Inverting Eq. (13) leads to $`{\displaystyle \frac{M}{\mu _m}}`$ $`=`$ $`1+{\displaystyle \frac{\alpha _s^{(n_f)}(\mu _m)}{\pi }}{\displaystyle \frac{4}{3}}+\left({\displaystyle \frac{\alpha _s^{(n_f)}(\mu _m)}{\pi }}\right)^2[{\displaystyle \frac{307}{32}}+2\zeta _2+{\displaystyle \frac{2}{3}}\zeta _2\mathrm{ln}2{\displaystyle \frac{1}{6}}\zeta _3`$ $`+n_l({\displaystyle \frac{71}{144}}{\displaystyle \frac{1}{3}}\zeta _2)+{\displaystyle \frac{4}{3}}{\displaystyle \underset{1in_l}{}}\mathrm{\Delta }\left({\displaystyle \frac{m_i}{\mu _m}}\right)]+({\displaystyle \frac{\alpha _s^{(n_f)}(\mu _m)}{\pi }})^3z_m^{inv,(3)}(\mu _m),`$ where for convenience $`\mu ^2=m^2`$ has been chosen. The numerical values of the coefficients $`z_m^{SI}`$ and $`z_m^{inv}`$ can also be found in Tab. 2. Their analytic expressions are easily obtained from Eqs. (13) and (15). Eq. (LABEL:eq:zminv) can be used to compute the on-shell quark mass if the corresponding mass in the $`\overline{\mathrm{MS}}`$ scheme is provided. In order to avoid large logarithms it is suggestive to use in a first step the renormalization group equation (9) and evaluate $`\mu _m`$. In a second step Eq. (LABEL:eq:zminv) is used for $`\mu =\mu _m`$. Also in the case when the on-shell mass is given it is advantageous to use Eq. (LABEL:eq:zminv) for the computation of the $`\overline{\mathrm{MS}}`$ mass. The reason is that Eq. (13) contains contributions from the ill-defined pole mass of the light quarks like, e.g., the strange quark. In the case of the top quark it is safe to use (13) as in general the contributions for the charm and strange quark masses can be neglected. Concerning the determination of the quark masses a crucial role is played by lattice calculations. There it is not possible to use directly the $`\overline{\mathrm{MS}}`$ scheme as it is tightly connected to dimensional regularization. Rather one has to use a prescription which is based on the so-called momentum subtraction scheme. In general these schemes have the disadvantages that they are not mass independent. Recently, however, a mass definition based on momentum subtraction — the regularization invariant (RI) mass — has been proposed which enjoys this feature . In the relation to the $`\overline{\mathrm{MS}}`$ mass has been evaluated to three-loop accuracy. It reads: $`{\displaystyle \frac{m^{(n_f)}(\mu )}{m^{RI}(\mu )}}`$ $`=`$ $`1+{\displaystyle \frac{\alpha _s^{(n_f)}(\mu )}{\pi }}\left[{\displaystyle \frac{4}{3}}\right]+\left({\displaystyle \frac{\alpha _s^{(n_f)}(\mu )}{\pi }}\right)^2\left[{\displaystyle \frac{995}{72}}+{\displaystyle \frac{19}{6}}\zeta _3+{\displaystyle \frac{89}{144}}n_f\right]`$ (18) $`+\left({\displaystyle \frac{\alpha _s^{(n_f)}(\mu )}{\pi }}\right)^3[{\displaystyle \frac{6663911}{41472}}+{\displaystyle \frac{408007}{6912}}\zeta _3{\displaystyle \frac{185}{36}}\zeta _5+({\displaystyle \frac{118325}{7776}}+{\displaystyle \frac{5}{12}}\zeta _4`$ $`{\displaystyle \frac{617}{216}}\zeta _3)n_f+({\displaystyle \frac{4459}{23328}}{\displaystyle \frac{1}{54}}\zeta _3)n_f^2].`$ Let us at this point consider an explicit example. For a given mass $`\mu _b=m_b^{(5)}(\mu _b)=3.97`$ GeV five-flavour running (cf. Eqs. (9) and (10)) is used in order to obtain $`m_b^{(5)}(M_Z)`$ and the on-shell mass is computed with the help of Eq. (LABEL:eq:zminv). Furthermore the value for $`m^{RI}(M_Z)`$ is evaluated. In Tab. 3 the results are listed for different number of loops. At the end of this section we want to summarize the different mass definitions introduced in this section in the following table: \| $`M`$ \| on-shell mass \| \| --- \| --- \| \| $`m(\mu )`$ \| $`\overline{\mathrm{MS}}`$ mass \| \| $`\mu _m`$ \| scale invariant mass \| \| $`\widehat{m}`$ \| renormalization group invariant mass \| \| $`m^{RI}`$ \| regularization invariant mass \| ## 4 Decoupling at flavour thresholds In MS-like renormalization schemes, the Appelquist-Carazzone decoupling theorem does not in general apply to quantities that do not represent physical observables, such as beta functions or coupling constants, i.e., quarks with masses much larger than the considered energy scale do not automatically decouple. The standard procedure to circumvent this problem is to render decoupling explicit by using the language of effective field theory. The formulae presented below are valid for QCD with $`n_l=n_f1`$ massless quark flavours and one heavy flavour $`h`$, with mass $`m_h`$ which is supposed to be much larger than the energy scale. Then, one constructs an effective $`n_l`$-flavour theory by requiring consistency with the full $`n_f`$-flavour theory at an energy scale comparable to $`m_h`$, the heavy-quark threshold $`\mu ^{(n_f)}=𝒪(m_h)`$. This leads to a nontrivial matching condition between the couplings and light masses, $`m_q`$, of the two theories. Although, $`\alpha _s^{(n_l)}(m_h)=\alpha _s^{(n_f)}(m_h)`$ and $`m_q^{(n_l)}(m_h)=m_q^{(n_f)}(m_h)`$ at leading and next-to-leading order, this relation does not generally hold at higher orders in the $`\overline{\mathrm{MS}}`$ scheme. At $`𝒪(\alpha _s^2)`$ the corresponding correction terms have been computed in . The connection between the strong coupling constant in the effective and the full theory is given by $`\alpha _s^{(nf1)}(\mu )`$ $`=`$ $`\zeta _g^2\alpha _s^{(n_f)}(\mu ),`$ (19) where $`\zeta _g`$ is known up to the three-loop order : $`\left(\zeta _g^{MS}\right)^2`$ $`=`$ $`1+{\displaystyle \frac{\alpha _s^{(n_f)}(\mu )}{\pi }}\left({\displaystyle \frac{1}{6}}\mathrm{ln}{\displaystyle \frac{\mu ^2}{m_h^2}}\right)+\left({\displaystyle \frac{\alpha _s^{(n_f)}(\mu )}{\pi }}\right)^2\left({\displaystyle \frac{11}{72}}{\displaystyle \frac{11}{24}}\mathrm{ln}{\displaystyle \frac{\mu ^2}{m_h^2}}+{\displaystyle \frac{1}{36}}\mathrm{ln}^2{\displaystyle \frac{\mu ^2}{m_h^2}}\right)`$ (20) $`+\left({\displaystyle \frac{\alpha _s^{(n_f)}(\mu )}{\pi }}\right)^3[{\displaystyle \frac{564731}{124416}}{\displaystyle \frac{82043}{27648}}\zeta _3{\displaystyle \frac{955}{576}}\mathrm{ln}{\displaystyle \frac{\mu ^2}{m_h^2}}+{\displaystyle \frac{53}{576}}\mathrm{ln}^2{\displaystyle \frac{\mu ^2}{m_h^2}}`$ $`{\displaystyle \frac{1}{216}}\mathrm{ln}^3{\displaystyle \frac{\mu ^2}{m_h^2}}+n_l({\displaystyle \frac{2633}{31104}}+{\displaystyle \frac{67}{576}}\mathrm{ln}{\displaystyle \frac{\mu ^2}{m_h^2}}{\displaystyle \frac{1}{36}}\mathrm{ln}^2{\displaystyle \frac{\mu ^2}{m_h^2}})].`$ In this equation the $`\overline{\mathrm{MS}}`$ mass $`m_h(\mu )`$ — indicated by the superscript MS — is chosen for the parameterization of the heavy quark mass and $`\mu `$ represents the renormalization scale. Often it is convenient to express $`\zeta _g`$ through the scale invariant mass, denoted by $`\mu _h=m_h(\mu _h)`$: $`\left(\zeta _g^{SI}\right)^2`$ $`=`$ $`1+{\displaystyle \frac{\alpha _s^{(n_f)}(\mu )}{\pi }}\left({\displaystyle \frac{1}{6}}\mathrm{ln}{\displaystyle \frac{\mu ^2}{\mu _h^2}}\right)+\left({\displaystyle \frac{\alpha _s^{(n_f)}(\mu )}{\pi }}\right)^2\left({\displaystyle \frac{1}{36}}\mathrm{ln}^2{\displaystyle \frac{\mu ^2}{\mu _h^2}}{\displaystyle \frac{19}{24}}\mathrm{ln}{\displaystyle \frac{\mu ^2}{\mu _h^2}}+{\displaystyle \frac{11}{72}}\right)`$ (21) $`+\left({\displaystyle \frac{\alpha _s^{(n_f)}(\mu )}{\pi }}\right)^3[{\displaystyle \frac{1}{216}}\mathrm{ln}^3{\displaystyle \frac{\mu ^2}{\mu _h^2}}{\displaystyle \frac{131}{576}}\mathrm{ln}^2{\displaystyle \frac{\mu ^2}{\mu _h^2}}+{\displaystyle \frac{1}{1728}}\mathrm{ln}{\displaystyle \frac{\mu ^2}{\mu _h^2}}(6793+281n_l)`$ $`{\displaystyle \frac{82043}{27648}}\zeta _3+{\displaystyle \frac{564731}{124416}}{\displaystyle \frac{2633}{31104}}n_l].`$ Transforming the heavy quark mass into the on-shell scheme leads to $`\left(\zeta _g^{OS}\right)^2`$ $`=`$ $`1+{\displaystyle \frac{\alpha _s^{(n_f)}(\mu )}{\pi }}\left({\displaystyle \frac{1}{6}}\mathrm{ln}{\displaystyle \frac{\mu ^2}{M_h^2}}\right)+\left({\displaystyle \frac{\alpha _s^{(n_f)}(\mu )}{\pi }}\right)^2\left({\displaystyle \frac{7}{24}}{\displaystyle \frac{19}{24}}\mathrm{ln}{\displaystyle \frac{\mu ^2}{M_h^2}}+{\displaystyle \frac{1}{36}}\mathrm{ln}^2{\displaystyle \frac{\mu ^2}{M_h^2}}\right)`$ (22) $`+\left({\displaystyle \frac{\alpha _s^{(n_f)}(\mu )}{\pi }}\right)^3[{\displaystyle \frac{58933}{124416}}{\displaystyle \frac{2}{3}}\zeta _2(1+{\displaystyle \frac{1}{3}}\mathrm{ln}2){\displaystyle \frac{80507}{27648}}\zeta _3{\displaystyle \frac{8521}{1728}}\mathrm{ln}{\displaystyle \frac{\mu ^2}{M_h^2}}`$ $`{\displaystyle \frac{131}{576}}\mathrm{ln}^2{\displaystyle \frac{\mu ^2}{M_h^2}}{\displaystyle \frac{1}{216}}\mathrm{ln}^3{\displaystyle \frac{\mu ^2}{M_h^2}}+n_l({\displaystyle \frac{2479}{31104}}+{\displaystyle \frac{\zeta _2}{9}}+{\displaystyle \frac{409}{1728}}\mathrm{ln}{\displaystyle \frac{\mu ^2}{M_h^2}})].`$ In practical applications also the inverted formulae are needed which read for Eqs. (20), (21) and (22): $`{\displaystyle \frac{1}{\left(\zeta _g^{MS}\right)^2}}`$ $`=`$ $`1+{\displaystyle \frac{\alpha _s^{(n_l)}(\mu )}{\pi }}\left({\displaystyle \frac{1}{6}}\mathrm{ln}{\displaystyle \frac{\mu ^2}{m_h^2}}\right)+\left({\displaystyle \frac{\alpha _s^{(n_l)}(\mu )}{\pi }}\right)^2\left({\displaystyle \frac{11}{72}}+{\displaystyle \frac{11}{24}}\mathrm{ln}{\displaystyle \frac{\mu ^2}{m_h^2}}+{\displaystyle \frac{1}{36}}\mathrm{ln}^2{\displaystyle \frac{\mu ^2}{m_h^2}}\right)`$ (23) $`+\left({\displaystyle \frac{\alpha _s^{(n_l)}(\mu )}{\pi }}\right)^3[{\displaystyle \frac{564731}{124416}}+{\displaystyle \frac{82043}{27648}}\zeta _3+{\displaystyle \frac{2645}{1728}}\mathrm{ln}{\displaystyle \frac{\mu ^2}{m_h^2}}+{\displaystyle \frac{167}{576}}\mathrm{ln}^2{\displaystyle \frac{\mu ^2}{m_h^2}}`$ $`+{\displaystyle \frac{1}{216}}\mathrm{ln}^3{\displaystyle \frac{\mu ^2}{m_h^2}}+n_l({\displaystyle \frac{2633}{31104}}{\displaystyle \frac{67}{576}}\mathrm{ln}{\displaystyle \frac{\mu ^2}{m_h^2}}+{\displaystyle \frac{1}{36}}\mathrm{ln}^2{\displaystyle \frac{\mu ^2}{m_h^2}})],`$ $`{\displaystyle \frac{1}{\left(\zeta _g^{SI}\right)^2}}`$ $`=`$ $`1+{\displaystyle \frac{\alpha _s^{(n_l)}(\mu )}{\pi }}\left({\displaystyle \frac{1}{6}}\mathrm{ln}{\displaystyle \frac{\mu ^2}{\mu _h^2}}\right)+\left({\displaystyle \frac{\alpha _s^{(n_l)}(\mu )}{\pi }}\right)^2\left({\displaystyle \frac{11}{72}}+{\displaystyle \frac{19}{24}}\mathrm{ln}{\displaystyle \frac{\mu ^2}{\mu _h^2}}+{\displaystyle \frac{1}{36}}\mathrm{ln}^2{\displaystyle \frac{\mu ^2}{\mu _h^2}}\right)`$ (24) $`+\left({\displaystyle \frac{\alpha _s^{(n_l)}(\mu )}{\pi }}\right)^3[{\displaystyle \frac{564731}{124416}}+{\displaystyle \frac{82043}{27648}}\zeta _3+{\displaystyle \frac{2191}{576}}\mathrm{ln}{\displaystyle \frac{\mu ^2}{\mu _h^2}}+{\displaystyle \frac{511}{576}}\mathrm{ln}^2{\displaystyle \frac{\mu ^2}{\mu _h^2}}`$ $`+{\displaystyle \frac{1}{216}}\mathrm{ln}^3{\displaystyle \frac{\mu ^2}{\mu _h^2}}+n_l({\displaystyle \frac{2633}{31104}}{\displaystyle \frac{281}{1728}}\mathrm{ln}{\displaystyle \frac{\mu ^2}{\mu _h^2}})],`$ $`{\displaystyle \frac{1}{\left(\zeta _g^{OS}\right)^2}}`$ $`=`$ $`1+{\displaystyle \frac{\alpha _s^{(n_l)}(\mu )}{\pi }}\left({\displaystyle \frac{1}{6}}\mathrm{ln}{\displaystyle \frac{\mu ^2}{M_h^2}}\right)+\left({\displaystyle \frac{\alpha _s^{(n_l)}(\mu )}{\pi }}\right)^2\left({\displaystyle \frac{7}{24}}+{\displaystyle \frac{19}{24}}\mathrm{ln}{\displaystyle \frac{\mu ^2}{M_h^2}}+{\displaystyle \frac{1}{36}}\mathrm{ln}^2{\displaystyle \frac{\mu ^2}{M_h^2}}\right)`$ (25) $`+\left({\displaystyle \frac{\alpha _s^{(n_l)}(\mu )}{\pi }}\right)^3[{\displaystyle \frac{58933}{124416}}+{\displaystyle \frac{2}{3}}\zeta _2(1+{\displaystyle \frac{1}{3}}\mathrm{ln}2)+{\displaystyle \frac{80507}{27648}}\zeta _3+{\displaystyle \frac{8941}{1728}}\mathrm{ln}{\displaystyle \frac{\mu ^2}{M_h^2}}`$ $`+{\displaystyle \frac{511}{576}}\mathrm{ln}^2{\displaystyle \frac{\mu ^2}{M_h^2}}+{\displaystyle \frac{1}{216}}\mathrm{ln}^3{\displaystyle \frac{\mu ^2}{M_h^2}}+n_l({\displaystyle \frac{2479}{31104}}{\displaystyle \frac{\zeta _2}{9}}{\displaystyle \frac{409}{1728}}\mathrm{ln}{\displaystyle \frac{\mu ^2}{M_h^2}})].`$ The decoupling relations (20)–(25) have to be applied whenever a flavour threshold is to be crossed. At this point we briefly want to comment on the order of $`\alpha _s`$ which has to be used for the running, respectively, the decoupling if the analysis should be consistent. If the $`\mu `$ evolution of $`\alpha _s^{(n_f)}(\mu )`$ is to be performed at $`N+1`$ loops, i.e., with the highest coefficient in Eq. (1) being $`\beta _N^{(n_f)}`$, then consistency requires the matching conditions to be implemented in terms of $`N`$-loop formulae. Then, the residual $`\mu `$ dependence of physical observables will be of order $`N+2`$. As an example let us compute $`\alpha _s^{(4)}(M_c)`$ from $`\alpha _s^{(5)}(M_Z)=0.118`$. Let us furthermore consider the on-shell definition of the heavy quark, $`M_b`$, i.e. we use Eq. (22) for the analysis. For the scale $`\mu `$ in (19) where the matching is performed we choose $`\mu _{th}=M_b`$. Assuming four-loop accuracy for the beta function and (as a consequence) three-loop accuracy for the matching computation one would proceed as follows (see Appendix for a description of the procedures): ``` In[2]:= (alsmuth = AlphasExact[asMz/.NumDef,Mz/.NumDef,Mb/.NumDef,5,4]) Out[2]= 0.2169467 In[3]:= (alsmuthp = DecAsDownOS[alsmuth,Mb/.NumDef,Mb/.NumDef,4,4]) Out[3]= 0.2163396 In[4]:= (alsMc = AlphasExact[alsmuthp,Mb/.NumDef,Mc/.NumDef,4,4]) Out[4]= 0.337848 ``` Finally one arrives at $`\alpha _s^{(4)}(M_c)=0.338`$. These steps are summarized in the function AsRunDec\[\] where the corresponding call would read ``` In[5]:= AsRunDec[asMz/.NumDef,Mz/.NumDef,Mc/.NumDef,4] Out[5]= 0.337848 ``` Note that the loop-argument of the function DecAsDownOS\[\] (last argument) refers to the order used for the running, i.e. in the considered case the “4” means that the three-loop relation is used for the decoupling. In this example the effect of the decoupling is quite small. It is actually comparable to the uncertainty from using different methods for the running (cf. Tab. 1). However, one has to remember that for the matching scale the heavy quark mass itself has been used, whence all logarithms in Eq. (22) vanish. A different choice would lead to a different result for $`\alpha _s^{(4)}(M_c)`$. On the other hand, on general grounds, the decoupling procedure should not depend on the choice of that scale, respectively, the dependence should become weaker when going to higher orders. In Tab. 4 the dependence of $`\alpha _s^{(4)}(M_c)`$ on the number of loops is shown. For the matching scale $`M_Z`$, $`M_b`$ and $`1`$ GeV has been chosen. It can be clearly seen that the four-loop analysis provides the most stable values — even in the case when the matching is performed at the a high scale like the $`Z`$ boson mass. This is expected on general grounds as physical results should not depend on the matching scale. We should mention that in case the $`\overline{\mathrm{MS}}`$ definition for the heavy quark is used in a first step $`m_h(\mu _{th})`$ has to be evaluated. The corresponding formulae can be found in Section 3. They are also implemented as Mathematica procedures and described in the Appendix. In Fig. 1 it is demonstrated that the inclusion of the four-loop coefficient $`\beta _3`$ accompanied by the three-loop matching leads to an independence of $`\mu _{th}=\mu ^{(5)}`$ over a very broad range . The plot shows the dependence of $`\alpha _s^{(5)}(M_Z)`$ on the matching scale (denoted by $`\mu ^{(5)}`$) where $`\alpha _s^{(4)}(M_\tau )`$ is used as starting point. Our procedure to get the different curves is as follows. We first calculate $`\alpha _s^{(4)}(\mu ^{(5)})`$ by exactly integrating Eq. (1) with the initial condition $`\alpha _s^{(4)}(M_\tau )=0.36`$, then obtain $`\alpha _s^{(5)}(\mu ^{(5)})`$ from Eqs. (25) with $`M_b=4.7`$ GeV, and finally compute $`\alpha _s^{(5)}(M_Z)`$ with Eq. (1). For consistency, $`N`$-loop evolution must be accompanied by $`(N1)`$-loop matching, i.e. if we omit terms of $`𝒪(\alpha _s^{N+2})`$ on the right-hand side of Eq. (1), we need to discard those of $`𝒪(\alpha _s^N)`$ in Eq. (25) at the same time. In Fig. 1, the variation of $`\alpha _s^{(5)}(M_Z)`$ with $`\mu ^{(5)}/M_b`$ is displayed for the various levels of accuracy, ranging from one-loop to four-loop evolution. For illustration, $`\mu ^{(5)}`$ is varied by almost two orders of magnitude. While the leading-order result exhibits a strong logarithmic behaviour, it stabilizes as we go to higher orders. The four-loop curve is almost flat for $`\mu ^{(5)}\mathrm{\Gamma }>\mathrm{\hspace{0.33em}1}`$ GeV. Besides the $`\mu ^{(5)}`$ dependence of $`\alpha _s^{(5)}(M_Z)`$, also its absolute normalization is significantly affected by the higher orders. At the central matching scale $`\mu ^{(5)}=M_b`$, we encounter a rapid, monotonic convergence behaviour. Fig. 1 can immediately be reproduced with the help of the procedure AlL2AlH\[\] described in the Appendix. Up to now only the decoupling of the coupling constant has been considered. However, also the (relatively) lighter quark masses undergo a decoupling procedure when crossing a flavour threshold. If we define the connection between the quark mass in the effective and full theory through $`m_q^{(n_f1)}`$ $`=`$ $`\zeta _mm_q^{(n_f)},`$ (26) the decoupling constant $`\zeta _m`$ is given by $`\zeta _m^{MS}`$ $`=`$ $`1+\left({\displaystyle \frac{\alpha _s^{(n_f)}(\mu )}{\pi }}\right)^2({\displaystyle \frac{89}{432}}{\displaystyle \frac{5}{36}}\mathrm{ln}{\displaystyle \frac{\mu ^2}{m_h^2}}+{\displaystyle \frac{1}{12}}\mathrm{ln}^2{\displaystyle \frac{\mu ^2}{m_h^2}})+\left({\displaystyle \frac{\alpha _s^{(n_f)}(\mu )}{\pi }}\right)^3[{\displaystyle \frac{2951}{2916}}`$ (27) $`{\displaystyle \frac{407}{864}}\zeta _3+{\displaystyle \frac{5}{4}}\zeta _4{\displaystyle \frac{1}{36}}B_4+\left({\displaystyle \frac{311}{2592}}{\displaystyle \frac{5}{6}}\zeta _3\right)\mathrm{ln}{\displaystyle \frac{\mu ^2}{m_h^2}}+{\displaystyle \frac{175}{432}}\mathrm{ln}^2{\displaystyle \frac{\mu ^2}{m_h^2}}`$ $`+{\displaystyle \frac{29}{216}}\mathrm{ln}^3{\displaystyle \frac{\mu ^2}{m_h^2}}+n_l({\displaystyle \frac{1327}{11664}}{\displaystyle \frac{2}{27}}\zeta _3{\displaystyle \frac{53}{432}}\mathrm{ln}{\displaystyle \frac{\mu ^2}{m_h^2}}{\displaystyle \frac{1}{108}}\mathrm{ln}^3{\displaystyle \frac{\mu ^2}{m_h^2}})],`$ where $`B_4`$ $`=`$ $`16\text{Li}_4\left({\displaystyle \frac{1}{2}}\right){\displaystyle \frac{13}{2}}\zeta _44\zeta _2\mathrm{ln}^22+{\displaystyle \frac{2}{3}}\mathrm{ln}^42`$ (28) $``$ $`\mathrm{1.762\hspace{0.17em}800}.`$ Note that all three quantities in Eq. (26) depend on the renormalization scale $`\mu `$. Again it turns out to be useful to consider in addition to (27) the quantities where the scale invariant and the on-shell mass, respectively, has been used for the parameterization of the heavy quark: $`\zeta _m^{SI}`$ $`=`$ $`1+\left({\displaystyle \frac{\alpha _s^{(n_f)}(\mu )}{\pi }}\right)^2({\displaystyle \frac{89}{432}}{\displaystyle \frac{5}{36}}\mathrm{ln}{\displaystyle \frac{\mu ^2}{\mu _h^2}}+{\displaystyle \frac{1}{12}}\mathrm{ln}^2{\displaystyle \frac{\mu ^2}{\mu _h^2}})+\left({\displaystyle \frac{\alpha _s^{(n_f)}(\mu )}{\pi }}\right)^3[{\displaystyle \frac{2951}{2916}}`$ (29) $`{\displaystyle \frac{407}{864}}\zeta _3+{\displaystyle \frac{5}{4}}\zeta _4{\displaystyle \frac{1}{36}}B_4+\left({\displaystyle \frac{1031}{2592}}{\displaystyle \frac{5}{6}}\zeta _3\right)\mathrm{ln}{\displaystyle \frac{\mu ^2}{\mu _h^2}}+{\displaystyle \frac{319}{432}}\mathrm{ln}^2{\displaystyle \frac{\mu ^2}{\mu _h^2}}`$ $`+{\displaystyle \frac{29}{216}}\mathrm{ln}^3{\displaystyle \frac{\mu ^2}{\mu _h^2}}+n_l({\displaystyle \frac{1327}{11664}}{\displaystyle \frac{2}{27}}\zeta _3{\displaystyle \frac{53}{432}}\mathrm{ln}{\displaystyle \frac{\mu ^2}{\mu _h^2}}{\displaystyle \frac{1}{108}}\mathrm{ln}^3{\displaystyle \frac{\mu ^2}{\mu _h^2}})],`$ $`\zeta _m^{OS}`$ $`=`$ $`1+\left({\displaystyle \frac{\alpha _s^{(n_f)}(\mu )}{\pi }}\right)^2({\displaystyle \frac{89}{432}}{\displaystyle \frac{5}{36}}\mathrm{ln}{\displaystyle \frac{\mu ^2}{M_h^2}}+{\displaystyle \frac{1}{12}}\mathrm{ln}^2{\displaystyle \frac{\mu ^2}{M_h^2}})+\left({\displaystyle \frac{\alpha _s^{(n_f)}(\mu )}{\pi }}\right)^3[{\displaystyle \frac{1871}{2916}}`$ (30) $`{\displaystyle \frac{407}{864}}\zeta _3+{\displaystyle \frac{5}{4}}\zeta _4{\displaystyle \frac{1}{36}}B_4+\left({\displaystyle \frac{121}{2592}}{\displaystyle \frac{5}{6}}\zeta _3\right)\mathrm{ln}{\displaystyle \frac{\mu ^2}{M_h^2}}+{\displaystyle \frac{319}{432}}\mathrm{ln}^2{\displaystyle \frac{\mu ^2}{M_h^2}}`$ $`+{\displaystyle \frac{29}{216}}\mathrm{ln}^3{\displaystyle \frac{\mu ^2}{M_h^2}}+n_l({\displaystyle \frac{1327}{11664}}{\displaystyle \frac{2}{27}}\zeta _3{\displaystyle \frac{53}{432}}\mathrm{ln}{\displaystyle \frac{\mu ^2}{M_h^2}}{\displaystyle \frac{1}{108}}\mathrm{ln}^3{\displaystyle \frac{\mu ^2}{M_h^2}})].`$ The corresponding inverted relations read $`{\displaystyle \frac{1}{\zeta _m^{MS}}}`$ $`=`$ $`1+\left({\displaystyle \frac{\alpha _s^{(n_l)}(\mu )}{\pi }}\right)^2({\displaystyle \frac{89}{432}}+{\displaystyle \frac{5}{36}}\mathrm{ln}{\displaystyle \frac{\mu ^2}{m_h^2}}{\displaystyle \frac{1}{12}}\mathrm{ln}^2{\displaystyle \frac{\mu ^2}{m_h^2}})+\left({\displaystyle \frac{\alpha _s^{(n_l)}(\mu )}{\pi }}\right)^3[{\displaystyle \frac{2951}{2916}}`$ (31) $`+{\displaystyle \frac{407}{864}}\zeta _3{\displaystyle \frac{5}{4}}\zeta _4+{\displaystyle \frac{1}{36}}B_4+\left({\displaystyle \frac{133}{2592}}+{\displaystyle \frac{5}{6}}\zeta _3\right)\mathrm{ln}{\displaystyle \frac{\mu ^2}{m_h^2}}{\displaystyle \frac{155}{432}}\mathrm{ln}^2{\displaystyle \frac{\mu ^2}{m_h^2}}`$ $`{\displaystyle \frac{35}{216}}\mathrm{ln}^3{\displaystyle \frac{\mu ^2}{m_h^2}}+n_l({\displaystyle \frac{1327}{11664}}+{\displaystyle \frac{2}{27}}\zeta _3+{\displaystyle \frac{53}{432}}\mathrm{ln}{\displaystyle \frac{\mu ^2}{m_h^2}}+{\displaystyle \frac{1}{108}}\mathrm{ln}^3{\displaystyle \frac{\mu ^2}{m_h^2}})],`$ $`{\displaystyle \frac{1}{\zeta _m^{SI}}}`$ $`=`$ $`1+\left({\displaystyle \frac{\alpha _s^{(n_l)}(\mu )}{\pi }}\right)^2({\displaystyle \frac{89}{432}}+{\displaystyle \frac{5}{36}}\mathrm{ln}{\displaystyle \frac{\mu ^2}{\mu _h^2}}{\displaystyle \frac{1}{12}}\mathrm{ln}^2{\displaystyle \frac{\mu ^2}{\mu _h^2}})+\left({\displaystyle \frac{\alpha _s^{(n_l)}(\mu )}{\pi }}\right)^3[{\displaystyle \frac{2951}{2916}}`$ (32) $`+{\displaystyle \frac{407}{864}}\zeta _3{\displaystyle \frac{5}{4}}\zeta _4+{\displaystyle \frac{1}{36}}B_4+\left({\displaystyle \frac{853}{2592}}+{\displaystyle \frac{5}{6}}\zeta _3\right)\mathrm{ln}{\displaystyle \frac{\mu ^2}{\mu _h^2}}{\displaystyle \frac{299}{432}}\mathrm{ln}^2{\displaystyle \frac{\mu ^2}{\mu _h^2}}`$ $`{\displaystyle \frac{35}{216}}\mathrm{ln}^3{\displaystyle \frac{\mu ^2}{\mu _h^2}}+n_l({\displaystyle \frac{1327}{11664}}+{\displaystyle \frac{2}{27}}\zeta _3+{\displaystyle \frac{53}{432}}\mathrm{ln}{\displaystyle \frac{\mu ^2}{\mu _h^2}}+{\displaystyle \frac{1}{108}}\mathrm{ln}^3{\displaystyle \frac{\mu ^2}{\mu _h^2}})],`$ $`{\displaystyle \frac{1}{\zeta _m^{OS}}}`$ $`=`$ $`1+\left({\displaystyle \frac{\alpha _s^{(n_l)}(\mu )}{\pi }}\right)^2({\displaystyle \frac{89}{432}}+{\displaystyle \frac{5}{36}}\mathrm{ln}{\displaystyle \frac{\mu ^2}{M_h^2}}{\displaystyle \frac{1}{12}}\mathrm{ln}^2{\displaystyle \frac{\mu ^2}{M_h^2}})+\left({\displaystyle \frac{\alpha _s^{(n_l)}(\mu )}{\pi }}\right)^3[{\displaystyle \frac{1871}{2916}}`$ (33) $`+{\displaystyle \frac{407}{864}}\zeta _3{\displaystyle \frac{5}{4}}\zeta _4+{\displaystyle \frac{1}{36}}B_4+\left({\displaystyle \frac{299}{2592}}+{\displaystyle \frac{5}{6}}\zeta _3\right)\mathrm{ln}{\displaystyle \frac{\mu ^2}{M_h^2}}{\displaystyle \frac{299}{432}}\mathrm{ln}^2{\displaystyle \frac{\mu ^2}{M_h^2}}`$ $`{\displaystyle \frac{35}{216}}\mathrm{ln}^3{\displaystyle \frac{\mu ^2}{M_h^2}}+n_l({\displaystyle \frac{1327}{11664}}+{\displaystyle \frac{2}{27}}\zeta _3+{\displaystyle \frac{53}{432}}\mathrm{ln}{\displaystyle \frac{\mu ^2}{M_h^2}}+{\displaystyle \frac{1}{108}}\mathrm{ln}^3{\displaystyle \frac{\mu ^2}{M_h^2}})].`$ As an example we compute $`m_c^{(5)}(M_Z)`$ for different number of loops and different matching points $`\mu _{th}`$. The results can be found in Tab. 5 where $`\mu _{th}=M_Z,M_b`$, and 1 GeV has been chosen. It can clearly be seen that the four-loop result provides the most stable values for $`m_c^{(5)}(M_Z)`$. A similar analysis as in Fig. 1 may be performed for the light-quark masses as well. For illustration, let us investigate how the $`\mu ^{(5)}`$ dependence of the relation between $`\mu _c=m_c^{(4)}(\mu _c)`$ and $`m_c^{(5)}(M_Z)`$ changes under the inclusion of higher orders in evolution and matching. As typical input parameters, we choose $`\mu _c=1.2`$ GeV, $`M_b=4.7`$ GeV, and $`\alpha _s^{(5)}(M_Z)=0.118`$. We first evolve $`m_c^{(4)}(\mu )`$ from $`\mu =\mu _c`$ to $`\mu =\mu _{th}=\mu ^{(5)}`$ via Eq. (10), then obtain $`m_c^{(5)}(\mu ^{(5)})`$ from Eqs. (33), and finally evolve $`m_c^{(5)}(\mu )`$ from $`\mu =\mu ^{(5)}`$ to $`\mu =M_Z`$ via Eq. (10). In all steps, $`\alpha _s^{(n_f)}(\mu )`$ is evaluated with the same values of $`n_f`$ and $`\mu `$ as $`m_c^{(n_f)}(\mu )`$. In Fig. 2, we show the resulting values of $`m_c^{(5)}(M_Z)`$ corresponding to $`N`$-loop evolution with $`(N1)`$-loop matching for $`N=1,\mathrm{},4`$. Similarly to Fig. 1, we observe a rapid, monotonic convergence behaviour at the central matching scale $`\mu ^{(5)}=M_b`$. Again, the prediction for $`N=4`$ is remarkably stable under the variation of $`\mu ^{(5)}`$ as long as $`\mu ^{(5)}\mathrm{\Gamma }>\mathrm{\hspace{0.33em}1}`$ GeV. Fig. 2 can easily be reproduced with the help of the procedure mL2mH\[\] (see Appendix). If one chooses to perform the running of $`\alpha _s(\mu )`$ with the help of $`\mathrm{\Lambda }`$ it is useful to have an equation at hand which relates this parameter in the full and effective theory. Combining Eqs. (4), (5) and (21) one obtains $`\beta _0^{}\mathrm{ln}{\displaystyle \frac{\mathrm{\Lambda }^2}{\mathrm{\Lambda }^2}}`$ $`=`$ $`(\beta _0^{}\beta _0)l_h+(b_1^{}b_1)\mathrm{ln}l_hb_1^{}\mathrm{ln}{\displaystyle \frac{\beta _0^{}}{\beta _0}}`$ (34) $`+{\displaystyle \frac{1}{\beta _0l_h}}\left[b_1(b_1^{}b_1)\mathrm{ln}l_h+b_1^2b_1^2b_2^{}+b_2+{\displaystyle \frac{11}{72}}\right]`$ $`+{\displaystyle \frac{1}{(\beta _0l_h)^2}}\{{\displaystyle \frac{b_1^2}{2}}(b_1^{}b_1)\mathrm{ln}^2l_h+b_1[b_1^{}(b_1^{}b_1)`$ $`+b_2^{}b_2{\displaystyle \frac{11}{72}}]\mathrm{ln}l_h+{\displaystyle \frac{1}{2}}(b_1^3b_1^3b_3^{}+b_3)`$ $`+b_1^{}(b_1^2+b_2^{}b_2{\displaystyle \frac{11}{72}})+{\displaystyle \frac{564731}{124416}}{\displaystyle \frac{82043}{27648}}\zeta _3{\displaystyle \frac{2633}{31104}}n_l\},`$ where $`l_h=\mathrm{ln}(\mu _h^2/\mathrm{\Lambda }^2)`$ and the primed quantities refer to the $`(n_f1)`$-flavour effective theory. In this equation $`\mu _h`$ has been chosen for the matching scale which is particularly convenient, since it eliminates the renormalization group logarithms in (21). This choice is furthermore justified with the help of Figs. 1 and 2 where it can be seen that, in higher orders, the actual value of the matching scale does not matter as long as it is comparable to the heavy-quark mass. In Eq. (34) the four different powers in $`l_h`$ correspond to the different loop orders. Whereas at one-loop accuracy only the linear term in $`l_h`$ has to be taken into account at four-loop order also the $`1/l_h^2`$ contribution has to be considered. Eq. (34) is implemented in the procedure DecLambdaDown\[\]. For completeness we also display the inverted relation of Eq. (34): $`\beta _0\mathrm{ln}{\displaystyle \frac{\mathrm{\Lambda }^2}{\mathrm{\Lambda }^2}}`$ $`=`$ $`(\beta _0\beta _0^{})l_h^{}+(b_1b_1^{})\mathrm{ln}l_h^{}b_1\mathrm{ln}{\displaystyle \frac{\beta _0}{\beta _0^{}}}`$ (35) $`+{\displaystyle \frac{1}{\beta _0^{}l_h^{}}}\left[b_1^{}(b_1b_1^{})\mathrm{ln}l_h^{}+b_1^2b_1^2b_2+b_2^{}{\displaystyle \frac{11}{72}}\right]`$ $`+{\displaystyle \frac{1}{(\beta _0^{}l_h^{})^2}}\{{\displaystyle \frac{b_1^2}{2}}(b_1b_1^{})\mathrm{ln}^2l_h^{}+b_1^{}[b_1(b_1b_1^{})`$ $`+b_2b_2^{}+{\displaystyle \frac{11}{72}}]\mathrm{ln}l_h^{}+{\displaystyle \frac{1}{2}}(b_1^3b_1^3b_3+b_3^{})`$ $`+b_1(b_1^2+b_2b_2^{}+{\displaystyle \frac{11}{72}}){\displaystyle \frac{564731}{124416}}+{\displaystyle \frac{82043}{27648}}\zeta _3+{\displaystyle \frac{2633}{31104}}n_l\},`$ with $`l_h^{}=\mathrm{ln}(\mu _h^2/(\mathrm{\Lambda }^{})^2)`$. It is realized in the procedure DecLambdaUp\[\]. At this point we would like to mention that next to the coupling constant and quark masses also the gauge parameter and the quark and gluon fields obey decoupling relations. The corresponding equations and results can be found in . ## 5 Description of the main procedures In this section we describe the procedures which are most important for the practical applications, namely the combined running and decoupling of the strong coupling and the conversion of the on-shell mass to the $`\overline{\mathrm{MS}}`$ one and vice versa. In RunDec.m some masses and couplings are set to default values which are used if they are not specified explicitly. They are collected in the set NumDef and read (also the corresponding symbol used in RunDec is given): $`\begin{array}{cccccc}\mathrm{𝙼𝚝𝚊𝚞}:\hfill & M_\tau =1.777\text{ GeV},\hfill & \mathrm{𝙼𝚌}:\hfill & M_c=1.6\text{ GeV},\hfill & \mathrm{𝙼𝚋}:\hfill & M_b=4.7\text{ GeV},\hfill \\ \mathrm{𝙼𝚝}:\hfill & M_t=175\text{ GeV},\hfill & \mathrm{𝚖𝚞𝚌}:\hfill & \mu _c=1.2\text{ GeV},\hfill & \mathrm{𝚖𝚞𝚋}:\hfill & \mu _b=3.97\text{ GeV},\hfill \\ \mathrm{𝙼𝚣}:\hfill & M_Z=91.18\text{ GeV},\hfill & \mathrm{𝚊𝚜𝙼𝚣}:\hfill & \alpha _s^{(5)}(M_Z)=0.118.\hfill & & \end{array}`$ (39) The following procedure computes $`\alpha _s^{(m)}(\mu )`$ where $`\alpha _s^{(n)}(\mu _0)`$ is used as input parameter. As input only $`\alpha _s(\mu _0)`$, $`\mu _0`$, $`\mu `$ and the number of loops have to be specified. Both $`n`$ and $`m`$ are determined according to the values of the quark masses given in NumDef. In case $`nm`$ the heavy quarks are consistently decoupled at the heavy quark scale itself where for the mass definition the on-shell scheme is used. * AsRunDec: + input: $`\alpha _s^{(n)}(\mu _0)`$, $`\mu _0`$, $`\mu `$, number of loops + output: $`\alpha _s^{(m)}(\mu )`$ + uses: AlphasExact\[\], AlL2AlH\[\] and AlH2AlL\[\] + comments: The decoupling is performed automatically at the pole mass of the heavy quark where the values defined in Numdef are taken. If $`\mu `$ is lower than $`M_c`$, $`m=3`$ is chosen, i.e. the strange quark is not decoupled. + example: In order to compute $`\alpha _s^{(6)}(500\text{GeV})=0.952`$ with four-loop accuracy if $`\alpha _s^{(5)}(M_Z)=0.118`$ is given one has to use the command `AsRunDec[asMz/.NumDef,Mz/.NumDef,500,4]`. The conversion of the on-shell mass $`M`$ to the $`\overline{\mathrm{MS}}`$ mass, $`m`$, can be computed with the help of the procedure mOS2mMS\[\]: * mOS2mMS: + input: $`M`$, $`n_f`$, number of loops + output: $`m^{(n_f)}(M)`$ + uses: AsRunDec\[\] and mOS2mMS\[\] (from the appendix) + comments: The relation is implemented up to order $`\alpha _s^3`$ (three loops). $`n_f`$ is the number of active flavours. $`\alpha _s^{(n_f)}(M)`$ is evaluated at the scale $`M`$ where $`\alpha _s^{(5)}(M_Z)`$, as defined in Numdef, serves as a starting point. For the running and decoupling the procedure AsRunDec is used. + example: In the case of the top quark the $`\overline{\mathrm{MS}}`$ mass $`m_t(M_t)=164.6`$ GeV is obtained via `mOS2mMS[175,6,3]` where $`M_t=175`$ GeV has been chosen. The inverted relation is implemented in * mMS2mOS: + input: $`\mu _m=m^{(n_f)}(\mu _m)`$, $`n_f`$, number of loops + output: $`M`$ + uses: AsRunDec\[\] and mMS2mOS\[\] (from the appendix) + comments: The relation is implemented up to order $`\alpha _s^3`$ (three loops). $`n_f`$ is the number of active flavours. $`\alpha _s^{(n_f)}(\mu _m)`$ is evaluated at the scale $`\mu _m`$ where $`\alpha _s^{(5)}(M_Z)`$, as defined in Numdef, serves as a starting point. For the running and decoupling the procedure AsRunDec is used. + example: In the case of the top quark the on-shell mass $`M_t=174.7`$ GeV is obtained via `mMS2mOS[165,6,3]` where $`m_t(M_t)=165`$ GeV has been chosen. In the above procedures all quarks lighter than the one under consideration are assumed to be massless. More specialized procedures providing more freedom in the choice of parameters and the running presentations can be found in the Appendix. ## Acknowledgments We would like to thank R. Harlander for carefully reading the manuscript and for many valuable comments. Also comments from G. Rodrigo and T. van Ritbergen are gratefully acknowledged. This work was supported by DFG under Contract Ku 502/8-1 (DFG-Forschergruppe “Quantenfeldtheorie, Computeralgebra und Monte-Carlo-Simulationen”). ## Appendix: Detailed presentation of the Mathematica modules contained in RunDec In the following we list the procedures contained in the program package RunDec and provide a brief description. The order of the parameters specified in the field input corresponds to the order required in the Mathematica procedures. The precision used for most of the numerical evaluations is controlled with the variable $NumPrec. For the procedures involving, e.g., numerical solutions of differential equations or recursive solutions of equations the default precision of Mathematica is kept which is for all practical purposes more than enough. Note that often the precision requested for with $NumPrec can not be reached when the input data are only known to a few digits. AsRunDec\[\] is not listed as it can already be found in Section 5. ### Procedures related to the strong coupling constant * LamExpl: + input: $`\alpha _s^{(n_f)}(\mu )`$, $`\mu `$, $`n_f`$, number of loops + output: $`\mathrm{\Lambda }^{(n_f)}`$ + uses: Eq. (4) + comments: + example: From the knowledge of $`\alpha _s^{(5)}(M_Z)=0.118`$ the computation of $`\mathrm{\Lambda }^{(5)}=0.2089`$ to three-loop accuracy proceeds as follows: `LamExpl[asMz/.NumDef,Mz/.NumDef,5,3]`. * LamImpl: + input: $`\alpha _s^{(n_f)}(\mu )`$, $`\mu `$, $`n_f`$, number of loops + output: $`\mathrm{\Lambda }^{(n_f)}`$ + uses: Eq. (5) + comments: Solves Eq. (5) numerically for $`\mathrm{\Lambda }^{(n_f)}`$. + example: If $`\alpha _s^{(5)}(M_Z)=0.118`$ is given the computation of $`\mathrm{\Lambda }^{(5)}=0.2083`$ to three-loop accuracy proceeds as follows: `LamImpl[asMz/.NumDef,Mz/.NumDef,5,3]`. * AlphasLam: + input: $`\mathrm{\Lambda }^{(n_f)}`$, $`\mu `$, $`n_f`$, number of loops + output: $`\alpha _s^{(n_f)}(\mu )`$ + uses: Eq. (5) + comments: An explicit warning is printed on the screen if the ratio $`\mu /\mathrm{\Lambda }^{(n_f)}`$ is too small. + example: For $`\mathrm{\Lambda }^{(5)}=0.208`$ and $`M_b=4.7`$ GeV the value of $`\alpha _s^{(5)}(M_b)=0.2163`$ is obtained to three-loop accuracy with `AlphasLam[0.208,4.7,5,3]`. * AlphasExact: + input: $`\alpha _s^{(n_f)}(\mu _0)`$, $`\mu _0`$, $`\mu `$, $`n_f`$, number of loops + output: $`\alpha _s^{(n_f)}(\mu )`$ + uses: Eq. (1) + comments: Solves the differential equation numerically using $`\alpha _s(\mu _0)`$ as initial condition. An explicit warning is printed on the screen if the ratio $`\mu /\mathrm{\Lambda }^{(n_f)}`$ is too small where $`\mathrm{\Lambda }^{(n_f)}`$ is obtained with the help of LamExpl\[\]. + example: $`\alpha _s^{(5)}(M_b)=0.2167`$ is computed from $`\alpha _s^{(5)}(M_Z)`$ through `AlphasExact[asMz/.NumDef,Mz/.NumDef,Mb/.NumDef,5,3]` where the three-loop formulae are used. ### Procedures relating different mass definitions * mOS2mMS: + input: $`M`$, $`\{\text{M}_q\}`$, $`\alpha _s^{(n_f)}(\mu )`$, $`\mu `$, $`n_f`$, number of loops + output: $`m^{(n_f)}(\mu )`$ + uses: Eq. (13) and Tab. 2 + comments: The relation is implemented up to order $`\alpha _s^3`$ (three loops). $`\{\text{M}_q\}`$. is a set of light quark masses which can also be empty. For consistency reasons their values must correspond to the on-shell mass. Note that the name of the procedure is the same as the one introduced in Section 5. The distinction is only in the number of the arguments. + example: The $`\overline{\mathrm{MS}}`$ mass corresponding to the on-shell top quark mass of 175 GeV is computed via `mOS2mMS[175,{},0.107,175,6,3]` where $`\alpha _s^{(6)}(175\text{GeV})=0.107`$ has been chosen. The result reads $`m_t(175\text{GeV})=164.64`$ GeV. Terms up to order $`\alpha _s^3`$ have been used and light quark mass effects have been neglected. * mMS2mOS: + input: $`m^{(n_f)}(\mu )`$, $`\{m_q\}`$, $`\alpha _s^{(n_f)}(\mu )`$, $`\mu `$, $`n_f`$, number of loops + output: $`M`$ + uses: Eq. (LABEL:eq:zminv) for general $`\mu `$ and Tab. 2 + comments: The relation is implemented up to order $`\alpha _s^3`$ (three loops). In this case the light quark masses are defined in the $`\overline{\mathrm{MS}}`$ scheme $`\{m_q\}`$. Note that the name of the procedure is the same as the one introduced in Section 5. The distinction is only in the number of the arguments. + example: The on-shell mass corresponding to the $`\overline{\mathrm{MS}}`$ top quark mass $`m_t(175\text{GeV})=165`$ GeV is computed via `mMS2mOS[165,{},0.107,175,6,3]` where $`\alpha _s^{(6)}(175\text{GeV})=0.107`$ has been chosen. The result reads $`M_t=175.35`$ GeV. Terms up to order $`\alpha _s^3`$ have been used and light quark mass effects have been neglected. * mOS2mMSrun: + input: $`M`$, $`\{\text{M}_q\}`$, $`\alpha _s^{(n_f)}(\mu )`$, $`\mu `$, $`n_f`$, number of loops + output: $`m^{(n_f)}(\mu )`$ + uses: AlphasExact\[\], mOS2mSI\[\] and mMS2mMS\[\] + comments: In a first step $`\mu _m`$ is computed and afterwards $`m^{(n_f)}(\mu )`$ is evaluated. The usage is identical to mOS2mMS\[\]. + example: (analog to mOS2mMS\[\]) * mMS2mOSrun: + input: $`m^{(n_f)}(\mu )`$, $`\{m_q\}`$, $`\alpha _s^{(n_f)}(\mu )`$, $`\mu `$, $`n_f`$, number of loops + output: $`M`$ + uses: AlphasLam\[\], LamImpl\[\], mMS2mMS\[\] and mMS2mOS\[\] + comments: In a first step $`\mu _m`$ is computed. Then Eq. (LABEL:eq:zminv) only has to be used for $`\mu =\mu _m`$. The usage is identical to mMS2mOS\[\]. + example: (analog to mMS2mOS\[\]) * mOS2mMSit: + input: $`M`$, $`\{m_q\}`$, $`\alpha _s^{(n_f)}(\mu )`$, $`\mu `$, $`n_f`$, number of loops + output: $`m^{(n_f)}(\mu )`$ + uses: Eq. (LABEL:eq:zminv) for general $`\mu `$ and Tab. 2 + comments: For the computation Eq. (LABEL:eq:zminv) is used in order to avoid the on-shell masses of the light quark masses $`\{m_q\}`$. The usage is identical to mOS2mMS\[\]. + example: (analog to mOS2mMS\[\]). * mOS2mSI: + input: $`M`$, $`\{\text{M}_q\}`$, $`\alpha _s^{(n_f)}(M)`$, $`n_f`$, number of loops + output: $`\mu _m=m^{(n_f)}(\mu _m)`$ + uses: Eq. (16) and Tab. 2 + comments: The scale invariant mass is computed from the on-shell mass. + example: In the case of the bottom quark, the mass $`\mu _b=3.97`$ GeV is evaluated via `mOS2mSI[Mb/.NumDef,{1.6},0.217,5,3]` where $`\alpha _s^{(5)}(M_b)=0.217`$ has been chosen. In the mass relation terms up to order $`\alpha _s^3`$ have been used and quark mass effects arising from $`M_c/M_b`$ with $`M_c=1.6`$ GeV have been taken into account. * mMS2mMS: + input: $`m^{(n_f)}(\mu _0)`$, $`\alpha _s^{(n_f)}(\mu _0)`$, $`\alpha _s^{(n_f)}(\mu )`$, $`n_f`$, number of loops + output: $`m^{(n_f)}(\mu )`$ + uses: Eqs. (9) and (10) + comments: + example: From $`m_b(M_b)=3.85`$ GeV one finds $`m_b(M_Z)=2.69`$ GeV with the help of `mMS2mMS[3.85,0.217,asMz/.NumDef,5,4])` where $`\alpha _s^{(5)}(M_b)=0.217`$ and $`\alpha _s^{(5)}(M_Z)=0.217`$ has been used. For the running the four-loop expressions have been used. * mMS2mSI: + input: $`m^{(n_f)}(\mu )`$, $`\alpha _s^{(n_f)}(\mu )`$, $`\mu `$, $`n_f`$, number of loops + output: $`\mu _m=m(\mu _m)`$ + uses: Eqs. (9) and (10) + comments: The scale invariant mass is computed from the $`\overline{\mathrm{MS}}`$ mass. + example: $`\mu _b=3.97`$ GeV is computed from the input $`m_b(M_b)=3.85`$ GeV, $`M_b=4.7`$ GeV and $`\alpha _s^{(5)}(M_b)=0.217`$ via the command `mMS2mSI[3.85,0.217,4.7,5,4]`. For the running the four-loop expressions have been used. * mMS2mRI: + input: $`m^{(n_f)}(\mu )`$, $`\alpha _s^{(n_f)}(\mu )`$, $`n_f`$, number of loops + output: $`m^{RI}`$ + uses: inverted equation of (18) + comments: The relation is implemented up to order $`\alpha _s^3`$ (three loops). + example: The regularization invariant mass, $`m_b^{RI}(M_Z)`$, corresponding to the $`\overline{\mathrm{MS}}`$ bottom quark mass $`m_b(M_Z)=2.695`$ GeV is computed via `mMS2mRI[2.695,asMz/.NumDef,5,3]` where $`\alpha _s^{(5)}(M_Z)=0.118`$ has been chosen. The result reads $`m_b^{RI}(M_Z)=2.872`$ GeV where terms up to order $`\alpha _s^3`$ have been used. * mRI2mMS: + input: $`m^{RI}(\mu )`$, $`\alpha _s^{(n_f)}(\mu )`$, $`n_f`$, number of loops + output: $`m^{(n_f)}(\mu )`$ + uses: Eq. (18) + comments: The relation is implemented up to order $`\alpha _s^3`$ (three loops). + example: The $`\overline{\mathrm{MS}}`$ mass corresponding to the regularization invariant top quark mass of $`m_t^{RI}=175`$ GeV is computed via `mRI2mMS[175,0.107,175,6,3]` where $`\alpha _s^{(6)}(175\text{GeV})=0.107`$ has been chosen. The result reads $`m_t(175\text{GeV})=165.6`$ GeV where terms up to order $`\alpha _s^3`$ have been used. * mMS2mRGI: + input: $`m^{(n_f)}(\mu )`$, $`\alpha _s^{(n_f)}(\mu )`$, $`n_f`$, number of loops + output: $`\widehat{m}`$ + uses: Eq (11) + comments: + example: The renormalization group invariant bottom quark mass corresponding to the $`\overline{\mathrm{MS}}`$ mass $`m_b^{(5)}(M_Z)=2.69`$ GeV is computed via `mMS2mRGI[2.69,asMz/.NumDef,5,4]` where $`\alpha _s^{(5)}(M_Z)=0.118`$ has been chosen. The result reads $`m_b^{RGI}=14.25`$ GeV assuming four-loop accuracy. * mRGI2mMS: + input: $`\widehat{m}`$, $`\alpha _s^{(n_f)}(\mu )`$, $`n_f`$, number of loops + output: $`m^{(n_f)}(\mu )`$ + uses: Eq (11) + comments: + example: The $`\overline{\mathrm{MS}}`$ mass corresponding to the renormalization group invariant bottom quark mass of 14.25 GeV is computed via `mRGI2mMS[14.25,asMz/.NumDef,5,4]` where $`\alpha _s^{(5)}(M_Z)=0.118`$ has been chosen. The result reads $`m_b(M_Z)=2.69`$ GeV assuming four-loop accuracy. ### Decoupling of the strong coupling and the masses At this point we once again want to stress, that the argument specifying the number of loops refers to the accompanied running, i.e. if “2” is chosen the decoupling relation is used to one-loop order. Furthermore, for the argument ruling the number of active flavours the number of light quarks, $`n_l=n_f1`$, is chosen. * DecAsUpOS: + input: $`\alpha _s^{(n_l)}(\mu _{th})`$, $`M_{th}`$, $`\mu _{th}`$, $`n_l`$, number of loops + output: $`\alpha _s^{(n_l+1)}(\mu _{th})`$ + uses: Eq. (25) + comments: For the heavy mass the on-shell definition is used. + example: The computation of $`\alpha _s^{(6)}(M_Z)=0.1169`$ from the knowledge of $`\alpha _s^{(5)}(M_Z)=0.118`$ proceeds via `DecAsUpOS[asMz/.NumDef,175,Mz/.NumDef,5,4]` where $`M_t=175`$ GeV has been chosen and terms of order $`\alpha _s^3`$ (indicated by the “4” in the last argument) have been included. * DecAsDownOS: + input: $`\alpha _s^{(n_l+1)}(\mu _{th})`$, $`M_{th}`$, $`\mu _{th}`$, $`n_l`$, number of loops + output: $`\alpha _s^{(n_l)}(\mu _{th})`$ + uses: Eq. (22) + comments: For the heavy mass the on-shell definition is used. + example: The computation of $`\alpha _s^{(5)}(200\text{GeV})=0.1047`$ from the knowledge of $`\alpha _s^{(6)}(200\text{GeV})=0.105`$ proceeds via `DecAsDownOS[0.105,175,200,6,4]` where $`M_t=175`$ GeV has been chosen and terms of order $`\alpha _s^3`$ (indicated by the “4” in the last argument) have been included. * DecAsUpMS: + input: $`\alpha _s^{(n_l)}(\mu _{th})`$, $`m_{th}(\mu _{th})`$, $`\mu _{th}`$, $`n_l`$, number of loops + output: $`\alpha _s^{(n_l+1)}(\mu _{th})`$ + uses: Eq. (23) + comments: The heavy mass is evaluated in the $`\overline{\mathrm{MS}}`$ scheme at the scale $`\mu _{th}`$. + example: The computation of $`\alpha _s^{(6)}(M_Z)=0.1170`$ from the knowledge of $`\alpha _s^{(5)}(M_Z)=0.118`$ proceeds via `DecAsUpMS[asMz/.NumDef,165,Mz/.NumDef,5,4]` where $`m_t(M_Z)=165`$ GeV has been chosen and terms of order $`\alpha _s^3`$ (indicated by the “4” in the last argument) have been included. * DecAsDownMS: + input: $`\alpha _s^{(n_l+1)}(\mu _{th})`$, $`m_{th}(\mu _{th})`$, $`\mu _{th}`$, $`n_l`$, number of loops + output: $`\alpha _s^{(n_l)}(\mu _{th})`$ + uses: Eq. (20) + comments: The heavy mass is evaluated in the $`\overline{\mathrm{MS}}`$ scheme at the scale $`\mu _{th}`$. + example: The computation of $`\alpha _s^{(5)}(200\text{GeV})=0.1048`$ from the knowledge of $`\alpha _s^{(6)}(200\text{GeV})=0.105`$ proceeds via `DecAsDownMS[0.105,165,200,6,4]` where $`m_t(M_Z)=165`$ GeV has been chosen and terms of order $`\alpha _s^3`$ (indicated by the “4” in the last argument) have been included. * DecAsUpSI: + input: $`\alpha _s^{(n_l)}(\mu _{th})`$, $`\mu _{m_{th}}`$, $`\mu _{th}`$, $`n_l`$, number of loops + output: $`\alpha _s^{(n_l+1)}(\mu _{th})`$ + uses: Eq. (24) + comments: Here the scale invariant mass $`\mu _{m_{th}}`$ is chosen for heavy mass. + example: (analog to DecAsUpMS\[\]) * DecAsDownSI: + input: $`\alpha _s^{(n_l+1)}(\mu _{th})`$, $`\mu _{m_{th}}`$, $`\mu _{th}`$, $`n_l`$, number of loops + output: $`\alpha _s^{(n_l)}(\mu _{th})`$ + uses: Eq. (21) + comments: Here the scale invariant mass $`\mu _{m_{th}}`$ is chosen for heavy mass. + example: (analog to DecAsDownMS\[\]) * DecMqUpOS: + input: $`m_q^{(n_l)}(\mu _{th})`$, $`\alpha _s^{(n_l)}(\mu _{th})`$, $`M_{th}`$, $`\mu _{th}`$, $`n_l`$, number of loops + output: $`m_q^{(n_l+1)}(\mu _{th})`$ + uses: Eq. (33) + comments: For the heavy mass the on-shell definition is used. + example: The computation of $`m_b^{(6)}(M_Z)=2.697`$ GeV from $`m_b^{(5)}(M_Z)=2.7`$ GeV with order $`\alpha _s^3`$ accuracy is performed via `DecMqUpOS[2.7,asMz/.NumDef,175,Mz/.NumDef,5,4]`. Here, $`\alpha _s^{(5)}(M_Z)=0.118`$ and $`M_t=175`$ GeV have been used. * DecMqDownOS: + input: $`m_q^{(n_l+1)}(\mu _{th})`$, $`\alpha _s^{(n_l+1)}(\mu _{th})`$, $`M_{th}`$, $`\mu _{th}`$, $`n_l`$, number of loops + output: $`m_q^{(n_l)}(\mu _{th})`$ + uses: Eq. (30) + comments: For the heavy mass the on-shell definition is used. + example: The computation of $`m_c^{(4)}(M_Z)=0.583`$ GeV from $`m_c^{(5)}(M_Z)=0.58`$ GeV with order $`\alpha _s^3`$ accuracy is performed via `DecMqDownOS[0.58,asMz/.NumDef,4.7,Mz/.NumDef,5,4]`. Here, $`\alpha _s^{(5)}(M_Z)=0.118`$ and $`M_b=4.7`$ GeV have been used. * DecMqUpMS: + input: $`m_q^{(n_l)}(\mu _{th})`$, $`\alpha _s^{(n_l)}(\mu _{th})`$, $`m_{th}(\mu _{th})`$, $`\mu _{th}`$, $`n_l`$, number of loops + output: $`m_q^{(n_l+1)}(\mu _{th})`$ + uses: Eq. (31) + comments: The heavy mass is evaluated in the $`\overline{\mathrm{MS}}`$ scheme at the scale $`\mu _{th}`$. + example: (analog to DecMqUpOS\[\]) * DecMqDownMS: + input: $`m_q^{(n_l+1)}(\mu _{th})`$, $`\alpha _s^{(n_l+1)}(\mu _{th})`$, $`m_{th}(\mu _{th})`$, $`\mu _{th}`$, $`n_l`$, number of loops + output: $`m_q^{(n_l)}(\mu _{th})`$ + uses: Eq. (27) + comments: The heavy mass is evaluated in the $`\overline{\mathrm{MS}}`$ scheme at the scale $`\mu _{th}`$. + example: (analog to DecMqDownOS\[\]) * DecMqUpSI: + input: $`m_q^{(n_l)}(\mu _{th})`$, $`\alpha _s^{(n_l)}(\mu _{th})`$, $`\mu _{m_th}`$, $`\mu _{th}`$, $`n_l`$, number of loops + output: $`m_q^{(n_l+1)}(\mu _{th})`$ + uses: Eq. (32) + comments: Here the scale invariant mass $`\mu _{m_{th}}`$ is chosen for heavy mass. + example: (analog to DecMqUpOS\[\]) * DecMqDownSI: + input: $`m_q^{(n_l+1)}(\mu _{th})`$, $`\alpha _s^{(n_l+1)}(\mu _{th})`$, $`\mu _{m_th}`$, $`\mu _{th}`$, $`n_l`$, number of loops + output: $`m_q^{(n_l)}(\mu _{th})`$ + uses: Eq. (29) + comments: Here the scale invariant mass $`\mu _{m_{th}}`$ is chosen for heavy mass. + example: (analog to DecMqUpMS\[\]) * DecLambdaUp: + input: $`\mathrm{\Lambda }^{(n_l)}`$, $`\mu _{m_th}`$, $`n_l`$, number of loops + output: $`\mathrm{\Lambda }^{(n_l+1)}`$ + uses: Eq. (34) + comments: For the heavy mass the scale invariant mass $`\mu _{m_{th}}`$ is used. + example: From $`\mathrm{\Lambda }^{(4)}=0.2876`$ one can compute $`\mathrm{\Lambda }^{(5)}=0.208`$ with the help of `DecLambdaUp[0.287,3.97,4,4]` where $`\mu _b=3.97`$ and four-loop accuracy has been chosen. * DecLambdaDown: + input: $`\mathrm{\Lambda }^{(n_l+1)}`$, $`\mu _{m_th}`$, $`n_l`$, number of loops + output: $`\mathrm{\Lambda }^{(n_l)}`$ + uses: Eq. (35) + comments: For the heavy mass the scale invariant mass $`\mu _{m_{th}}`$ is used. + example: From $`\mathrm{\Lambda }^{(5)}=0.208`$ one can compute $`\mathrm{\Lambda }^{(4)}=0.208`$ with the help of `DecLambdaDown[0.208,3.97,4,4]` where $`\mu _b=3.97`$ and four-loop accuracy has been chosen. ### Miscellaneous procedures The following modules provide some simple examples which mostly combine the modules described above. “L” stands for low and “H” for high. The condition $`l<h`$ is assumed in all four procedures. * AlL2AlH: + input: $`\alpha _s^{(l)}(\mu _1)`$, $`\mu _1`$, $`\{\{n_{f_1},M_{th_1},\mu _{th_1}\},\{n_{f_2},M_{th_2},\mu _{th_2}\},\mathrm{}\}`$, $`\mu _2`$, number of loops + output: $`\alpha _s^{(h)}(\mu _2)`$ + uses: AlphasExact\[\] and DecAsUpOS\[\] + comments: The set in the third argument may contain several triples indicating the number of flavours, the heavy (on-shell) quark mass and the scale at which the decoupling is performed. + examples: 1. For the computation of $`\alpha _s^{(6)}(500\text{GeV})=0.0952`$ from $`\alpha _s^{(4)}(M_c=1.6\text{GeV})=0.338`$ to $`𝒪(\alpha _s^3)`$ accuracy the input would look as follows: `AlL2AlH[0.338,1.6,{{5,4.7,5},{6,175,200}},500,4]` Here, the matching is performed at 5 GeV and 200 GeV, respectively. 2. Fig. 1 can be reproduced with the help of the following input `AlL2AlH[0.36,1.777,{{5,4.7,mu5}},91.187,l]` where $`l=1,2,3,4`$ corresponds to the number of loops and `mu5/4.7` is the scale on the abscissa. * AlH2AlL: + input: $`\alpha _s^{(h)}(\mu _1)`$, $`\mu _1`$, $`\{\{n_{f_1},M_{th_1},\mu _{th_1}\},\{n_{f_2},M_{th_2},\mu _{th_2}\},\mathrm{}\}`$, $`\mu _2`$, number of loops + output: $`\alpha _s^{(l)}(\mu _2)`$ + uses: AlphasExact\[\] and DecAsDownOS\[\] + comments: The set in the third argument may contain several triples indicating the number of flavours, the heavy (on-shell) quark mass and the scale at which the decoupling is performed. + example: Consider the inverse order of the first example of the previous procedure. The input `AlH2AlL[0.0952,500,{{6,175,200},{5,4.7,5}},1.6,4]`, indeed leads to $`\alpha _s^{(4)}(1.6\text{GeV})=0.338`$. * mL2mH: + input: $`m_q^{(l)}(\mu _1)`$, $`\alpha _s^{(l)}(\mu _1)`$, $`\mu _1`$, $`\{\{n_{f_1},M_{th_1},\mu _{th_1}\},\{n_{f_2},M_{th_2},\mu _{th_2}\},\mathrm{}\}`$, $`\mu _2`$, number of loops + output: $`m_q^{(h)}(\mu _2)`$ + uses: AlphasExact\[\], mMS2mMS\[\] DecMqUpOS\[\] and DecAsUpOS\[\] + comments: The set in the fourth argument may contain several triples indicating the number of flavours, the heavy (on-shell) quark mass and the scale at which the decoupling is performed. + example: Using $`\alpha _s^{(4)}(1.2\text{GeV})=0.403`$ and $`m_c^{(4)}(1.2\text{GeV})=1.2`$ GeV one finds $`m_c^{(5)}(M_Z)=0.580`$ GeV with the help of `mL2mH[1.2,0.403,1.2,{{5,4.7,5.0}},Mz/.NumDef,4]`. The decoupling of $`M_b=4.7`$ GeV is performed at $`5.0`$ GeV. In this way the results of Fig. 2 can be reproduced. * mH2mL: + input: $`m_q^{(h)}(\mu _1)`$, $`\alpha _s^{(h)}(\mu _1)`$, $`\mu _1`$, $`\{\{n_{f_1},M_{th_1},\mu _{th_1}\},\{n_{f_2},M_{th_2},\mu _{th_2}\},\mathrm{}\}`$, $`\mu _2`$, number of loops + output: $`m_q^{(l)}(\mu _2)`$ + uses: AlphasExact\[\], mMS2mMS\[\] DecMqDownOS\[\] and DecAsDownOS\[\] + comments: The set in the fourth argument may contain several triples indicating the number of flavours, the heavy (on-shell) quark mass and the scale at which the decoupling is performed. + example: Using $`\alpha _s^{(5)}(M_Z)=0.118`$ and $`m_c^{(5)}(M_Z)=0.580`$ GeV one finds $`m_c^{(4)}(1.2\text{GeV})=1.20`$ GeV with the help of `mH2mL[0.580,asMz/.NumDef,Mz/.NumDef,{{5,4.7,5.0}},1.2,4]`. The decoupling of $`M_b=4.7`$ GeV is performed at $`5.0`$ GeV.
warning/0004/nlin0004022.html	ar5iv	text	# 1 Introduction ## 1 Introduction <br> The purpose of this article is to investigate the Lie point symmetries of a large class of “generalized Toda field theories”. The class is characterized by the equation $$u_{n,xy}=F_n,F_n=\underset{m=nn_1}{\overset{n+n_2}{}}K_{nm}\mathrm{exp}\left(\underset{l=mn_3}{\overset{m+n_4}{}}H_{ml}u_l\right)$$ (1.1) where $`K`$ and $`H`$ are some real constant matrices and $`n_1,\mathrm{},n_4`$ are some finite non-negative integers. The range of $`n`$ may be infinite, semi-infinite or finite, hence the matrices $`K`$ and $`H`$ may also be infinite, semi-infinite, or finite. If the range of $`n`$ is finite, $`K`$ and $`H`$ may be rectangular, not necessarily square. We assume that all the rows in $`H`$ are different, that $`H`$ contains no zero rows and $`K`$ no zero columns. In all the cases we assume that the range of the interaction on the right hand side of eq. (1.1) is finite, hence the finite summation limits in both sums. “Generalized Toda lattices” are obtained from eq. (1.1) by symmetry reduction, using translational invariance, i.e. restricting to solutions of the form $`u_n(x,y)=w_n(t)`$ where $`t=x+\lambda y`$. Toda lattices and their generalisations, Toda field theories, represent one of the most interesting, rich and fruitful developments in the realm of completely integrable systems. The original Toda lattice was introduced by M. Toda who found analytical solitons and periodic solutions in a discrete lattice with an exponential potential involving nearest neighbour interactions. It was also found that the Toda lattice admits a Lax representation and all the usual attributes of integrability . The Toda lattice was generalized to integrable lattices related to the root systems of simple Lie algebras \- . The considered lattices can be finite, infinite, semi-infinite, or periodic. The attractive features of Toda lattices have been generalized to two space dimensions in several different ways \- . All of them can be recovered from eq. (1.1) by specifying the matrices $`K`$ and $`H`$. Thus, the Mikhailov-Fordy-Gibbons field theories (for infinitely many fields) $$u_{n,xy}=e^{u_{n1}u_n}e^{u_nu_{n+1}}$$ (1.2) are obtained by putting $`H_{nn1}H_{nn}=1`$, $`K_{nn}=K_{nn+1}=1`$ and all other components to zero. A class of Toda field theories $$u_{n,xy}=\underset{m=nn_1}{\overset{n+n_2}{}}K_{nm}e^{u_m},$$ (1.3) studied by Leznov and Saveliev , Olive, Turok and others \- (usually for a finite number of fields $`u_n`$) are obtained by setting $`H=I`$ and taking $`K`$ to be the Cartan matrix of a semisimple Lie algebra (or an affine one). A further class of Toda field theories, also studied by Leznov and Saveliev , by Bilal and Gervais , and Babelon and Bonora (for a finite number of fields) can be written as $$u_{n,xy}=\mathrm{exp}\underset{l=mn_3}{\overset{m+n_4}{}}H_{nl}u_l$$ (1.4) and is obtained by taking $`K=I`$ and $`H`$ as a Cartan matrix. In this article we will be interested in point symmetries of the system (1.1), rather than in questions of integrability, or explicit solutions. The symmetries we are interested in will include conformal invariance, whenever it is present, and gauge invariance, not however higher, or generalized symmetries, be they local, or not. In Section 2 we consider infinite Toda field theories, i.e. take $`\mathrm{}<n<\mathrm{}`$. In this case eq. (1.1) can be viewed as a differential-difference equation. Continuous Lie symmetries of such equations have been studied using several different approaches \- . We shall follow that of Ref. \- , using both the “intrinsic method” and the “differential equation method” . In Section 3 we turn to finite Toda field theories, when we have $`1nN<\mathrm{}`$ in eq. (1.1). Eq. (1.1) in this case represents a system of $`N`$ differential equations and its point symmetries can be obtained in a standard manner . We first obtain general results, then specify the matrices $`H`$ and $`K`$ in several different ways. Section 4 is devoted to semi-infinite Toda field theories, i.e. $`0n<\mathrm{}`$. Again we first obtain general results, then specify the matrices $`H`$ and $`K`$, inforcing the cut-off at $`n=0`$ in several different ways. Some conclusions are drawn in Section 5. ## 2 Symmetries of Generalized $`\mathrm{}`$Toda Field Theories ### 2.1 <br>General Results Let us consider eq. (1.1) with $`n`$ in the range $`\mathrm{}<n<\mathrm{}`$. We follow the “differential equation method” described in Ref. and look for transformations of the form $$\stackrel{~}{\stackrel{}{x}}=\mathrm{\Lambda }_g(\stackrel{}{x},\{u_k\}),\stackrel{~}{u}_n=\mathrm{\Omega }_g(\stackrel{}{x},n,\{u_k\}),\stackrel{~}{n}=n,$$ (2.1) where we have used the notation $`\stackrel{}{x}(x,y)`$, $`\stackrel{~}{\stackrel{}{x}}(\stackrel{~}{x},\stackrel{~}{y})`$, taking solutions of eq. (1.1) into solutions. The notation $`\{u_k\}`$ indicates that the new variables can depend on all the fields $`\{u_k\}_{k𝐙}`$. The Lie group transformation (2.1) is generated by a Lie algebra of vector fields of the form $$\widehat{v}=\xi (x,y,\{u_k\})_x+\eta (x,y,\{u_k\})_y+\underset{j=\mathrm{}}{\overset{\mathrm{}}{}}\varphi _j(x,y,\{u_k\})_{u_j}.$$ (2.2) The prolongation of this vector field is constructed in the same manner as for differential equations (albeit an infinite system of them). For a general equation of the form $$E_n=u_{n,xy}F_n(x,y,\{u_k\})=0,$$ (2.3) we require $$pr^{\left(2\right)}\widehat{v}E_n_{E_n=0}=0.$$ (2.4) It was shown quite generally that for eq (2.3) with $`F_n`$ any sufficiently smooth function depending on at least one function $`u_k`$, $`kn`$, the vector field (2.2) satisfying eq. (2.4) will have the form $$\xi =\xi (x),\eta =\eta (y),\varphi _n=\underset{k=\mathrm{}}{\overset{\mathrm{}}{}}A_{nk}u_k+B_n(x,y),$$ (2.5) where $`A=\{A_{n\alpha }\}`$ is a constant (infinite) matrix. The functions in eq. (2.5) must satisfy a remaining determining equation, namely $$\begin{array}{c}B_{n,xy}(\xi _x+\eta _y)F_n+\underset{\alpha =\mathrm{}}{\overset{\mathrm{}}{}}A_{n\alpha }F_\alpha \xi F_{n,x}\eta F_{n,y}\hfill \\ \\ \underset{\alpha =\mathrm{}}{\overset{\mathrm{}}{}}\left(\underset{\beta =\mathrm{}}{\overset{\mathrm{}}{}}A_{\alpha \beta }u_\beta +B_\alpha \right)F_{n,u_\alpha }=0,\hfill \end{array}$$ (2.6) where $`F_{n,u_\alpha }`$ is the derivative of $`F_n`$ with respect to the variable $`u_\alpha `$. Let us now specify the function $`F_n`$ to be a sum of exponentials as in eq. (1.1). There are three types of terms in eq. (2.6): those independent of $`u_n`$, linear in $`u_n`$ times exponentials and pure exponentials. Each type of term must vanish separately. Since $`H`$ has no zero rows we get the determining equations $$B_{n,xy}=0,$$ (2.7) $$\underset{\alpha =\mathrm{}}{\overset{\mathrm{}}{}}A_{\alpha m}F_{n,u_\alpha }=0,$$ (2.8) $$(\xi _x+\eta _y)F_n+\underset{\alpha =\mathrm{}}{\overset{\mathrm{}}{}}A_{n\alpha }F_\alpha \underset{\alpha =\mathrm{}}{\overset{\mathrm{}}{}}B_\alpha F_{n,u_\alpha }=0.$$ (2.9) Eq. (2.8) can be rewritten as $$\underset{\alpha \beta }{}K_{n\beta }H_{\beta \alpha }A_{\alpha m}\mathrm{exp}\left(\underset{\gamma }{}H_{\beta \gamma }u_\gamma \right)=0.$$ (2.10) All exponentials in eq. (2.10) are linearly independent (since all rows in $`H`$ are different), so the equation must hold for each $`\beta `$ separately and the exponentials can be dropped. Moreover, the factor $`K_{n\beta }`$ can be dropped (since $`K`$ has no zero column). We find that eq. (2.8) in this case implies an equation for the matrix $`A`$, namely $$\underset{\alpha =\mathrm{}}{\overset{\mathrm{}}{}}H_{n\alpha }A_{\alpha m}=0,$$ (2.11) or in matrix form $`HA=0`$ (however, the matrices are infinite). Let us now turn to eq. (2.9) and make use of the finite range of the interaction $`F_n`$ in eq. (1.1). We have $$\frac{F_n}{u_k}=0,n+n_u<k\mathrm{or}k<nn_d$$ (2.12) for some non-negative integers $`n_u`$ and $`n_d`$. In eq. (2.9) all exponentials, obtained after substituing for $`F_n`$ from eq. (1.1), are linearly independent. This allows us to split eq. (2.9) into two types of equations. These are obtained as coefficients of $`\mathrm{exp}\left(_lH_{ml}u_l\right)`$, with $`m[nn_1,n+n_2]`$ and with $`m`$ outside this interval, respectively. Thus we have: $$\begin{array}{c}\hfill K_{nm}\left[(\xi _x+\eta _y)+\underset{\alpha =mn_3}{\overset{m+n_4}{}}B_\alpha H_{m\alpha }\right]+\underset{\rho =mn_1}{\overset{m+n_2}{}}A_{n\rho }K_{\rho m}=0,\\ \\ \hfill m[nn_1,n+n_2],\end{array}$$ (2.13) $$\underset{\rho =mn_1}{\overset{m+n_2}{}}A_{n\rho }K_{\rho m}=0,m\overline{)}[nn_1,n+n_2].$$ (2.14) We shall show that eq. (2.14) actually holds for all values of $`m`$ so that eq. (2.13) can be simplified. To do this, we view eq. (2.11) as a difference equation for $`A_{\alpha m}`$. To make this explicit we restrict $`H`$ and $`K`$ to be band matrices, with finite bands of constant width $$H_{nm}=H_{n,n+\sigma }=\{\begin{array}{c}h_\sigma (n)\sigma [p_1,p_2]\hfill \\ 0\sigma \overline{)}[p_1,p_2]\hfill \end{array},h_{p_1}(n)0,h_{p_2}(n)0.$$ (2.15) Similarly $$K_{nm}=K_{m+\sigma ,m}=\{\begin{array}{c}k_\sigma (m)\sigma [q_1,q_2]\hfill \\ 0\sigma \overline{)}[q_1,q_2]\hfill \end{array},k_{q_1}(m)0,k_{q_2}(m)0.$$ (2.16) In these notations we see that eq. (2.11) is a linear difference equation for $`A_{\sigma m}`$ with $`p_1p_2+1`$ terms $$\underset{\sigma =p_1}{\overset{p_2}{}}h_\sigma (n)A_{\sigma +n,m}=0.$$ (2.17) Equation (2.17) determines the dependence of $`A_{nm}`$ on $`n`$. Indeed the linear difference equation $$\underset{\sigma =p_1}{\overset{p_2}{}}h_\sigma (n)\psi _{\sigma +n}=0$$ (2.18) has $`p_2p_1`$ linearly independent solutions, a basis of which we denote by $`\left\{\psi _n^j,j=1,2,\mathrm{},p_2p_1\right\}`$. Thus, we have $$A_{nm}=\underset{j=1}{\overset{p_2p_1}{}}\psi _n^jC_{jm},$$ (2.19) where $`C_{jm}`$ are constants to be determined by the remaining determining equations (2.13) and (2.14). In order to analyze them, let us define the quantities $$Q_{nm}=\underset{\sigma =mn_1}{\overset{m+n_2}{}}A_{n\sigma }K_{\sigma m}.$$ From eq. (2.14) we have $`Q_{nm}=0`$ for $`m`$ “sufficiently far away” from $`n`$. But, by using the expansion (2.19), we get $$Q_{nm}=\underset{j=1}{\overset{p_2p_1}{}}\psi _n^j\underset{\sigma =mn_1}{\overset{m+n_2}{}}C_{j\sigma }K_{\sigma m}$$ which, because of the linear independency of the $`\psi _n^j`$, implies $$\underset{\sigma =mn_1}{\overset{m+n_2}{}}C_{j\sigma }K_{\sigma m}=0$$ (2.20) for all values of $`m`$, since this relation does not depend on $`n`$ and the index $`m`$ is no longer tied to $`n`$. In other words, if $`Q_{nm}=0`$ holds for certain values of $`n`$ and $`m`$, as in eq. (2.14), then that equation must hold for all values. As in the case of eq. (2.17), we introduce a solution basis $`\left\{\varphi _m^l,l=1,\mathrm{},q_2q_1\right\}`$ for the equation $$\underset{\sigma =q_1}{\overset{q_2}{}}k_\sigma \left(m\right)\varphi _{\sigma +m}=0.$$ (2.21) The general solution of eq. (2.20) now is $$C_{jm}=\underset{l=1}{\overset{q_2q_1}{}}q_{jl}\varphi _m^l,$$ where $`q_{jl}`$ are arbitrary constants. The expression (2.19) for $`A_{nm}`$ is replaced by $$A_{nm}=\underset{j=1}{\overset{p_2p_1}{}}\underset{l=1}{\overset{q_2q_1}{}}q_{jl}\psi _n^j\varphi _m^l.$$ (2.22) A further consequence is that the last term in eq. (2.13) can be dropped. Then, using the general solution for eq. (2.7) $$B_n(x,y)=\beta _n\left(x\right)+\gamma _n\left(y\right),$$ we separate the $`x`$ from the $`y`$ dependence in eq. (2.13) and reduce it to two inhomogeneous difference equations for $`\beta _n\left(x\right)`$ and $`\gamma _n\left(y\right)`$. The general solutions of which are $$\beta _n\left(x\right)=\stackrel{p_2p_1}{\underset{j=1}{}}r_j\left(x\right)\psi _n^jb_n\xi _x(x),\gamma _n\left(x\right)=\stackrel{p_2p_1}{\underset{j=1}{}}s{}_{j}{}^{}\left(y\right)\psi _n^jb_n\eta _y(y),$$ (2.23) where $`b_n`$ is an arbitrarily chosen solution of the inhomogeneous difference equation $$\underset{\sigma =p_1}{\overset{p_2}{}}h_\sigma (n)b_{\sigma +n}=1.$$ (2.24) Furthermore, in eq. (2.23) the functions $`r_j\left(x\right)`$ and $`s_j\left(y\right)`$ are arbitrarily chosen. Finally, we obtain the following theorem. ###### Theorem 1 Consider all the generalized Toda theories of the form (1.1) for infinitely many fields $`u_n(x,y)`$, where the coupling matrices $`H`$ and $`K`$ satisfy eqs. (2.15) and (2.16). Their Lie point symmetry algebra is infinite-dimensional and a basis for it is given by the following vector fields: $$\widehat{X}(\xi )=\xi (x)_x\xi _x\left(x\right)\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}b_n_{u_n},\widehat{Y}(\eta )=\eta (y)_y\eta _y\left(y\right)\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}b_n_{u_n},$$ (2.25) $$\widehat{U}_j\left(r_j\right)=r_j\left(x\right)\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}\psi _n^j_{u_n},\widehat{V}_j\left(s_j\right)=s_j\left(y\right)\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}\psi _n^j_{u_n}\left(j=1,\mathrm{},p_2p_1\right),$$ (2.26) $$\widehat{Z}_{jl}=\left(\underset{m=\mathrm{}}{\overset{\mathrm{}}{}}\varphi _m^lu_m\right)\left(\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}\psi _n^j_{u_n}\right)\left(j=1,\mathrm{},p_2p_1;l=1,\mathrm{},q_2q_1\right).$$ (2.27) The functions $`\xi (x),\eta (y),r_j\left(x\right)`$ and $`s_j\left(y\right)`$ are arbitrary, all the other quantities are determined by solving the linear difference eqs. (2.18), (2.21) and (2.24). As far as interpretation is concerned, we see that the generalized $`\mathrm{}`$ Toda lattice (1.1) is always conformally invariant, since the vector fields (2.25) generate arbitrary reparametrizations of $`x`$ and $`y`$, accompanied by appropriate transformations of the fields $`u_n`$. More specifically, the conformal transformations leaving eq. (1.1) invariant are $$\begin{array}{c}\stackrel{~}{x}=\psi (x,\lambda ),\stackrel{~}{y}=\chi (y,\lambda ),\hfill \\ \\ \stackrel{~}{u}_n(\stackrel{~}{x},\stackrel{~}{y})=u_n(x,y)b_n\mathrm{ln}(\frac{\text{d}\psi }{\text{d}x}\frac{\text{d}\chi }{\text{d}y}),\hfill \end{array}$$ (2.28) where $`\psi (x,\lambda )`$ and $`\chi (y,\lambda )`$ are arbitrary functions of $`x`$ and $`y`$, related to $`\xi (x)`$ and $`\eta (y)`$ by the relations $$\begin{array}{c}\stackrel{~}{x}=\psi (x,\lambda )=T^1(\lambda +T(x)),\hfill \\ \\ \stackrel{~}{y}=\chi (y,\lambda )=S^1(\lambda +S(y)),\hfill \end{array}$$ (2.29) with $$T(x)=_0^x\frac{\text{d}s}{\xi (s)},S(y)=_0^y\frac{\text{d}t}{\eta (t)}.$$ (2.30) The vector fields $`\widehat{U}_j\left(r\right)`$ and $`\widehat{V}_j\left(s\right)`$ generate gauge transformations: certain functions obtained by integrating the vector fields can be added to any solution. Formally, the operators $`\widehat{Z}_{jl}`$ generate linear transformations among components of solutions. However, the sums are over infinite range, so convergence problems may arise. Moreover, we have $$_{xy}\left(\underset{m}{}\varphi _m^lu_m\right)=0$$ (2.31) as a consequence of eq. (2.21). In other words, if the equation (2.21) admits non trivial solutions, than one can always perform a linear transformation among the $`u_n`$’s, in such a way $`q_2q_1`$ new fields $`v_l=_m\varphi _m^lu_m`$, satifying the wave equation $`_x_yv_l=0`$, are replaced in the Toda system. As stated in Theorem 1, the problem of finding all symmetries of eq. (1.1) reduces to solving the recursion relations (2.18), (2.21) and (2.24). In general, this may not be possible analytically in closed form. Well developed techniques exist for solving homogeneous and inhomogeneous difference equations with constant coefficients . This is the case that occurs for all generalized Toda field theories that we found in the literature: $`h_\sigma \left(n\right)`$ and $`k_\sigma \left(m\right)`$ do not depend on $`n`$ and $`m`$, respectively. The nonzero commutation relations for the symmetry algebra of the generalized $`\mathrm{}`$Toda theory (1.1) are: $$\begin{array}{cc}\hfill [\widehat{X}(\xi _1),\widehat{X}(\xi _2)]=\widehat{X}(\xi _1\xi _{2,x}\xi _{1,x}\xi _2),& [\widehat{Y}(\eta _1),\widehat{Y}(\eta _2)]=\widehat{Y}(\eta _1\eta _{2,y}\eta _{1,y}\eta _2),\hfill \\ & \\ \hfill [\widehat{X}(\xi ),\widehat{U}_j\left(r\right)]=\widehat{U}_j\left(\xi r_x\right),& [\widehat{Y}(\eta ),\widehat{V}_j\left(s\right)]=\widehat{V}_j\left(\eta s_y\right),\hfill \\ & \\ \hfill [\widehat{X}(\xi ),\widehat{Z}_{jl}]=\widehat{U}_j\left(\xi _x\underset{n}{}b_n\varphi _n^l\right),& [\widehat{Y}(\eta ),\widehat{Z}_{jl}]=\widehat{V}_j\left(\eta _y\underset{n}{}b_n\varphi _n^l\right),\hfill \\ & \\ \hfill [\widehat{U}_a\left(r\right),\widehat{Z}_{jl}]=\widehat{U}_j\left(r\underset{m}{}\varphi _m^l\psi _m^a\right),& [\widehat{V}_a\left(s\right),\widehat{Z}_{jl}]=\widehat{V}_j\left(s\underset{m}{}\varphi _m^l\psi _m^a\right),\hfill \\ & \\ \hfill [\widehat{Z}_{ab},\widehat{Z}_{cd}]=\left(\underset{m}{}\varphi _m^d\psi _m^a\right)\widehat{Z}_{cb}& \left(\underset{m}{}\varphi _m^b\psi _m^c\right)\widehat{Z}_{ad}.\hfill \end{array}$$ (2.32) The algebra of vector fields $`\widehat{Z}_{jl}`$ is finite dimensional (its dimension is $`d=\left(p_2p_1\right)\times \left(q_2q_1\right)`$). However, its isomorphism class cannot be determined without specifying the functions $`\varphi _m^l`$ and $`\psi _n^j`$, i.e. the matrices $`H`$ and $`K`$ in (1.1). In all examples in the literature, we have either $`d=1`$, or $`d=0`$. It is however easy to invent examples in which $`\left\{\widehat{Z}_{jl}\right\}`$ is simple, semisimple, solvable, or whatever we postulate a priori. The overall structure of the obtained Lie algebra is $$\left(\left\{\widehat{X}\right\}\left\{\widehat{Y}\right\}\right)+\left(\left\{\widehat{Z}\right\}+\left(\widehat{U}\widehat{V}\right)\right).$$ (2.33) If $`\left\{\widehat{Z}\right\}`$ is solvable, then (2.33) amounts to a Levi decomposition, since both $`\left\{\widehat{X}\right\}`$ and $`\left\{\widehat{Y}\right\}`$ are centerless Virasoro algebras and hence simple. We recall that the Levi theorem does not hold for infinite-dimensional Lie algebras and a Levi decomposition does not necessarily exist. Let us sum up the general results obtained so far for the symmetries of the generalized $`\mathrm{}`$Toda field theories (1.1) under the constraints imposed in Theorem 1. 1. The theory is always conformally invariant, since the inhomogeneous equation (2.24) always has a solution. 2. The theory allows gauge transformations $`\widehat{U}`$ and $`\widehat{V}`$ if $`p_2p_11`$. 3. The transformations of type $`\widehat{Z}`$ exist if $`\left(p_2p_1\right)\left(q_2q_1\right)1.`$ ### 2.2 <br>Special cases 1. The Mikhailov-Fordy-Gibbons two dimensional $`\mathrm{}`$-Toda system (1.2) We have $$h_1\left(n\right)=h_0\left(n\right)=1,\mathrm{and}k_1\left(n\right)=k_0\left(n\right)=1,$$ (2.34) so $`p_2p_1=q_2q_1=1`$. From eqs. (2.18) and (2.21) we have $$\psi _m=\varphi _m=1.$$ Equations (2.23) and (2.24) in this case imply $$\beta _n=\beta (x)+n\xi _x,\gamma _n=\gamma (y)+n\eta _y.$$ From Theorem 1 we now obtain all symmetries of eq (1.2), namely $`\widehat{X}(\xi )=\xi (x)_x+\xi _x{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}n_{u_n},\widehat{Y}(\eta )`$ $`=`$ $`\eta (y)_y+\eta _y{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}n_{u_n},`$ $`\widehat{U}=\beta (x){\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}_{u_n},\widehat{V}`$ $`=`$ $`\gamma (y){\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}_{u_n},`$ (2.35) $`\widehat{Z}=\left({\displaystyle \underset{m=\mathrm{}}{\overset{\mathrm{}}{}}}u_m\right)\left({\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}_{u_n}\right).`$ The generators $`\widehat{X},\widehat{Y},\widehat{U}`$ and $`\widehat{V}`$ were obtained in ref. using the so called “intrinsic method”. The generator $`\widehat{Z}`$ was not obtained there and cannot be obtained by the intrinsic method. 2. The Toda field theory (1.3) We take $`H=I`$. Then equations (2.18), (2.21) and (2.24) in this case imply $$\beta _m=\xi _x,,\gamma _m=\eta _y,A_{nm}=0.$$ The theory is only conformally invariant $$\widehat{X}(\xi )=\xi (x)_x\xi _x\underset{n}{}_{u_n},\widehat{Y}(\eta )=\eta (y)_y\eta _y\underset{n}{}_{u_n}$$ (2.36) and no further symmetries are obtained. 3. The Toda field theories (1.4) We take $`K=I`$ and relation (2.21) implies $$A_{nm}=0.$$ The remaining equations (2.24) cannot be solved explicitely for general $`h_\sigma (m)`$, but as said above, we can easily deal with in the constant coefficients case. As an example, let us restrict to the case when $`H`$ is the $`A_{\mathrm{}}`$ Cartan matrix (This is the $`A_N`$ Cartan matrix for $`N\mathrm{}`$, where the limit is taken symmetrically from a fixed, but not extremal, vertex in the corresponding Dynkin diagram). Thus we have $$h_1=h_{+1}=1,h_0=2,$$ (2.37) the solutions (2.23) become $$\beta _n=\frac{n^2}{2}\xi _x+nr_2(x)+r_1(x),\gamma _n=\frac{n^2}{2}\eta _y+ns_2(y)+s_1(y).$$ (2.38) The symmetry algebra is $`\widehat{X}(\xi )`$ $`=`$ $`\xi (x)_x+{\displaystyle \frac{1}{2}}\xi _x{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}n^2_{u_n},\widehat{Y}(\eta )=\eta (y)_y+{\displaystyle \frac{1}{2}}\eta _y{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}n^2_{u_n},`$ $`\widehat{U}_1\left(r_1\right)`$ $`=`$ $`r_1\left(x\right){\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}_{u_n},\widehat{V}_1\left(s_1\right)=s_1\left(y\right){\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}_{u_n},`$ (2.39) $`\widehat{U}_2\left(r_2\right)`$ $`=`$ $`r_2\left(x\right){\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}n_{u_n},\widehat{V}_2\left(s_2\right)=s_2\left(y\right){\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}n_{u_n},`$ where $`\xi (x)`$, $`\eta (y)`$, $`r_1(x)`$, $`s_1(y)`$, $`r_2(x)`$ and $`s_2(y)`$ are arbitrary smooth functions. ## 3 Symmetries of Finite Generalized Toda Field Theories ### 3.1 <br>General Results In this case we have a system of $`N`$ partial differential equations in $`N`$ fields $`u_n(x,y)`$, namely $$u_{n,xy}=F_n,F_n=\underset{m=1}{\overset{M}{}}K_{nm}\mathrm{exp}\left(\underset{l=1}{\overset{N}{}}H_{ml}u_l\right)\left(1nN\right).$$ (3.1) The “coupling constant” matrices $`H`$ and $`K`$ satisfy $`H^{M\times N}`$ and $`K^{N\times M}`$. The system (3.1) could arise in a quite general field theory with Lagrangian $$=\frac{1}{2}\underset{m,n=1}{\overset{N}{}}\kappa _{mn}_xu_m_yu_n\underset{m=1}{\overset{M}{}}c_m\mathrm{exp}\left(\underset{l=1}{\overset{N}{}}H_{ml}u_l\right)\left(c_m0\right),$$ (3.2) with $$K=L^1H^TC,L=\frac{\kappa +\kappa ^T}{2},C=\mathrm{diag}(c_1,\mathrm{},c_N).$$ (3.3) Some general considerations concerning the system (3.1) are in order. First, if either $`K`$, or $`H`$ (or both) allow an inverse, or at least a left inverse, then this system can be simplified. Indeed, let $`K^1`$ exist. We put $`u_n=_mK_{nm}\rho _m`$ and obtain $$\rho _{m,xy}=\mathrm{exp}\left(\underset{l=1}{\overset{M}{}}\left(HK\right)_{ml}\rho _l\right),1mM.$$ (3.4) Conversely, let $`H^1`$ exist and put $`w_j=_lH_{jl}u_l`$, we obtain $$w_{m,xy}=\underset{j=1}{\overset{M}{}}\left(HK\right)_{mj}e^{w_j},1mM.$$ (3.5) In other words, one of the matrices $`H`$, or $`K`$ can be normalized to $`I_M`$, if it is left invertible. The second comment is that the system (3.1) with $`K=I`$ admits Lie-Bäcklund transformations, and in this sense is completely integrable, if the matrix $`H`$ is a Cartan, or a generalized Cartan matrix . We mention that in the case of the infinite Toda field theories the matrices $`H`$ and $`K`$ in general have nontrivial kernels, are hence not invertible and we cannot normalize them. Let us now turn to the Lie point symmetries of the system (3.1). We write a general element of the symmetry algebra in the form (2.2) (with the sum in the range $`1nN`$), apply its prolongation to eq. (3.1) as in eq. (2.4). From the determining equations we find that for any $`F_n`$ in eq. (3.1), in complete analogy with the $`\mathrm{}`$-Toda theory, a general element of the symmetry algebra will have the form (2.5), the summation being from $`1`$ to $`N`$. Two determining equations remain and they depend on the specific form of $`F_n`$ in eq. (3.1). Making use of the fact that all the exponentials are linearly independent (no two rows in $`H`$ coincide) and that the matrix $`K`$ has no zero column, we reduce the remaining determining equations to two matrix relations $$HA=0,$$ (3.6) $$\left[(A(\xi _x+\eta _y)I)K\right]_{nm}=K_{nm}\left(HB\right)_m(1nN,\mathrm{\hspace{0.33em}\hspace{0.33em}1}mM).$$ (3.7) We multiply eq. (3.7) by $`H`$ from the left and use (3.6) to obtain $$\left(\xi _x+\eta _y\right)\left(HK\right)_{km}=\left(HK\right)_{km}\left(HB\right)_mk,m.$$ (3.8) If the matrix $`HK`$ has no zero column, then we obtain $$HB=\left(\xi _x+\eta _y\right)\overline{\mathrm{𝟏}}_M,$$ (3.9) where $`\overline{\mathrm{𝟏}}_M=(1,\mathrm{},1)^T^M`$, and from eq. (3.7) $$AK=0.$$ (3.10) Thus, matrix $`A`$ must satisfy the same two homogeneous equations (3.6) and (3.10) as in the infinite case. Furthermore, if $`\overline{\mathrm{𝟏}}_M`$ is in the image of $`H`$, then we define $`𝐛_N^N`$ to be an arbitrarily chosen (but specified) solution of the inhomogeneous equation $$H𝐛_N=\overline{\mathrm{𝟏}}_M.$$ (3.11) The results of these considerations can be summed up as follows ###### Theorem 2 Consider the generalized Toda field theories (3.1) with a finite number of fields $`N`$. Assume that all rows in $`H`$ are different and that the matrix $`HK`$ has no zero column. Then 3 types of symmetries can occur and they depend on the properties of the fundamental spaces of the matrices $`H`$ and $`K`$. The symmetries are of the same form as in Theorem 1, except that all summations range from 1 to $`N`$. However, if $`\overline{\mathrm{𝟏}}_MIm(H)`$, then $`\xi `$ and $`\eta `$ are arbitrary functions of $`x`$ and $`y`$, respectively, and the theory is conformally invariant. The quantities $`b_n`$ are the components of the vector $`𝐛_N`$, itself an arbitrary solution of eq. (3.11). Otherwise, if $`\overline{\mathrm{𝟏}}_MIm(H)`$, the theory is invariant only under the Poincaré group, generated by $$\widehat{P}_1=_x,\widehat{P}_2=_y,\widehat{L}=x_xy_y.$$ (3.12) Gauge transformations exist only if $`H`$ is not invertible. Analogously to the formulas (2.26), $`r_j`$ and $`s_j`$ are arbitrary functions and the vectors $`\psi ^j`$ span $`Ker\left(H\right)`$. Finally, the vectors $`\varphi ^l`$ span the left kernel of $`K`$. If this space is not zero, then $`\mathrm{dim}\left(Ker\left(K^T\right)\right)\times \mathrm{dim}\left(Ker\left(H\right)\right)`$ symmetries of the form (2.27) are admitted. From Theorem 2, contrary to the case of infinitely many fields, conformal invariance is not a priori guaranteed, but it imposes restrictions on the image of $`H`$. Gauge symmetries exist only if the matrix $`H`$ has a nonzero kernel. ### 3.2 <br>Special cases 1. The Mikhailov-Fordy-Gibbons Toda theory and generalizations Consider the field equation $$𝐔_{xy}=\frac{\mu ^2}{\beta }\underset{i=1}{\overset{N}{}}\frac{\alpha _i}{\alpha _i^2}\mathrm{exp}\left(\beta \alpha _i𝐔\right),$$ (3.13) where $`𝐔=(u_1,\mathrm{},u_N)`$ is an N-ple of real fields and $`(\alpha _1,\mathrm{},\alpha _N)`$ denote the simple roots of a classical simple finite Lie algebra. Equations (3.13) above take the form (1.2) for all $`n`$ satisfying $`N_0nN1`$. For $`n=N`$ we obtain $$u_{N,xy}=\mathrm{exp}\left(u_{N1}u_N\right).$$ (3.14) The equations for $`1n<N_0`$ are different for each Cartan series. The number $`N_0`$ is equal to 2 for $`A_N,B_N,C_N`$, and 3 for $`D_N`$. For the $`A_N`$ algebra we have $$u_{1,xy}=\mathrm{exp}\left(u_1u_2\right).$$ (3.15) Conformal and gauge transformations are exactly the same as given in eq. (2.35) (except that the summations are from $`1`$ to $`N`$). For the $`B_N`$ algebra we have $$u_{1,xy}=\mathrm{exp}\left(u_1\right)\mathrm{exp}\left(u_1u_2\right).$$ (3.16) Conformal transformations are as in eq. (2.35) (with the same comment about the summations) and there is no gauge invariance. For the $`C_N`$ algebra we have $$u_{1,xy}=\mathrm{exp}\left(u_1u_2\right)+2\mathrm{exp}\left(2u_1\right).$$ (3.17) The only symmetry is conformal invariance, generated by $`\widehat{X}(\xi )`$ $`=`$ $`\xi (x)_x+\xi _x{\displaystyle \underset{n=1}{\overset{N}{}}}\left(n{\displaystyle \frac{1}{2}}\right)_{u_n},`$ $`\widehat{Y}(\eta )`$ $`=`$ $`\eta (y)_y+\eta _y{\displaystyle \underset{n=1}{\overset{N}{}}}\left(n{\displaystyle \frac{1}{2}}\right)_{u_n}`$ (3.18) Finally, for the $`D_N`$ algebra we have $`u_{1,xy}`$ $`=`$ $`\mathrm{exp}\left(u_1u_2\right)\mathrm{exp}\left(u_1u_2\right),`$ $`u_{2,xy}`$ $`=`$ $`\mathrm{exp}\left(u_1u_2\right)+\mathrm{exp}\left(u_1u_2\right)\mathrm{exp}\left(u_2u_3\right).`$ (3.19) Again, the only symmetry is conformal invariance, in this case generated by $`\widehat{X}(\xi )`$ $`=`$ $`\xi (x)_x+\xi _x{\displaystyle \underset{n=1}{\overset{N}{}}}\left(n1\right)_{u_n},`$ $`\widehat{Y}(\eta )`$ $`=`$ $`\eta (y)_y+\eta _y{\displaystyle \underset{n=1}{\overset{N}{}}}\left(n1\right)_{u_n}.`$ (3.20) We mention that the infinite system (1.2) can also be reduced to the finite one by imposing periodicity $`u_{N+1}=u_1`$. In this case $`\overline{\mathrm{𝟏}}_N`$ is not contained in $`Im(H)`$ and there is no conformal invariance. Thus, the symmetry is given by the two dimensional Poincaré algebra (3.12) and by the gauge generators given in (2.35). 2. The Toda field theory (1.3) The symmetries are the same in the finite case as in the infinite one, namely the conformal transformations generated by (2.36) (for any finite matrix $`k`$). 3. The finite Toda theories (1.4) Since the Cartan matrix $`H`$ is invertible, this theory is equivalent to that described by eq. (1.3) in the sense of eqs. (3.4) and (3.5). Hence this theory is always and only conformally invariant. However, the generators of the vector fields take a slightly different form, which we report for a subsequent discussion. For the $`A_N`$ algebra the generators are given by $$\widehat{W}=\xi (x)_x+\eta (y)_y+\frac{1}{2}\left(\xi _x+\eta _y\right)\underset{n=1}{\overset{N}{}}n\left(nN1\right)_{u_n}.$$ (3.21) For the $`B_N`$ algebra, the symmetry generator is given by $`\widehat{W}=\xi (x)_x+\eta (y)_y{\displaystyle \frac{1}{4}}(\xi _x+\eta _y)\times `$ (3.22) $`\left\{N\left(N+1\right)_{u_1}+2{\displaystyle \underset{n=2}{\overset{N}{}}}\left[N\left(N+1\right)n\left(n1\right)\right]_{u_n}\right\}.`$ (3.23) For the $`C_N`$ algebra, the symmetry generator is given by $$\widehat{W}=\xi (x)_x+\eta (y)_y+\frac{1}{2}\left(\xi _x+\eta _y\right)\underset{n=1}{\overset{N}{}}\left[n\left(n2\right)N^2+1\right]_{u_n}.$$ (3.24) Finally, for the $`D_N`$ algebra ($`N4`$), one has $`\widehat{W}=\xi (x)_x+\eta (y)_y{\displaystyle \frac{1}{4}}(\xi _x+\eta _y)\times `$ (3.25) $`\left\{N\left(N1\right)\left(_{u_1}+_{u_2}\right)+2{\displaystyle \underset{n=3}{\overset{N}{}}}\left[N\left(N1\right)\left(n2\right)\left(n1\right)\right]_{u_n}\right\}.`$ (3.26) ## 4 Symmetries of Generalized Semi-Infinite Toda Field Theories ### 4.1 <br>General Results Let us now restrict the range of the discrete variable $`n`$ to be $`1n<\mathrm{}`$. Both the equations (1.1) of the generalized Toda field theories, and their symmetries will be modified. The matrices $`H`$ and $`K`$ will no longer be pure band matrices but will have the form $$H=\left(\begin{array}{cccccc}H_{1,1}\hfill & \begin{array}{ccc}\mathrm{}\hfill & \mathrm{}\hfill & \mathrm{}\hfill \end{array}\hfill & H_{1,N}\hfill & & & \hfill \\ \mathrm{}\hfill & \begin{array}{ccc}\mathrm{}\hfill & \mathrm{}\hfill & \mathrm{}\hfill \end{array}\hfill & \mathrm{}\hfill & & & \hfill \\ H_{M,1}\hfill & \begin{array}{ccc}\mathrm{}\hfill & \mathrm{}\hfill & \mathrm{}\hfill \end{array}\hfill & H_{M,N}\hfill & & & \hfill \\ & H_{M+1,M+1+p_1}\hfill & \begin{array}{ccc}\mathrm{}\hfill & \mathrm{}\hfill & \mathrm{}\hfill \end{array}\hfill & \mathrm{}\hfill & H_{M+1,M+1+p_2}\hfill & \hfill \\ & & H_{M+2,M+2+p_1}\hfill & \mathrm{}\hfill & \begin{array}{cccc}\mathrm{}\hfill & \mathrm{}\hfill & \mathrm{}\hfill & \mathrm{}\hfill \end{array}\hfill & \begin{array}{cc}H_{M+2,M+2+p_2}\hfill & \end{array}\hfill \\ & & & \mathrm{}\hfill & \begin{array}{cccc}\mathrm{}\hfill & & & \mathrm{}\hfill \end{array}\hfill & \begin{array}{ccccc}\mathrm{}\hfill & & & & \mathrm{}\hfill \end{array}\hfill \end{array}\right),$$ (4.1) where $`M+p_1NM+p_2`$ and the void entries are equal to zero. Similarly, the matrix $`K`$ takes the form $$K=\left(\begin{array}{cccccc}K_{1,1}\hfill & \mathrm{}\hfill & K_{1,N^{}}\hfill & & & \\ \mathrm{}\hfill & \mathrm{}\hfill & \mathrm{}\hfill & K_{N^{}+1+q_1,N^{}+1}\hfill & & \\ K_{M^{},1}\hfill & \mathrm{}\hfill & K_{M^{},N^{}}\hfill & \mathrm{}\hfill & K_{N^{}+2+q_1,N^{}+2}\hfill & \\ & & & \mathrm{}\hfill & \mathrm{}\hfill & \mathrm{}\hfill \\ & & & K_{N^{}+1+q_2,N^{}+1}\hfill & \mathrm{}\hfill & \mathrm{}\hfill \\ & & & & K_{N^{}+2+q_2,N^{}+2}\hfill & \mathrm{}\hfill \\ & & & & & \mathrm{}\hfill \end{array}\right)$$ (4.2) where $`N^{}+q_1M^{}N^{}+q_2`$. Although one could easily construct non trivial models, which do not fit in the given scheme, they seem quite artificial and, moreover, all the cases which we found in the literature satisfy the above restrictions. We denote by $`\stackrel{~}{H}`$ and $`\stackrel{~}{K}`$ respectively, the $`M\times N`$ and $`M^{}\times N^{}`$ matrices, which can be extracted by taking the first $`M`$ rows and the first $`N`$ columns from $`H`$ and, in turn, the first $`M^{}`$ rows and the first $`N^{}`$ columns from $`K`$. The symmetry algebra of the semi-infinite Toda field theory equation can either be obtained directly, ab initio, or we can obtain it from the infinite case of Section 2, by adding appropriate boundary conditions and requiring that they be invariant. As above, the functions $`\xi \left(x\right)`$, $`\eta \left(y\right)`$, $`A_{mn}`$ and $`B_n(x,y)`$ must satisfy the remaining determining equations (2.7) - (2.9). Following the same reasoning as in the finite case (see Section 3), we obtain the analogs of all the relations (3.6) - (3.10), where now all the labels and summations range from 1 to $`\mathrm{}`$ (i.e. we take $`N\mathrm{}`$ in all formulas). The key equation of the discussion is eq. (3.9) and its associated homogeneous system. Here, we separate the problem into the finite subsystems $`\stackrel{~}{H}\stackrel{~}{𝐁}`$ $`=`$ $`0,`$ (4.3) $`\stackrel{~}{H}\stackrel{~}{𝐁}`$ $`=`$ $`\left(\xi _x+\eta _y\right)\overline{\mathrm{𝟏}}_M,`$ (4.4) where $`\stackrel{~}{𝐁}=(B_1,\mathrm{},B_N)`$, and a difference linear equation, which we can put again in the form (2.18), or (2.24) respectively, for $`nM+1`$. The eq. (4.3) has $`Ker\left(\stackrel{~}{H}\right)`$ as its solution space. On the other hand, the difference equation (2.18) has a $`\left(p_2p_1\right)`$-dimensional solution space, the elements of which have the form $$B_n=\underset{j=1}{\overset{p_2p_1}{}}\alpha _j\psi _n^j,nM+1+p_1,$$ (4.5) in terms of the basis $`\left\{\psi _n^j\right\}`$. Moreover, the difference eq. (2.18) has only the zero solution in the case $`p_1=p_2`$. But, because of the imposed restrictions on the form of $`H`$, in such a case the components of the vector $`\stackrel{~}{𝐁}`$ are decoupled from the remaining $`(B_{N+1},\mathrm{})`$. This means that the semi-infinite homogeneous linear system $`HB=0`$ has zero-dimensional kernel only if both the finite system (4.3) and the homogeneous difference eq. (2.18) do. Assuming now that $`p_1<p_2`$ and, moreover, that $`M+p_1+1N`$, the components $`(B_{M+1+p_1},\mathrm{},B_N)`$ have to satisfy both the finite linear eq. (4.3) and the difference eq. (2.18). Substituting the representation (4.5) into (4.3), we get $`N\mathrm{dim}\left(Ker\left(\stackrel{~}{H}\right)\right)`$ constraints on the $`\left\{\alpha _i\right\}_{i=1,\mathrm{},p_2p_1}`$. Thus, if it results that $$MN+p_2+dim\left(Ker\left(\stackrel{~}{H}\right)\right)=n_0>0,$$ (4.6) then the semi-infinite homogeneous system $`HB=0`$ admits a $`n_0`$-dimensional kernel, spanned by the set of linearly independent functions $`\left\{\chi _n^j\right\}_{j=1,\mathrm{},n_0}`$. The above result implies that, if the constraint (4.6) holds, then the semi-infinite Toda model defined (4.1) and (4.2) possesses a symmetry group of gauge transformations, generated by the $`2\times n_0`$ vector fields $$\widehat{U}_j\left(r_j\right)=r_j\left(x\right)\underset{n=1}{\overset{\mathrm{}}{}}\chi _n^j_{u_n},\widehat{V}_j\left(s_j\right)=s_j\left(y\right)\underset{n=1}{\overset{\mathrm{}}{}}\chi _n^j_{u_n}\left(j=1,\mathrm{},p_2p_1\right).$$ (4.7) As in the finite case, a semi-infinite theory is conformally invariant if the inhomogeneous eq. (3.9) ( for semi-infinite matrices) has a solution. Thus, now we must require that the vector $`\overline{\mathrm{𝟏}}=(1,1,\mathrm{})`$ be contained in $`Im\left(H\right)`$. But, as outlined above, the problem is reduced to finding a solution of the eq. (4.4) and of the difference eq. (2.24). The former equation is solved if $$\overline{\mathrm{𝟏}}_MIm\left(\stackrel{~}{H}\right).$$ (4.8) For the difference eq. (2.24) a solution always exists as seen in Sec. 2. Hence the structure of the matrix $`H`$ shown in (4.1) garantees that a solution of the total inhomogeneous system always exists, once eq. (4.8) is satisfied ici. In conclusion, the condition (4.8) is not only necessary, but also sufficient to ensure the conformal invariance of the given Toda theories. Finally, an analysis similar to the study of the gauge invariance can be performed for the $`\widehat{Z}`$-type transformations, which exist if a common solution of the two semi-infinite homogeneous systems $$HA=0,AK=0$$ (4.9) can be found. Thus, we are lead to the following theorem ###### Theorem 3 Consider the semi-infinite Toda field theory (1.1), with $`H`$ and $`K`$ given by (4.1) and (4.2), respectively, and with all rows of $`H`$ different. Moreover, let $`HK`$ have no zero columns. Then, the symmetry algebra depends on the fundamental spaces of the finite dimensional submatrices $`\stackrel{~}{H}`$ and $`\stackrel{~}{K}`$, on the solutions of the difference eqs. (2.18) and (2.24) for $`nM+1`$ and, finally, on the solutions of the difference eq. (2.21) for $`mN^{}+1`$. The theory is conformally invariant if the condition (4.8) holds. The corresponding generators take the form (2.25). Otherwise, if (4.8) does not hold, the symmetry reduces to the Poincaré group generated by (3.12). A gauge transformation group, involving $`2n_0`$ arbitrary functions of one variable, exists if the relation (4.6) holds. The algebra generators take the form (4.7). Finally, $`\widehat{Z}`$-type gauge transformations exist if not only (4.6) holds, but also the supplementary condition $$N^{}M^{}+q_2+\mathrm{dim}\left(Ker\left(\stackrel{~}{K}^T\right)\right)=m_0>0$$ (4.10) is satisfied. In such a case they form a Lie algebra of dimension $`m_0\times n_0`$. ### 4.2 <br>Special cases Now let us consider the same three examples as in the previous Sections. 1. Mikhailov-Fordy-Gibbons field theories All examples of Section 3.2 can be generalized to the semi-infinite case, simply allowing $`N`$ to go to $`\mathrm{}`$ for each classical Lie algebra. The equations labeled by $`1nN_0`$ are explicitly given by (3.15), (3.16), (3.17) and (3.19), respectively. Moreover, for $`iN_0`$ the equations are the same as in the infinite case, i.e. eq. (1.2). For the $`A_\mathrm{}+`$ algebra (We use this notation in order to distinguish this semi-infinite model from the previously introduced $`A_{\mathrm{}}`$ infinite one) we have $`M=N=M^{}=N^{}=0`$ and hence the symmetries are exactly the same as in the infinite and in the finite cases (see eq. (2.35)), where the summations are over the appropriate range. For the $`B_{\mathrm{}}`$ algebra one has $`\stackrel{~}{H}=\stackrel{~}{K}=\left(1\right)`$, then also $`M=N=M^{}=N^{}=1`$, as one can see from (3.16). Theorem 3 allows to establish that there are no gauge transformations of any kind and the generators of the conformal transformations are the same as given in (2.35). From eq. (3.17) one sees that $`\stackrel{~}{H}=\stackrel{~}{K}=\left(2\right)`$ for the $`C_{\mathrm{}}`$ algebra, then $`M=N=M^{}=N^{}=1`$. Thus, Theorem 3 establishes that only the conformal invariance is admitted. Its generators have the same form as in eq. (3.18), where the summation is over the positive integers. Finally, for the $`D_{\mathrm{}}`$ algebra one has $$\stackrel{~}{H}=\left(\begin{array}{cc}1\hfill & 1\hfill \\ 1\hfill & 1\hfill \end{array}\right)=\stackrel{~}{K}^T.$$ Theorem 3 implies that only conformal transformations leave the system invariant and their generators are obtained by taking the limit $`N\mathrm{}`$ in the formulas (3.20). 2. The semi-infinite Toda field theory (1.3) The discussion is very simple. Indeed, since $`H`$ is the indentity matrix, there are no gauge transformations. Moreover, the generators of the conformal transformations in the infinite, semi-infinite and finite cases take always the same form (2.36), where the summations are over the appropriate range. 3. The semi-infinite Toda field theories (1.4) As opposed to the finite case, the matrix $`H`$ is no longer invertible, so now these theories are not equivalent to the ones given by (1.3). First, we observe that, since $`K`$ is the identity matrix, there are no $`\widehat{Z}`$-type transformations. For any classical Lie algebra, extended to $`N\mathrm{}`$, the recursive part of the systems, i.e. the equations labeled by $`nN_0`$ as defined in Sec. 3.2, are always the same as in the infinite case discussed in Sec. 2.2.3. The solution of the corresponding difference equations for $`B_n\left(nN_0\right)`$ , that is (2.18) and (2.24), are the same as in (2.38) and the generators are as in (2.39). However, for $`1n<N_0`$ the equations provide constraints of the form (4.3) and (4.4). The application of the Theorem 3 implies * All the semi-infinite systems (1.4) are conformally invariant. * All the semi-infinite systems (1.4) have $`n_0=1`$, as defined in (4.6), hence a gauge transformation algebra of the form (4.7) exists, with $`j=1`$. In the $`A_\mathrm{}+`$ case the $`\widehat{X}`$ and $`\widehat{Y}`$ conformal symmetries survive as in eq. (2.39), and so do $`\widehat{U}_2`$ and $`\widehat{V}_2`$ do. However the symmetries $`\widehat{U}_1`$ and $`\widehat{V}_1`$ are no longer present. In the $`B_{\mathrm{}}`$ case the generators $`\widehat{X},\widehat{Y}`$ and $`\widehat{U}_2,\widehat{V}_2`$ combine together to give the new conformal symmetry generators $$\widehat{X}=\xi \left(x\right)_x+\frac{1}{2}\xi _x\underset{n=1}{\overset{\mathrm{}}{}}n\left(n1\right)_{u_n},\widehat{Y}=\eta \left(y\right)_y+\frac{1}{2}\eta _y\underset{n=1}{\overset{\mathrm{}}{}}n\left(n1\right)_{u_n}.$$ (4.11) The remaining gauge invariance is generated by $$\widehat{U}\left(r\right)=r\left(x\right)\left[_{u_1}+2\underset{n=2}{\overset{\mathrm{}}{}}_{u_n}\right],\widehat{V}\left(s\right)=s\left(y\right)\left[_{u_1}+2\underset{n=2}{\overset{\mathrm{}}{}}_{u_n}\right].$$ (4.12) For the $`C_{\mathrm{}}`$ algebra the symmetry algebra is $$\begin{array}{c}\widehat{X}=\xi \left(x\right)_x+\frac{1}{2}\xi _x\underset{n=1}{\overset{\mathrm{}}{}}n\left(n2\right)_{u_n},\hfill \\ \widehat{Y}=\eta \left(y\right)_y+\frac{1}{2}\eta _y\underset{n=1}{\overset{\mathrm{}}{}}n\left(n2\right)_{u_n},\hfill \\ \widehat{U}\left(r\right)=r\left(x\right)\underset{n=1}{\overset{\mathrm{}}{}}_{u_n},\widehat{V}\left(s\right)=s\left(y\right)\underset{n=1}{\overset{\mathrm{}}{}}_{u_n}.\hfill \end{array}$$ (4.13) Finally, for the $`D_{\mathrm{}}`$ algebra one has $$\begin{array}{c}\widehat{X}=\xi \left(x\right)_x+\frac{1}{2}\xi _x\underset{n=1}{\overset{\mathrm{}}{}}\left(n1\right)\left(n2\right)_{u_n},\hfill \\ \\ \widehat{Y}=\eta \left(y\right)_y+\frac{1}{2}\eta _y\underset{n=1}{\overset{\mathrm{}}{}}\left(n1\right)\left(n2\right)_{u_n},\hfill \\ \\ \widehat{U}\left(r\right)=r\left(x\right)\left[_{u_1}+_{u_2}+2\underset{n=3}{\overset{\mathrm{}}{}}_{u_n}\right],\hfill \\ \\ \widehat{V}\left(s\right)=s\left(y\right)\left[_{u_1}+_{u_2}+2\underset{n=3}{\overset{\mathrm{}}{}}_{u_n}\right].\hfill \end{array}$$ (4.14) The formulas for the semi-infinite models (1.4) are consistent with those obtained in the finite case in Sec. 3.2.3. The generators of the conformal invariance, in each case, are simply obtainable by dropping all terms involving $`N`$. Conversely, the terms proportional to a power of $`N`$ provide us with the gauge invariance generators in the semi-infinite extensions. In this limit, the functions $`r=\xi _x`$ and $`s=\eta _y`$ must be considered as new linearly independent functions. ## 5 Conclusions We have introduced the generalized Toda system (1.1) and investigated its Lie point symmetry group. It turned out that in the infinite case $`\left(\mathrm{}<n<\mathrm{}\right)`$ these systems are always invariant under an infinite dimensional group of conformal transformations. It is also gauge invariant, if a certain homogeneous linear difference equation (i.e. eq. (2.18)) has non trivial solutions. Further gauge transformations exist if another linear homogeneous difference equation (i.e. eq. (2.21)) also has nontrivial solutions. If we restrict the range of $`n`$ to $`1n<\mathrm{}`$, in some cases the symmetry group remains the same, or is reduced to a subgroup of the original symmetry group. However, in other cases (see (4.12) and (4.14)) the symmetry group does not coincide with a Lie subgroup. In the finite case, with $`1nN`$, the symmetry group remains the same as in the semi-infinite case, or it is reduced further. In some situations (see Theorem 2 and 3) the infinite dimensional conformal symmetry group is reduced to the Poincaré group in two dimensions (see eq. (3.12)). These results were obtained directly, that is by analyzing the determining equations for the symmetries for all types of systems: infinite, semi-infinite and finite. The question to which we plan to devote a separate article is the application of the infinite generalized Toda systems. In particular we will establish the degree to which the symmetries of the semi-infinite and finite Toda systems are “inherited” from those of the infinite systems. In other words we plan to discuss symmetry breaking by boundary or periodicity conditions of the infinite chains. One of the surprising results obtained in the present work is that the class of the conformally invariant Toda field theories is much larger than the class of the completely integrable models. Indeed, the existence of a Lax pair imposes severe algebraic restrictions on the matrices $`H`$ and $`K`$ (see for instance ). ## <br>Acknowledgments This work is part of a project supported by the NATO CRG 960717, by the Italian INFN and by the project SINTESI of the Italian Ministry of the University and the Scientific Research. L.M. would like to thank the Centre de Recherches Mathématiques of the Université de Montréal for its warm hospitality. S.L. and P.W. would like to thank the Dipartimento di Fisica - Universitá di Lecce for its hospitality. The research of P.W. was partly supported by grants from NSERC and from FCAR. S.L. acknowledges a PhD scholarship from FCAR.
warning/0004/hep-ex0004015.html	ar5iv	text	# CURRENT STATUS OF THE K2K LONG-BASELINE NEUTRINO-OSCILLATION EXPERIMENTaafootnote aTalk at XXXVth Rencontres de Moriond “Electroweak interactions and unified theories”, Les Arcs, Savoie, France, March 11-18, 2000. The transparencies used in this talk can be found in http://neutrino.kek.jp/õyama/public.html ## 1 Introduction The K2K experiment$`^{\mathrm{?},\mathrm{?}}`$ is the first long-baseline neutrino-oscillation experiment using an artificial neutrino beam. Almost a pure $`\nu _\mu `$ beam from $`\pi ^+`$ decays is generated in the KEK 12-GeV/c Proton Synchrotron, and is detected in Super-Kamiokande (SK) 250km away. Neutrino oscillation can be examined from characteristics of the neutrino events observed in SK. The nominal sensitive region in the neutrino-oscillation parameters is $`\mathrm{\Delta }m^2>3\times 10^3`$eV<sup>2</sup>, which covers the parameter region suggested by the atmospheric neutrino anomaly observed by several underground experiments,$`^{\mathrm{?},\mathrm{?},\mathrm{?}}`$ and confirmed by SK.$`^\mathrm{?}`$ The sensitive regions on the neutrino oscillation parameters for $`\nu _\mu \nu _\tau `$ oscillation and $`\nu _e\nu _\mu `$ oscillation are shown in Figure 1. ## 2 Neutrino beam and detectors The K2K experiment consists of (1)a proton synchrotron and neutrino beam line including magnetic horns$`^\mathrm{?}`$ and various beam monitors, (2)two front detectors (1kt water Cherenkov detector (1KT) and so-called Fine Grained detector (FGD)) 300m downstream of the target and (3)SK as a far detector. The detailed design and performance of those components were already presented in Ref.2, and are not discussed in this report. Instead, the design of the neutrino beam line and the property of the neutrino beam are given in Table 1. The design, performance and event rate of the detectors are summarized in Table 2. A schematic view of the K2K front detector is also shown in Figure 2. ## 3 Summary of the data in 1999 The K2K experiment was successfully started in early 1999. The first neutrino beam was generated on January 27. After machine studies and beam tuning, the first physics run was started on March 3. In 1999, about 100 days of physics data-taking was scheduled. However, the successful physics data-taking was only for 39.4 days, due to some problems with the neutrino beam line. On June 19, the first neutrino event was observed in SK. In the 1999 run, the integrated proton intensity was $`7.20\times 10^{18}`$ p.o.t. (protons on target), which is about 7% of the goal of the experiment, $`10^{20}`$ p.o.t.$`^\mathrm{?}`$ The neutrino beam direction has been confirmed to agree with the direction of the SK detector within 0.3 mrad based on the data of the beam monitors.$`^\mathrm{?}`$ Because the absolute flux and energy spectrum of the neutrino beam are expected to be almost the same within 3 mrad, the adjustment of the neutrino beam direction is sufficient. The selection of neutrino events in SK employs the time difference between the neutrino beam and each event. Considering the neutrino beam duration (1.1$`\mu `$sec) and accuracy of the absolute time determination ($`<0.3\mu `$sec), events within a 1.5$`\mu `$sec time window covering the neutrino beam period are selected. A total of 12 events have been found, which includes events in the inner counter, in the outer counter and interactions in the surrounding rock. Among them, 3 events were detected in the fiducial volume of the inner counter. Because the expected atmospheric neutrino background in the fiducial volume within the neutrino beam period is calculated to be $`2\times 10^4`$ events, the 3 events in the fiducial volume are a clear signal of neutrinos from KEK. ## 4 Present status of data analysis Strategies concerning oscillation searches at K2K are summarized as follows. The $`\nu _\mu \nu _\tau `$ oscillation can be examined by a disappearance of neutrino events in SK because the energy of the neutrino beam is smaller than the $`\tau `$ production threshold. To recognize a reduction of neutrino events efficiently, expected event number without oscillation must be accurately estimated using the observed event numbers in the front detectors. In addition, the neutrino energy spectrum in SK should be distorted in the case of oscillation because the oscillation probability depends on the energy of the neutrinos. Therefore, the expected neutrino energy spectrum for no oscillation must be calculated from an extrapolation of the spectrum in the front detectors. An examination of the $`\nu _e\nu _\mu `$ oscillation is an appearance search. A possible excess of electron neutrino events in SK is direct evidence of the $`\nu _e\nu _\mu `$ oscillation, because the beam from KEK is almost pure muon neutrinos in the case of no oscillation, and because the particle identification capability in SK has already been proved to be excellent.$`^{\mathrm{?},\mathrm{?}}`$ In order to attain a better oscillation sensitivity, the fraction of electron neutrinos in the original neutrino beam, which is estimated to be about 1% from a Monte-Carlo simulation, must be experimentally measured as precisely as possible. Therefore, the $`\nu _e/\nu _\mu `$ ratio observed in the front detectors is a key point of the $`\nu _e\nu _\mu `$ oscillation analysis. The following three subsections discuss the present status of the front detector analysis along these three strategies, i.e. (1)absolute event number, (2)shape of the neutrino energy spectrum, and (3)$`\nu _e/\nu _\mu `$ ratio. At present, those studies are impossible in SK because the statistics is not sufficient. ### 4.1 absolute event numbers The expected event numbers in SK ($`N_{exp}^{SK}`$) are calculated based on the observed event numbers in the front detectors, and extrapolation of the front detector to SK. $`N_{exp}^{SK}`$ is written as $$N_{exp}^{SK}=\frac{N_{obs}^{FD}\times N_{cal}^{SK}}{N_{cal}^{FD}},$$ where $`N_{obs}^{FD}`$ is the observed event numbers in the front detectors; $`N_{cal}^{SK}`$ and $`N_{cal}^{FD}`$ are the calculated event numbers in SK and the front detectors, respectively, using the same simulation program. We employed neutrino interactions in 50.3 tons fiducial volume of 1KT for $`N_{obs}^{FD}`$ and $`N_{cal}^{FD}`$. In addition, neutrino interactions in the scintillating fiber tracker$`^\mathrm{?}`$ and the muon chamber were also used to examine the consistency between the observations in the front detectors. The results are summarized in Table 3. $`N_{obs}^{FD}/N_{cal}^{FD}`$, which shows an agreement of observed event number with the simulation, is found to be $`0.840.85`$ for three independent observations, and is consistent with each other. The expected event number in SK is calculated to be 12.3$`\genfrac{}{}{0pt}{}{+1.7}{1.9}`$ in 22.5ktons of the fiducial volume, and $``$31 events in the total volume. Although the observed event numbers, 3 in fiducial volume and 12 in the total volume, are considerably smaller than the expectations, nothing can be concluded about the neutrino oscillation at this stage because of poor statistics. ### 4.2 neutrino energy spectrum To determine the neutrino energy spectrum, quasi-elastic interactions of muon neutrinos, $`\nu _\mu N\mu N^{}`$, in the scintillating fiber tracker are employed. This is because most of the neutrino energy is transfered to the muons in quasi-elastic interactions, and the muon energy can be measured from the range in the muon chamber. The neutrino energy can be directly calculated from the muon energy with a small correction related to the scattering angle of the muon. It should also be noted that quasi-elastic scatterings are detected as single ring events in SK, and can be easily analyzed. The muon energy distribution for quasi-elastic interactions in the scintillating fiber tracker is shown in Fig.3-(a) along with expectations from a Monte-Carlo simulation. For a comparison, the expected neutrino energy spectrum in SK is shown in Fig.3-(b) together with the spectrum for two sets of oscillation parameters, $`\mathrm{\Delta }m^2=0.01`$eV<sup>2</sup> and $`\mathrm{\Delta }m^2=0.005`$eV<sup>2</sup>. The shape of the muon energy spectrum in Fig.3-(a) agrees with the spectrum from a Monte-Carlo simulation, and is somehow correlated with the expected neutrino energy distribution shown in Fig.3-(b). However, a calculation of neutrino energy spectrum from the muon energy distribution, and its extrapolation to SK is still being studied. ### 4.3 $`\nu _e/\nu _\mu `$ ratio The $`\nu _e/\nu _\mu `$ ratio in the front detectors are independently measured using 1KT and FGD. As reported in Ref.21, the e/$`\mu `$ identification analysis in water Cherenkov detectors employs a likelihood function which quantitatively evaluates the agreement of the Cherenkov ring patterns with electrons and with muons. The likelihood distribution for the neutrino beam obtained in the 1KT is shown in Fig.4-(a) together with an expectation from a Monte-Carlo simulation. The likelihood distribution for atmospheric neutrino interactions obtained by SK is shown in Fig.4-(b) for a comparison. In Fig.4-(b), two peaks corresponding to electron neutrinos and muon neutrinos can be clearly distinguished. On the other hand, in Fig.4-(a), most of the Cherenkov ring patterns are judged to be muons, and the agreement with the distribution for a Monte-Carlo simulation is excellent. Although the $`\nu _e/\nu _\mu `$ ratio experimentally obtained in 1KT is found to be as small as the Monte-Carlo expectation, the numerical results are still being analyzed. The identification of an electron neutrino event in the FGD uses the response of each detector component for electrons and muons. A typical muon neutrino event and an electron neutrino event are shown in Fig.5-(a) and Fig.5-(b), respectively. In Fig.5-(a), a muon is generated in the water target of the scintillating fiber tracker, and produces a clear and long single muon track. The muon penetrates the trigger counters and the lead-glass counters, and is stopped in the middle of the muon chamber. In Fig.5-(b), on the other hand, the electron produces an electromagnetic shower in the scintillating fiber tracker; three charged particles are generated as shown in the expanded view of Fig.5-(b). These charged particles lose all of their energies in the lead glass counters by producing a further electromagnetic shower. No particles escape from the lead-glass counters. Considering these characteristics, electron neutrino events are identified by (1)a large track multiplicity in the scintillating fiber tracker, (2)a large energy deposit in the lead glass counters, and (3)an absence of tracks in the muon chamber. Specially, the selection using the lead glass counters is efficient to separate muons and electrons because the energy deposit from a single muon in the lead-glasses is found to be 450$`\pm `$150 MeV from cosmic-ray muon data, and energy resolution for 1 GeV electron is about 10% from an electron beam test. At present, the selection criteria for the electromagnetic shower events is being tuned. Although a small $`\nu _e/\nu _\mu `$ ratio is indicated by both 1KT and FGD, the numerical results on the $`\nu _e/\nu _\mu `$ ratio are still being studied. ## 5 Summary and future prospect During 39.4 days of successful data-taking, a total intensity of $`7.20\times 10^{18}`$ protons on target were accumulated in 1999. In Super-Kamiokande, we observed 3 neutrino interactions in the fiducial volume, where the expectation is 12.3$`\genfrac{}{}{0pt}{}{+1.7}{1.9}`$. An oscillation analysis focusing on the absolute event number, the distortion of the neutrino energy spectrum, and the $`\nu _e/\nu _\mu `$ ratio is in progress. From January 2000 to March 2001, about 160 days of physics data taking are scheduled. If the data can be accumulated with 100% efficiency, we will obtain a total intensity of $`46\times 10^{18}`$ p.o.t., and about 70 events will be accumulated by end of March 2001. ## References
warning/0004/gr-qc0004061.html	ar5iv	text	# Entropies of Rotating Charged Black Holes from Conformal Field Theory at Killing Horizons ## I INTRODUCTION The statistical mechanical description of the Bekenstein-Hawking black hole entropy \- in terms of microscope states is an outstanding open question and much effort has been concentrated on the problem for some years -. Success seems to come with the paper of Strominger and Vafa which was followed by a host of others. It is well known since the work of Brown and Henneaux that a asymptotic symmetry group of AdS<sub>3</sub> is generated by a Virasoro algebra, and that therefore any consistent quantum theory of gravity on AdS<sub>3</sub> is conformal field theory. Using the result Strominger calculated the entropy of black holes whose near-horizon geometry is locally $`AdS_3`$ from the asymptotic growth of states. Precise numerical agreement with the Bekenstein-Hawking area formula for the entropy was found. In light of the work, one could statistically reinterpret the black hole entropy by establishing a relation conformal field theory on the boundary of related anti-de Sitter space. In order to overcome the limitations of Strominger’s method, such as the approach can only be used for 2+1 dimensional spacetime and it is based on an algebra of transformations at infinity, Carlip generalized Brown-Henneaux-Strominger’s approach by looking at the symmetries of the event horizon of an (n+1)-dimensional Schwarzschild-like black hole. This construction is valid for black hole in any dimension. In Ref. Carlip re-derived the central extension of the constraint algebra of general relativity by using manifestly covariant phase space methods - and a boundary which is a surface that look like a (local) Killing horizon. A natural set of boundary conditions leads to a Virasoro subalgebra with a calculable central charge. Then, by means of conformal field theory method, Carlip studied the statistical entropies of the Rindler space, static de Sitter space, Taub-NUT and Taub-Bolt spaces, and 2-dimensional dilaton gravity. However, at the moment the question whether or not the covariant phase space approach can be used for the stationary axisymmetric charged black holes which are described by solutions of the Einstein-Maxwell equations, such as the Kerr-Newman black hole and the Kerr-Newman-AdS black hole, still remains open. The aim of this paper is to settle the question. The paper is organized as follows: In Sec. II, by using the covariant phase techniques we extend Carlip’s investigation for vacuum case $`𝐋_{a_1a_2\mathrm{}a_n}=\frac{1}{16\pi G}ϵ_{a_1a_2\mathrm{}a_n}R`$ to a case including a cosmological term and electromagnetic fields, i.e., the Lagrangian n-form is described by $`𝐋_{a_1a_2\mathrm{}a_n}=\frac{1}{16\pi }ϵ_{a_1a_2\mathrm{}a_n}\left[\frac{1}{G}(R2\mathrm{\Lambda })+F^{ab}F_{ab}\right].`$ In Sec. III, the standard Virasoro subalgebras with corresponding central charges are constructed for the Kerr-Newman black hole and the Kerr-Newman-AdS black hole. The statistical entropies for these objects are then calculated by using Cardy formula. Some discussions and summaries are presented in the last section. ## II Algebra of diffeomorphism Lee, Wald, and Iyer showed that the variation of the Lagrangian defines the equation of motion n-form $`𝐄`$ and the symplectic potential (n-1)-form $`𝚯`$ via the equation $`\delta 𝐋=𝐄\delta \varphi +d𝚯`$ , where $`𝐋`$ is an n-form, $`𝐄\delta \varphi =𝐄_g^{ab}\delta g_{ab}+𝐄_\psi \delta \psi ,`$ $`\varphi =(g_{ab},\psi )`$ denotes an arbitrary collection of dynamical fields, and the equations of motion are taken to be $`𝐄_g^{ab}=0`$ and $`𝐄_\psi =0.`$ Let $`\xi ^a`$ be any smooth vector fields on the spacetime manifold $`𝐌`$, i. e., $`\xi ^a`$ is the infinitesimal generator of a diffeomorphism, we can define a Noether current (n-1)-form as $$𝐉[\xi ]=𝚯[\varphi ,_\xi \varphi ]\xi 𝐋,$$ (1) here and hereafter the “central dot” denotes the contraction of the vector field $`\xi ^a`$ into the first index of the differential form. By using the equations of motion a standard calculation shows that $`𝐉`$ is closed for all $`\xi ^a`$, i.e., $`d𝐉=0`$. Then we have $$𝐉=d𝐐,$$ (2) where $`𝐐`$ is a Noether charge (n-2)-form. From the variation of Noether current (n-1)-form, we know that the symplectic current (n-1)-form $`\omega [\varphi ,\delta _1\varphi ,\delta _2\varphi ]=\delta _2𝚯[\varphi ,\delta _1\varphi ]\delta _1𝚯[\varphi ,\delta _2\varphi ]`$ can be expressed as $$\omega [\varphi ,\delta \varphi ,_\xi \varphi ]=\delta 𝐉[\xi ]d(\xi 𝚯[\varphi ,\delta \varphi ]),$$ (3) and Hamilton’s equation of motion is given by $$\delta H[\xi ]=_C\omega [\varphi ,\delta \varphi ,_\xi \varphi ]=_C[\delta 𝐉[\xi ]d(\xi 𝚯[\varphi ,\delta \varphi ])].$$ (4) By using Eq. (2) and Carlip’s boundary conditions listed in Appendix A and defining a (n-1)-form $`𝐁`$ as $$\delta _C\xi 𝐁[\varphi ]=_C\xi 𝚯[\varphi .\delta \varphi ],$$ (5) the Hamiltonian can be expressed as $$H[\xi ]=_C(𝐐[\xi ]\xi 𝐁[\varphi ]).$$ (6) It is well-known that the Poisson bracket forms a standard “surface deformation algebra” $$\{H[\xi _1],H[\xi _2]\}=H[\{\xi _1,\xi _2\}]+K[\xi _1,\xi _2],$$ (7) where the central term $`K[\xi _1,\xi _2]`$ depends on the dynamical fields only through their boundary values. In this paper, we focus our attention to stationary axisymmetric charged black holes. So we take the Lagrangian n-form as $$𝐋_{a_1a_2\mathrm{}a_n}=\frac{1}{16\pi }ϵ_{a_1a_2\mathrm{}a_n}\left[\frac{1}{G}(R2\mathrm{\Lambda })+F^{ab}F_{ab}\right],$$ (8) where $`ϵ_{a_1a_2\mathrm{}a_n}`$ is a volume element (a continuous non-vanishing n-form), $`\mathrm{\Lambda }`$ is the cosmological constant, and $`F_{ab}`$ is the electromagnetic field strength tensor. By using the infinitesimal generator of a diffeomorphism, $`\xi `$, we know that the symplectic potential (n-1)-form is given by $$𝚯_{a_1a_2\mathrm{}a_{n1}}[g,_\xi g]=\frac{1}{4\pi }ϵ_{ca_1a_2\mathrm{}a_{n1}}\left\{\frac{1}{2G}(_e^{[e}\xi ^{c]}+R_e^c\xi ^e)+F^{dc}\left[F_{ed}\xi ^e+(\xi ^eA_e)_{;d}\right]\right\}.$$ (9) Eqs. (1) and (9) yields $`𝐉_{a_1a_2\mathrm{}a_{n1}}`$ $`=`$ $`{\displaystyle \frac{1}{8\pi G}}ϵ_{ca_1a_2\mathrm{}a_{n1}}[_e^{[e}\xi ^{c]}+(R_e^c{\displaystyle \frac{1}{2}}\delta _e^cR+\delta _e^c\mathrm{\Lambda })\xi ^e]`$ (11) $`{\displaystyle \frac{1}{4\pi }}ϵ_{ca_1a_2\mathrm{}a_{n1}}\left[{\displaystyle \frac{1}{4}}F^{bd}F_{bd}\delta _e^cF^{cd}F_{ed}\right]\xi ^e+{\displaystyle \frac{1}{4\pi }}ϵ_{ca_1a_2\mathrm{}a_{n1}}F^{ec}(\xi ^dA_d)_{;e}`$ $`=`$ $`{\displaystyle \frac{1}{4\pi }}ϵ_{ca_1a_2\mathrm{}a_{n1}}\left[{\displaystyle \frac{1}{2G}}_e^{[e}\xi ^{c]}+_e(^{[e}A^{c]}A_d\xi ^d)\right],`$ (12) in above calculation, we used the Einstein-Maxwell field equations in which the energy-momentum tensors is given by $`\frac{1}{4\pi }\left[\frac{1}{4}F^{bc}F_{bc}\delta _e^dF^{dc}F_{ec}\right]`$. From Eqs. (2) and (12) we have $$𝐐_{a_1a_2\mathrm{}a_{n2}}=\frac{1}{4\pi }ϵ_{bca_1\mathrm{}a_{n2}}\left[\frac{1}{4G}^b\xi ^c+(^bA^c)A_e\xi ^e\right].$$ (13) For a stationary axisymmetric charged black hole (such as the Kerr-Newman black hole and the Kerr-Newman AdS/dS black hole), the electromagnetic potential $`A_a`$, the electromagnetic field tensors $`F^{03}`$, and the Killing vector can be expressed respectively as $`A_a`$ $`=`$ $`(A_0,0,0,A_3)`$ (14) $`F^{03}`$ $`=`$ $`F^{30}=0.`$ (15) $`\chi _H^a`$ $`=`$ $`\chi _H^{(t)}+\chi _H^{(\phi )}=(1,0,0,\mathrm{\Omega }_H),`$ (16) where the vector $`\chi _H^{(t)}`$ correspond to time translation invariance, $`\chi _H^{(\phi )}`$ to rotational symmetry, and $`\mathrm{\Omega }_H=(g_{t\phi }/g_{\phi \phi })_H`$ is the angular velocity of the black hole. Using Eqs. (15), (16), (A6), and (A11) it is easy to show that $`{\displaystyle \frac{1}{4\pi }}ϵ_{bca_1\mathrm{}a_{n2}}(^bA^c)A_e\xi ^e0.\text{ at the horizon}`$ (17) Then, Eq. (13) is reduced to $$Q_{a_1a_2\mathrm{}a_{n2}}=\frac{1}{16\pi G}ϵ_{bca_1a_2\mathrm{}a_{n2}}^b\xi ^c.$$ (18) Denoting by $`\delta _\xi `$ the variation corresponding to diffeomorphism generated by $`\xi `$, for the Noether current $`𝐉[\xi ]`$ we have $$\delta _{\xi _2}𝐉[\xi _1]=\xi _2d𝐉[\xi _1]+d(\xi _2𝐉[\xi _1])=d[\xi _2(𝚯[\varphi ,_{\xi _1}\varphi ]\xi _1𝐋)].$$ (19) Substituting Eq. (19) into Eq. (4) and using Eq. (9) we get $`\delta _{\xi _2}H[\xi _1]`$ $`=`$ $`{\displaystyle _C}\left(\delta _{\xi _2}𝐉[\xi _1]d(\xi _1𝚯[\varphi ,\delta _{\xi _2}\varphi ])\right)`$ (20) $`=`$ $`{\displaystyle _C}\left(\xi _2𝚯[\varphi ,_{\xi _1}\varphi ]\xi _1𝚯[\varphi ,_{\xi _2}\varphi ]\xi _2\xi _1𝐋\right)`$ (21) $`=`$ $`{\displaystyle \frac{1}{16\pi G}}{\displaystyle _C}ϵ_{bca_1\mathrm{}a_{n2}}\left[\xi _2^b_d(^d\xi _1^c^c\xi _1^d)\xi _1^b_d(^d\xi _2^c^c\xi _2^d)\right]`$ (24) $`+{\displaystyle \frac{1}{8\pi }}{\displaystyle _C}ϵ_{bca_1\mathrm{}a_{n2}}\left\{\xi _2^bF^{dc}\left[F_{ed}\xi _1^e+(\xi _1^eA_e)_{;d}\right]\xi _1^bF^{dc}\left[F_{ed}\xi _2^e+(\xi _2^eA_e)_{;d}\right]\right\}`$ $`{\displaystyle \frac{1}{16\pi G}}{\displaystyle _C}ϵ_{bca_1\mathrm{}a_{n2}}\left[2R_d^c(\xi _1^b\xi _2^d\xi _2^b\xi _1^d)+\xi _2^b\xi _1^c𝐋\right].`$ At the horizon, by using Eqs. (15), (16) and (A5)- (A11) we know $`{\displaystyle _C}ϵ_{bca_1\mathrm{}a_{n2}}\xi _2^b\xi _1^c𝐋`$ (25) $`=`$ $`{\displaystyle _C}\widehat{ϵ}_{a_1\mathrm{}a_{n2}}𝐋\left[{\displaystyle \frac{\|\chi \|}{\rho }}𝒯_2\rho _c+\left({\displaystyle \frac{\rho }{\|\chi \|}}+t\rho \right)_2\chi _c\right](𝒯_1\chi ^c+_1\rho ^c)`$ (26) $`=`$ $`{\displaystyle _C}\widehat{ϵ}_{a_1\mathrm{}a_{n2}}𝐋\left[{\displaystyle \frac{\|\chi \|}{\rho }}𝒯_2_1\rho ^2+\left({\displaystyle \frac{\rho }{\|\chi \|}}+t\rho \right)_2𝒯_1\chi ^2\right]`$ (27) $`=`$ $`0,`$ (28) $`{\displaystyle _C}ϵ_{bca_1\mathrm{}a_{n2}}2R_d^c(\xi _1^b\xi _2^d\xi _2^b\xi _1^d)`$ (29) $`=`$ $`{\displaystyle _C}\widehat{ϵ}_{a_1\mathrm{}a_{n2}}R_d^c\left({\displaystyle \frac{1}{\kappa }}{\displaystyle \frac{\chi ^2}{\rho ^2}}\right)\left[{\displaystyle \frac{\|\chi \|}{\rho }}\rho _c\rho ^d\left({\displaystyle \frac{\rho }{\|\chi \|}}+t\rho \right)\chi _c\chi ^d\right](𝒯_1D𝒯_2𝒯_2D𝒯_1)`$ (30) $`=`$ $`0,`$ (31) and $`\xi ^bϵ_{bca_1a_2\mathrm{}a_{n2}}F^{dc}\left[F_{ed}\xi ^e+(\xi ^eA_e)_{;d}\right]`$ (32) $`=`$ $`\xi ^bϵ_{bca_1a_2\mathrm{}a_{n2}}F^{dc}\delta _\xi A_d`$ (33) $`=`$ $`\widehat{ϵ}_{a_1a_2\mathrm{}a_{n2}}\left[{\displaystyle \frac{\|\chi \|}{\rho }}𝒯\rho _c+\left({\displaystyle \frac{\rho }{\|\chi \|}}+t\rho \right)\chi _c\right]F^{dc}\delta _\xi A_d`$ (34) $`=`$ $`0.`$ (35) Therefore, Eq. (20) can be rewritten as $`\delta _{\xi _2}H[\xi _1]`$ $`=`$ $`{\displaystyle \frac{1}{16\pi G}}{\displaystyle _C}ϵ_{bca_1\mathrm{}a_{n2}}\left[\xi _2^b_d(^d\xi _1^c^c\xi _1^d)\xi _1^b_d(^d\xi _2^c^c\xi _2^d)\right].`$ (36) Since the “bulk” part of the generator $`H[\xi _1]`$ on the left side vanishes on shell, we can interpret the left side of Eq. (20) the variation of the boundary term $`J`$, i.e., $`\delta _{\xi _2}J[\xi _1]`$. On the other hand, the change in $`J[\xi _1]`$ under a surface deformation generated by $`J[\xi _2]`$ can be precisely described by Dirac bracket $`\{J[\xi _1],j[\xi _2]\}^{}`$ , that is, $$\delta _{\xi _2}J[\xi _1]=\{J[\xi _1],J[\xi _2]\}^{}=\frac{1}{16\pi G}_Cϵ_{bca_1\mathrm{}a_{n2}}\left[\xi _2^b_d(^d\xi _1^c^c\xi _1^d)\xi _1^b_d(^d\xi _2^c^c\xi _2^d)\right].$$ (37) Above discussions and Eq. (7) show the following relation on shell $$\{J[\xi _1],J[\xi _2]\}^{}=J[\{\xi _1,\xi _2\}]+K[\xi _1,\xi _2],$$ (38) Substituting Eqs. (A6), (A7), and (A11) into Eq. (37) we find that $`\{J[\xi _1],J[\xi _2]\}^{}`$ $`=`$ $`{\displaystyle \frac{1}{16\pi G}}{\displaystyle _C}\widehat{ϵ}_{a_1\mathrm{}a_{n2}}\left[{\displaystyle \frac{1}{\kappa }}(𝒯_1D^3𝒯_2𝒯_2D^3𝒯_1)2\kappa (𝒯_1D𝒯_2𝒯_2D𝒯_1)\right].`$ (39) It is also easy to show that $$\{\xi _1,\xi _2\}^a=(𝒯_1D𝒯_2𝒯_2D𝒯_1)\chi ^a+\frac{1}{\kappa }\frac{\chi ^2}{\rho ^2}D(𝒯_1D𝒯_2𝒯_2D𝒯_1)\rho ^a.$$ (40) On the other hand, the integrand of the right hand of Eq. (5) can be expressed as $$\xi ^b𝚯_{ba_1\mathrm{}a_{n2}}=\frac{1}{4\pi }\xi ^bϵ_{bca_1\mathrm{}a_{n2}}\left\{\frac{1}{2G}(_e^{[e}\xi ^{c]}+R_e^c\xi ^e)+F^{dc}\left[F_{ed}\xi ^e+(\xi ^eA_e)_{;d}\right]\right\}.$$ (41) The first two terms in the right hand of Eq. (41) can be treated as Carlip did in Ref. . And by using Eqs. (35) we know that the last two terms in Eq. (41) gives no contribution to $`K[\xi _1,\xi _2]`$. Making use of Eqs. (2), (18), (A6), (A7), and (A11) and replacing $`\xi ^a`$ in $`J`$ by $`\{\xi _1,\xi _2\}^a`$, we have $$J[\{\xi _1,\xi _2\}]=\frac{1}{16\pi G}_C\widehat{ϵ}_{a_1a_2\mathrm{}a_{n2}}\left[2\kappa (𝒯_1D𝒯_2𝒯_2D𝒯_1)\frac{1}{\kappa }D(𝒯_1D^2𝒯_2𝒯_2D^2𝒯_1)\right].$$ (42) The central term can then be obtained from Eqs.(38), (39), and (42), which is explicitly given by $$K[\xi _1,\xi _2]=\frac{1}{16\pi G}_C\widehat{ϵ}_{a_1a_2\mathrm{}a_{n2}}\frac{1}{\kappa }(D𝒯_1D^2𝒯_2D𝒯_2D^2𝒯_1).$$ (43) ## III Entropy of some rotating charged black holes In this section, lets us study statistical-mechanical entropies of the stationary axisymmetric black holes by using the constraint algebra constructed in the preceding section and conformal field theory methods. ### A Entropy of the Kerr-Newman black hole In Boyer-Lindquist coordinates, the metric of the Kerr-Newman black hole takes the form $$ds^2=\frac{\mathrm{\Delta }}{\varrho ^2}\left[dta\mathrm{sin}^2\theta d\phi \right]^2+\frac{\varrho ^2}{\mathrm{\Delta }}dr^2+\varrho ^2d\theta ^2+\frac{\mathrm{sin}^2\theta }{\varrho ^2}\left[adt(r^2+a^2)d\phi \right]^2,$$ (44) with $`\varrho ^2`$ $`=`$ $`r^2+a^2\mathrm{cos}^2\theta ,`$ (45) $`\mathrm{\Delta }`$ $`=`$ $`(rr_+)(rr_{}),`$ (46) where $`r_+=r_H=M+\sqrt{M^2Q^2a^2}`$, $`r_{}=M\sqrt{M^2Q^2a^2}`$, the parameter $`a`$ is related to the angular momentum, and $`M`$ and $`Q`$ represent the mass and electric charge of the black hole, respectively. The metric (44) is a solution of the Einstein-Maxwell field equations with an electromagnetic vector potential $`𝐀`$ $`=`$ $`{\displaystyle \frac{Qr}{\varrho ^2}}(dtasin^2\theta d\phi ),`$ (47) and associated field strength tensor $`𝐅=`$ $``$ $`{\displaystyle \frac{Q}{\varrho ^4}}(r^2a^2\mathrm{cos}^2\theta )e^0e^1`$ (48) $`+`$ $`{\displaystyle \frac{Q}{\varrho ^4}}(r^2a^2\mathrm{cos}^2\theta )e^2e^3.`$ (49) The Killing vector can be expressed as $$\chi _H^a=(1,0,0,\mathrm{\Omega }_H),$$ (50) where $`\mathrm{\Omega }_H=\left(\frac{g_{t\phi }}{g_{\phi \phi }}\right)_H=\frac{a}{r^2+a^2}`$ is the angular velocity of the black hole. A one-parameter group of diffeomorphism satisfying Eqs. (A10) and (40) can be taken as $$𝒯_n=\frac{1}{\kappa }exp\left[in(\kappa t+C_\alpha (\phi \mathrm{\Omega }_Ht))\right],$$ (51) where $`C_\alpha `$ is a arbitrary constant. Substituting Eq. (51) into central term (43) and using condition (A10) we obtain $$K[𝒯_m,𝒯_n]=\frac{iA_H}{8\pi G}m^3\delta _{m+n,0},$$ (52) where $`A_H=_C\widehat{ϵ}_{a_1a_2\mathrm{}a_{n2}}=4\pi (r_+^2+a^2)`$ is the area of the event horizon. Thus, Eq. (38) takes standard form of a Virasoro algebra $$i\{J[𝒯_m],J[𝒯_n]\}=(mn)J[𝒯_{m+n}]+\frac{c}{12}m^3\delta _{m+n,0},$$ (53) with central charge $`\frac{c}{12}=\frac{A_H}{8\pi G}.`$ The boundary term $`J[𝒯_0]`$ can easily be obtained by using Eqs (2), (13), and (51), which is given by $`J[𝒯_0]=\frac{A_H}{8\pi G}.`$ The number of states with a given eigenvalue $`\mathrm{}`$ of $`J[𝒯_0]`$ grows asymptotically for large $`\mathrm{}`$ as $$\rho (\mathrm{})exp\left\{2\pi \sqrt{\frac{c}{6}\left(\mathrm{}\frac{c}{24}\right)}\right\}=exp\left[\frac{A_H}{4G}\right],$$ (54) and the statistical entropy of the Kerr-Newman black hole is $$\mathrm{log}\rho (\mathrm{})\frac{A_H}{4G},$$ (55) which coincides with the standard Bekenstein-Hawking entropy. ### B Entropy of the Kerr-Newman-AdS black hole Carter constructed the Kerr-Newman-AdS black hole in four dimensions many years ago, which can be explicitly given by $$ds^2=\frac{\mathrm{\Delta }_r}{\varrho ^2}\left[dt\frac{a}{\mathrm{\Xi }}\mathrm{sin}^2\theta d\phi \right]^2+\frac{\varrho ^2}{\mathrm{\Delta }_r}dr^2+\frac{\varrho ^2}{\mathrm{\Delta }_\theta }d\theta ^2+\frac{\mathrm{sin}^2\theta \mathrm{\Delta }_\theta }{\varrho ^2}\left[adt\frac{(r^2+a^2)}{\mathrm{\Xi }}d\phi \right]^2,$$ (56) with $`\varrho ^2`$ $`=`$ $`r^2+a^2\mathrm{cos}^2\theta ,`$ (57) $`\mathrm{\Delta }_r`$ $`=`$ $`(r^2+a^2)(1+l^2r^2)2Mr+q^2+p^2,`$ (58) $`\mathrm{\Delta }_\theta `$ $`=`$ $`1l^2a^2\mathrm{cos}^2\theta ,`$ (59) $`\mathrm{\Xi }`$ $`=`$ $`1l^2a^2,`$ (60) where the parameter $`M`$ is related to the mass, $`a`$ to the angular momentum, $`q`$ is proportional to the electric charge, $`p`$ is proportional to the magnetic charge, and $`l^2=\mathrm{\Lambda }/3`$ ( where $`\mathrm{\Lambda }`$ is the (negative) cosmological constant). The event horizon is located at $`r=r_+`$, the largest root of the polynomial $`\mathrm{\Delta }_r`$. The metric (56) is a solution of the Einstein-Maxwell field equations with an electromagnetic vector potential is given by $`𝐀`$ $`=`$ $`{\displaystyle \frac{qr}{\varrho ^2\mathrm{\Xi }}}(dtasin^2\theta d\phi ){\displaystyle \frac{p\mathrm{cos}\theta }{\varrho ^2\mathrm{\Xi }}}[adt(r^2+a^2)d\phi ],`$ (61) and the associated field strength tensor is $`𝐅=`$ $``$ $`{\displaystyle \frac{1}{\varrho ^4}}[q(r^2a^2\mathrm{cos}^2\theta )+2pra\mathrm{cos}\theta ]e^0e^1`$ (62) $`+`$ $`{\displaystyle \frac{1}{\varrho ^4}}[q(r^2a^2\mathrm{cos}^2\theta )2pra\mathrm{cos}\theta ]e^2e^3.`$ (63) The Killing vector is $$\chi _H^a=_t+\mathrm{\Omega }_H_\phi ,$$ (64) where $`\mathrm{\Omega }_H=\left(\frac{g_{t\phi }}{g_{\phi \phi }}\right)_H=\frac{\mathrm{\Xi }a}{r_+^2+a^2}`$ is the angular velocity of the black hole. The analysis of the preceding subsection goes through with virtually no changes, yields a statistical entropy $$S=\frac{A_H}{4G}=\frac{\pi }{G}\frac{r_+^2+a^2}{\mathrm{\Xi }},$$ (65) which also coincides with its Bekenstein-Hawking entropy. ## IV summary and discussion By using the covariant phase techniques we extend Carlip’s investigation in Ref. to a case containing a cosmological term and a electromagnetic field. If the event horizon is treated as a boundary with Carlip’s constraint conditions, the central extension of the constraint algebra is worked out, and a standard Virasoro subalgebra with a corresponding central charge is constructed for the stationary axisymmetric charged black hole. The statistical entropies of the Kerr-Newman black hole and the Kerr-Newman-AdS black hole are then obtained by using Cardy formula and the results agree with their Bekenstein-Hawking entropies. Since the static charged black holes, such as the Reissner-Nordström black hole and Reissner-Nordström-AdS black hole are special case of the metrics (44) and (56), respectively, the above results are also valid for the static charged black holes. The results obtained in this paper support Carlip’s supposition: regardless of the details of a quantum theory of gravity, symmetries inherited from the classical theory may be sufficient to determine the asymptotic behavior of the density of states. ###### Acknowledgements. Jiliang Jing would like to thank Profs. Yongjiu Wang, Zhiming Tang, and Zongyang Sun for several helpful discussions. This work was supported in part by the National Nature Science Foundation of China under grant number 19975018. ## A Boundary Conditions In this section, we list the Carlip’s boundary conditions for convenience. As Carlip did in Ref. we define a “stretched horizon” $$\chi ^2=ϵ.$$ (A1) where $`\chi ^2=g_{ab}\chi ^a\chi ^b`$, $`\chi ^a`$ is a Killing vector. The result of the computation will be evaluated at the event horizon of the black hole by taking $`ϵ`$ to zero. Near the stretched horizon, one can introduce a vector orthogonal to the orbit of $`\chi ^a`$ by $$_a\chi ^2=2\kappa \rho _a,$$ (A2) where $`\kappa `$ is the surface gravity. Vector $`\rho ^a`$ satisfies conditions $`\chi ^a\rho _a={\displaystyle \frac{1}{\kappa }}\chi ^a\chi ^b_a\chi _b=0,`$ everywhere (A3) $`\rho ^a\chi ^a,`$ at the horizon (A4) To preserve the “asymptotic” structure at horizon, we impose Carlip’s boundary conditions $$\delta \chi ^2=0,\chi ^at^b\delta g_{ab}=0,\delta \rho _a=\frac{1}{2\kappa }_a(\delta \chi ^2)=0,at\chi ^2=0,$$ (A5) where $`t^a`$ is a any unit spacelike vector tangent to boundary $`𝐌`$ of the spacetime $`𝐌`$. And the surface deformation vector is suggested as the following form $$\xi ^a=\rho ^a+𝒯\chi ^a,$$ (A6) where functions $``$ and $`𝒯`$ satisfy $`={\displaystyle \frac{1}{\kappa }}{\displaystyle \frac{\chi ^2}{\rho ^2}}\chi ^a_a𝒯,`$ everywhere (A7) $`\rho ^a_a𝒯=0,`$ $`\text{at the horizon}.`$ (A8) Fixing the average value of $`\stackrel{~}{\kappa }`$ ($`\stackrel{~}{\kappa }=\frac{a^2}{\chi ^2}`$, $`a^a=\chi ^b_b\chi ^a`$ is the acceleration of an orbit of $`\chi ^a`$ ) over a cross section of the horizon $$\delta _C\widehat{ϵ}\left(\stackrel{~}{\kappa }\frac{\rho }{\|\chi \|}\kappa \right)=0,$$ (A9) where $`\kappa `$ is the surface gravity and $`\widehat{ϵ}`$ is the induced volume measure on $``$ ($``$ denote the (n-2)-dimensional intersection of the Cauchy surface $`C`$ with the Killing horizon $`\chi ^2=0`$). The technical role of the condition (A9) is to guarantee the existence of generators $`H[\xi ]`$. For a one-parameter group of diffeomorphism such that $`D𝒯_\alpha =\lambda _\alpha 𝒯_\alpha `$, ( $`D\chi ^a_a`$ ), condition (A9) in turn implies an orthogonality relation $$_C\widehat{ϵ}𝒯_\alpha 𝒯_\beta \delta _{\alpha +\beta }.$$ (A10) By using the other future-directed null normal vector $`N^a=k^a\alpha \chi ^at^a,`$ with $`k^a=\frac{1}{\chi ^2}\left(\chi ^a\frac{\|\chi \|}{\rho }\rho ^a\right)`$ and a normalization $`N_a\chi ^a=1`$, the volume element can be expressed as $$ϵ_{bca_1\mathrm{}a_{n2}}=\widehat{ϵ}_{a_1\mathrm{}a_{n2}}(\chi _bN_c\chi _cN_b)+\mathrm{}\mathrm{},$$ (A11) the omitted terms do not contribute to the integral.
warning/0004/math0004085.html	ar5iv	text	# Heat Invariant 𝐸₂ for Nonminimal Operator on Manifolds with Torsion ## 1 Introduction Determination of the internal structure of an object via the spectra of different radiations and waves around the object is one of the archetypal problems in physics. More restricted mathematical version of this problem may be formulated as follows. A manifold (bundle) equipped with such structures as metric, curvature, torsion, gauge fields etc. and an elliptic (pseudo)differential operator acting on this manifold are given. What information about the manifold can one obtain studying the spectral properties of the operator? M. Kac phrases the problem in the evocative title of his paper “Can one hear the shape of a drum?”. The answer to this radical question is negative. In 1964, J. Milnor found a pair of isospectral (with respect to the Laplace operator) but non-isometric tori in dimension sixteen, and recently, in 1999, D. Schüth constructed continuous isospectral families of metrics on the product of spheres $`S^4\times S^3\times S^3.`$<sup>1</sup><sup>1</sup>1After the Milnor’s result many examples of multiply connected isospectral manifolds have been constructed, but the Schüth’s construction is the first example of closed simply connected isospectral but non-isometric Riemannian manifolds. Nevertheless, many global geometric invariants, such as dimension, volume, and total scalar curvature, are known to be spectrally determined. Moreover, various manifolds such as round spheres of dimension less than or equal to six and 2-dimensional flat tori are uniquely determined by the spectra of the Laplacian acting on them. One of the most constructive approaches to study the spectral properties of operators on manifolds is investigation of the heat kernel expansion. This approach can be described briefly as follows. Starting with an elliptic operator $`A`$ of the order $`2r`$, acting on a bundle whose base is a compact close $`n`$-dimensional manifold $`M`$, and introducing an additional “time” variable $`t`$ one can construct the heat operator $`A\frac{}{t}`$. Then one can compute the short-time asymptotic expansion of the diagonal elements of the kernel of this heat operator: $$x\|e^{tA}\|x\underset{m0}{}\mathrm{E}_\mathrm{m}(\mathrm{x}\|\mathrm{A})\mathrm{t}^{{\scriptscriptstyle \frac{\mathrm{m}\mathrm{n}}{2\mathrm{r}}}},\mathrm{t}+0.$$ (1) The coefficients $`E_m(x\|A)`$ in this expansion are spectral invariants of the operator $`A`$, and encode information about the asymptotic properties of the spectrum. These coefficients are called the heat invariants or heat kernel coefficients. They are also widely known under the names Hadamard coefficients <sup>2</sup><sup>2</sup>2It was Hadamard who introduced these coefficients for scalar operator $`A`$ already in 1923 and established their essential properties., Hadamard-Minakshisundaram-DeWitt-Seeley, DeWitt-Seeley-Gilkey (HMDS or HAMIDEW, DWSG) coefficients, according to papers of these authors . The heat invariants $`E_m`$ are of fundamental importance in quantum field theory, quantum gravity, spectral geometry and topology of manifolds. Many quantities of interest (such as the effective action, Green function, anomalies in quantum field theory , the indices of elliptic operators and the invariants of manifolds in spectral geometry ) are expressed in terms of the heat invariants. Most papers devoted to computation of the heat invariants deal with so-called minimal operators whose leading term is a power of the Laplacian and symbol is a scalar (w.r.t space-time indices). A typical example of such operators is $$A=\text{ }\text{ }\text{ }\text{ }\text{ }+X.$$ Here $`\text{ }\text{ }\text{ }\text{ }\text{ }=g^{\mu \nu }D_\mu D_\nu `$, $`D_\mu `$ is a covariant derivative including generally different connections (affine and spinor connections, gauge fields), $`X`$ is a matrix in internal space, i.e., an operator acting in sections of the bundle. For minimal operators there are efficient enough methods for computing the heat invariants. In this paper we consider nonminimal operator of the form $$A^{\mu \nu }=g^{\mu \nu }\text{ }\text{ }\text{ }\text{ }\text{ }+aD^\mu D^\nu +X^{\mu \nu },$$ (2) where $`X^{\mu \nu }`$ is a tensor field (bundle indices are assumed implicitly), $`a`$ is a scalar parameter which should satisfy to the condition $`a<1`$ for the positive definiteness and, hence, for the ellipticity of operator (2). In recent years special cases of operator (2) have been encountered by physicists studying the quantization of gauge and gravitational fields in arbitrary gauges . For example, the quantization of Yang-Mills field in an arbitrary covariant background gauge leads to the operator $$A_{\mu \nu }^{ab}=\delta _{\mu \nu }\stackrel{ab}{\text{ }\text{ }\text{ }\text{ }\text{ }}\left(\frac{1}{\alpha }1\right)D_\mu ^{ac}D_\nu ^{cb}2f^{acb}G_{\mu \nu }^c,$$ where $`D_\mu `$ is a covariant derivative containing the external field potential $`A_\mu `$, $`G_{\mu \nu }`$ is a corresponding field strength, $`f^{abc}`$ are the structure constants of a corresponding Lie algebra and $`\alpha `$ is a scalar (gauge) parameter. Another example: the quantization of electro-magnetic field in an external gravitational field leads to the operator $$A_{\mu \nu }=g_{\mu \nu }\text{ }\text{ }\text{ }\text{ }\text{ }\left(\frac{1}{\alpha }1\right)D_\mu D_\nu +R_{\mu \nu }$$ (for an analogous operator in quantum gravity see ). The torsion is defined as antisymmetric part of affine (or linear) connection $$T^\lambda {}_{\mu \nu }{}^{}=\mathrm{\Gamma }^\lambda {}_{\nu \mu }{}^{}\mathrm{\Gamma }^\lambda {}_{\mu \nu }{}^{},$$ where $`\mathrm{\Gamma }^\lambda _{\mu \nu }`$ are connection coefficients. The Einstein’s General Relativity (see e.g. ) is based on a special connection called Levi-Civita connection, i.e., symmetric and compatible with metric affine connection. This connection can be expressed completely in terms of metric and is torsionless. The General Relativity well describes the interaction of the matter with the gravity as far as macroscopic bulk matter is considered. However, on the microscopic level, where the elementary particles posses such quantum property as spin, it seems necessary to take into account the influence of spin on the geometry of space-time. To describe the interaction of spinning particles with the gravitation, a gravitation theory should include the non-vanishing torsion. In 1922 Elie Cartan first pointed out that there is no a priori reason to assume an affine connection to be symmetric in the context of General Relativity. He proposed also a theory of gravitation with torsion which development is known now as the Einstein-Cartan theory. The torsion arises naturally in the different (based on Poincaré and affine groups) gauge theories of gravity developed in the recent years (see Refs. ). Moreover, all kinds of modern superstring theories (for the recent review see e. g. ), which allow to deduce the properties of space-time, also predict, along with the metric, the existence of torsion. In we computed the heat invariants for operator (2) up to $`E_4`$ but for manifolds without torsion. In this paper we consider more general (and computationally much more difficult) case of manifolds with torsion. ## 2 Algorithm and Implementation The algorithm we use was developed by V. Gusynin . This algorithm is based on the covariant generalization of the pseudodifferential calculus given by Widom . The main advantage of this algorithm is its universality. It can be applied to the wide class of pseudodifferential operators, in particular, to the nonminimal and higher-order operators intractable by other methods such as the DeWitt ansatz for heat kernel matrix elements. The algorithm has the following main features. For a positive elliptic operator $`A`$ the spectrum of which lies inside a contour $`C`$, the heat operator $`\mathrm{exp}(tA)`$ can be expressed in terms of the resolvent $`(A\lambda )^1`$ via the formula $$e^{tA}=_C\frac{id\lambda }{2\pi }e^{t\lambda }(A\lambda )^1.$$ (3) The pseudodifferential calculus method uses the following representation for the matrix elements of the resolvent $$G(x,x^{},\lambda )x\|\frac{1}{A\lambda }\|x^{}=\frac{d^nk}{(2\pi )^n\sqrt{g(x^{})}}e^{il(x,x^{},k)}\sigma (x,x^{},k;\lambda ),$$ (4) where $`\sigma (x,x^{},k;\lambda )`$ is an amplitude, $`l(x,x^{},k)`$ is a (real) phase function which is a biscalar with respect to general coordinate transformations, $`k`$ is a wave vector. The resolvent satisfies the equation $`(A\lambda )G=1`$ which leads to the equation for the amplitude: $$(A(x,D_\mu +iD_\mu l)\lambda )\sigma (x,x^{},k;\lambda )=I(x,x^{}),$$ (5) where $`I(x,x^{})`$ is a transport function having both bundle and Lorentz indices. In the pseudodifferential calculus, it is assumed that in the flat space the phase function has the form $`l=(xx^{})_\mu k^\mu .`$ The covariant analogue of the linearity of the function $`l`$ is based on the requirement that all higher-order symmetrized covariant derivatives of $`l(x,x^{},k)`$ vanish at the points $`x=x^{}`$, i. e., satisfy the infinite set of relations : $$[\{D_{\mu _1}\mathrm{}D_{\mu _m}\}l]=0,m>1,$$ (6) where $`\{\mathrm{}\}`$ means symmetrizing in all indices, and $`[\mathrm{}]`$ means transition to coincidence limit $`(x=x^{})`$. In an analogous way, the covariant transport function should satisfy the relations: $$[\{D_{\mu _1}\mathrm{}D_{\mu _m}\}I]=0,m1.$$ (7) Equations (6) and (7) together with the “initial conditions” $`[l]=0,[D_\mu l]=k_\mu `$ and $`[I]=1\mathrm{l}`$ (unit operator) allow one to compute the coincidence limits for nonsymmetrized covariant derivatives $`[D_{\mu _1}\mathrm{}D_{\mu _m}l]`$ and $`[D_{\mu _1}\mathrm{}D_{\mu _m}I]`$. These nonsymmetrized derivatives are obtained directly from (6) and (7) by reducing all terms to a unified index ordering with the help of the Ricci identity. The resulting expressions are universal polynomials in the torsion $`T^\lambda _{\mu \nu }`$, curvature tensor $`R^\lambda _{\mu \nu \eta }`$, gauge curvature $`W_{\mu \nu }`$ and their covariant derivatives. In fact, once computed and stored the coincidence limits $`[D_{\mu _1}\mathrm{}D_{\mu _m}l]`$ and $`[D_{\mu _1}\mathrm{}D_{\mu _m}I]`$ can be used in many calculations for different operators $`A`$. The functions $`l(x,x^{},k)`$ and $`I(x,x^{})`$, introduced with the help of formulas (6) and (7),<sup>3</sup><sup>3</sup>3The existence of these functions has been proved in . play an important role in the covariant pseudodifferential calculus called also intrinsic symbolic calculus . In fact, just these universal functions manifest the geometric properties of a base manifold and a bundle. Expanding the amplitude $`\sigma `$ in degrees of homogeneity of $`k`$: $$\sigma =\underset{m=1}{\overset{\mathrm{}}{}}\sigma _m(x,x^{},k;\lambda ),$$ we obtain the recursion equations for $`\sigma _m`$ from equation (5). For example, for operator (2) these recursion expressions take the form $`A^{\mu \lambda }\sigma _{0\lambda \nu }=I_\nu ^\mu ,`$ $`A^{\mu \lambda }\sigma _{1\lambda \nu }+i\left[g^{\mu \lambda }(l+2D^\eta lD_\eta )+a(D^\mu D^\lambda l+D^\mu lD^\lambda +D^\lambda lD^\mu )\right]`$ $`\times \sigma _{0\lambda \nu }=0,`$ $`\mathrm{}`$ $`A^{\mu \lambda }\sigma _{m\lambda \nu }+i\left[g^{\mu \lambda }(l+2D^\eta lD_\eta )+a(D^\mu D^\lambda l+D^\mu lD^\lambda +D^\lambda lD^\mu )\right]`$ $`\times \sigma _{(m1)\lambda \nu }+(g^{\mu \lambda }+aD^\mu D^\lambda +X^{\mu \lambda })\sigma _{(m2)\lambda \nu }=0,m2,`$ where the matrix $`A^{\mu \nu }=g^{\mu \nu }(D^\eta lD_\eta l\lambda )aD^\mu lD^\nu l`$ is the principal symbol for operator (2). Solving the recursion equations we obtain $`\sigma _m`$. The heat invariants are expressed in terms of integrals of the coincidence limits $`[\sigma _m]`$: $$E_m(x\|A)=\frac{d^nk}{(2\pi )^n\sqrt{g}}_C\frac{id\lambda }{2\pi }e^\lambda [\sigma _m](x,k,\lambda )J([\sigma _m]).$$ (8) The integrals in (8) can be expressed in terms of gamma and Gauss hypergeometric functions for a wide class of operators $`A`$. The typical integral of terms of the coincidence limit $`[\sigma _m]`$ takes the form $`J\left({\displaystyle \frac{k^{2p}k_{\mu _1}\mathrm{}k_{\mu _{2s}}}{(k^{2r}\lambda )^l[(1a)k^{2r}\lambda ]^m}}\right)=`$ $`g_{\{\mu _1\mathrm{}\mu _{2s}\}}{\displaystyle \frac{\mathrm{\Gamma }((p+s+n/2)/r)}{(4\pi )^{n/2}2^sr\mathrm{\Gamma }(n/2+s)\mathrm{\Gamma }(l+m)}}F(m,(p+s+n/2)/r;l+m;a),`$ where $`g_{\{\mu _1\mathrm{}\mu _{2s}\}}`$ is a symmetrized sum of products of metric tensors. Using the fact that $`m`$ and $`l`$ are whole numbers, one can express the hypergeometric function in (2) in terms of elementary functions with the help of the Gauss relation $$a(1z)F(a+1,b;c;z)=(ca)F(a1,b;c;z)+(2acaz+bz)F(a,b;c;z),$$ (9) and using then the formula $`F(1,b;m;z)=(m1)!{\displaystyle \frac{(z)^{1m}}{(1b)_{m1}}}\left[(1z)^{mb1}{\displaystyle \underset{k=0}{\overset{m2}{}}}{\displaystyle \frac{(bm+1)_k}{k!}}z^k\right],`$ $`m=1,2,\mathrm{},mb1,2,\mathrm{},`$ where $`(a)_k=a(a+1)\mathrm{}(a+k1)`$ is the Pochhammer symbol (shifted factorial). During simplification of tensor expressions we use various symmetry properties of the tensors $`R^\lambda {}_{\eta \mu \nu }{}^{},T^\lambda {}_{\mu \nu }{}^{},W_{\mu \nu }`$, and also the Ricci identity $`[D_\mu ,D_\nu ]\phi _{\lambda _1\mathrm{}\lambda _k}^{\eta _1\mathrm{}\eta _l}={\displaystyle \underset{i=1}{\overset{l}{}}}R^{\eta _i}{}_{\alpha \mu \nu }{}^{}\phi _{\lambda _1\mathrm{}\lambda _k}^{\eta _1\mathrm{}\eta _{i1}\alpha \eta _{i+1}\mathrm{}\eta _l}`$ $`{\displaystyle \underset{i=1}{\overset{k}{}}}R^\alpha {}_{\lambda _i\mu \nu }{}^{}\phi _{\lambda _1\mathrm{}\lambda _{i1}\alpha \lambda _{i+1}\mathrm{}\lambda _k}^{\eta _1\mathrm{}\eta _l}+T^\alpha {}_{\mu \nu }{}^{}D_{\alpha }^{}\phi _{\lambda _1\mathrm{}\lambda _k}^{\eta _1\mathrm{}\eta _l}+W_{\mu \nu }\phi _{\lambda _1\mathrm{}\lambda _k}^{\eta _1\mathrm{}\eta _l},`$ the Bianchi identities for both affine and gauge curvatures $`D_\alpha R^\beta {}_{\gamma \delta ϵ}{}^{}+D_\delta R^\beta {}_{\gamma ϵ\alpha }{}^{}+D_ϵR^\beta _{\gamma \alpha \delta }`$ $`+T^\lambda {}_{\alpha \delta }{}^{}R_{}^{\beta }{}_{\gamma ϵ\lambda }{}^{}+T^\lambda {}_{\delta ϵ}{}^{}R_{}^{\beta }{}_{\gamma \alpha \lambda }{}^{}+T^\lambda {}_{ϵ\alpha }{}^{}R_{}^{\beta }{}_{\gamma \delta \lambda }{}^{}=0,`$ $`D_\alpha W_{\beta \gamma }+D_\beta W_{\gamma \alpha }+D_\gamma W_{\alpha \beta }`$ $`+W_{\alpha \lambda }T^\lambda {}_{\beta \gamma }{}^{}+W_{\beta \lambda }T^\lambda {}_{\gamma \alpha }{}^{}+W_{\gamma \lambda }T^\lambda {}_{\alpha \beta }{}^{}=0,`$ and the cyclic identity $`R_{\beta \gamma \delta }^\alpha +R_{\gamma \delta \beta }^\alpha +R_{\delta \beta \gamma }^\alpha +D_\beta T^\alpha {}_{\gamma \delta }{}^{}+D_\gamma T^\alpha {}_{\delta \beta }{}^{}+D_\delta T^\alpha _{\beta \gamma }`$ $`+T^\alpha {}_{\beta \lambda }{}^{}T_{}^{\lambda }{}_{\gamma \delta }{}^{}+T^\alpha {}_{\gamma \lambda }{}^{}T_{}^{\lambda }{}_{\delta \beta }{}^{}+T^\alpha {}_{\delta \lambda }{}^{}T_{}^{\lambda }{}_{\beta \gamma }{}^{}=0.`$ The above algorithm has been implemented in the C language. The C code of total length about 11000 lines contains about 250 functions for different manipulations with tensors and scalars. These functions are gathered into two programs DWSGCOEF and COLIM. The COLIM program computes coincidence limits of the $`l(x,x^{},k)`$ and $`I(x,x^{})`$ functions and writes them to the disk. Once computed and stored<sup>4</sup><sup>4</sup>4For the operators of different tensor ranks $`A,A^{\mu \nu },\mathrm{}`$ the coincidence limits for the functions $`I,I^{\mu \nu },\mathrm{}`$ should be computed separately. the coincidence limits, being universal functions, can be used in many calculations for different operators $`A`$. The DWSGCOEF program computes $`E_m`$ coefficients by the following steps: 1. Reading input information (operator, order $`m`$, etc.) 2. Computing a set of asymptotic operators for constructing recursion equations. 3. Computing $`\sigma _m`$ with the help of the recursion equations. 4. Taking the coincidence limit $`[\sigma _m]`$. 5. Integrating $`[\sigma _m]`$ to obtain the coefficient $`E_m`$. 6. Substituting tensor expressions for $`[D_{\mu _1}\mathrm{}D_{\mu _k}l]`$ and $`[D_{\mu _1}\mathrm{}D_{\mu _k}I]`$ into $`E_m`$. 7. Reducing hypergeometric to elementary functions in the scalar coefficients ( $`C_i`$ in the formulas of Section 3) including in the heat invariants in the case of nonminimal or higher-order operator. Eliminating possible linear dependencies among these scalar coefficients<sup>5</sup><sup>5</sup>5Usually there are many dependencies among the scalar coefficients which are not seen in terms of hypergeometric functions, i.e., quite different hypergeometric expressions may be reduced sometimes to the same elementary function. to make the resulting formulas as compact as possible. 8. Output $`E_m`$ (and its Lorentz trace in the nonminimal case). To cut down the swelling of the intermediate expressions, we use term-by-term strategy, i.e., the most cumbersome Steps 4-6 are applied consecutively to single terms of $`\sigma _m`$ generated during the execution of Step 3. ## 3 Heat Invariant $`E_2`$ We present here the full expression for the coefficient $`E_2`$ and also its trace with respect to Lorentz indices for nonminimal operator (2) on a curved manifold with the torsion and gauge field manifested itself in the gauge curvature $`W_{\mu \nu }`$. We consider the case of arbitrary dimension $`n`$ and also the most important<sup>6</sup><sup>6</sup>6The Atiyah–Singer index of an elliptic operator on a manifold of dimension $`n`$ can be expressed in terms of an integral of $`E_n`$ over the manifold. for $`E_2`$ case $`n=2.`$ In the below formulas the indices $`\alpha ,\beta ,\gamma `$ and $`\mu ,\nu `$ are dummy and free, correspondingly. We use the following definition for the torsion trace: $`T_\mu =T^\alpha _{\alpha \mu }`$. ### 3.1 Full Expression $`E_2`$ $`=`$ $`(4\pi )^{\frac{n}{2}}\{C_1X^{\mu \nu }C_2(X^{\nu \mu }+g^{\mu \nu }X_\alpha {}_{}{}^{\alpha })+C_3(W^{\mu \nu }+{\displaystyle \frac{8}{3}}D_\alpha T^{\alpha \mu \nu }`$ $`+{\displaystyle \frac{19}{6}}T_\alpha T^{\alpha \mu \nu })+C_4R^{\mu \nu }C_5R^{\nu \mu }+C_6(D_\alpha T^{\mu \alpha \nu }+D_\alpha T^{\nu \alpha \mu })`$ $`+C_7T_{\alpha \beta }{}_{}{}^{\mu }T_{}^{\alpha \beta \nu }+C_8T_{\alpha \beta }{}_{}{}^{\mu }T_{}^{\beta \alpha \nu }+C_9\left(T_{\alpha \beta }{}_{}{}^{\mu }T_{}^{\nu \alpha \beta }+T_{\alpha \beta }{}_{}{}^{\nu }T_{}^{\mu \alpha \beta }\right)`$ $`C_{10}T^\mu {}_{\alpha \beta }{}^{}T_{}^{\nu \alpha \beta }+C_{11}D^\mu T^\nu C_{12}D^\nu T^\mu +C_{13}T_\alpha \left(T^{\mu \alpha \nu }+T^{\nu \alpha \mu }\right)`$ $`C_{14}T^\mu T^\nu +C_{15}g^{\mu \nu }\left(R+D_\alpha T^\alpha \right)+C_{16}g^{\mu \nu }T_\alpha T^\alpha `$ $`+(C_{16}C_{15})g^{\mu \nu }T_{\alpha \beta \gamma }T^{\beta \alpha \gamma }C_{17}g^{\mu \nu }T_{\alpha \beta \gamma }T^{\alpha \beta \gamma }\}.`$ Coefficients $`C_i`$ in arbitrary dimension $`n`$: $`C_1`$ $`=`$ $`{\displaystyle \frac{1}{a(n2)n(n+2)}}\{(1a)^{\frac{n}{2}}(3an6a+4n+4)+an^32an^2`$ $`3an+6a4n4\},`$ $`C_2`$ $`=`$ $`{\displaystyle \frac{1}{a(n2)n(n+2)}}\left\{(1a)^{\frac{n}{2}}(an+2a4)+an2a+4\right\},`$ $`C_3`$ $`=`$ $`{\displaystyle \frac{1}{a(n2)n}}\left\{(1a)^{1\frac{n}{2}}(an8)+3an8a+8\right\},`$ $`C_4`$ $`=`$ $`{\displaystyle \frac{1}{6a(n2)n(n+2)}}\{(1a)^{\frac{n}{2}}(17a^2n^2+34a^2n17an^2168an`$ $`268a+140n+256)53an^2+40an+268a140n256\},`$ $`C_5`$ $`=`$ $`{\displaystyle \frac{1}{6a(n2)n(n+2)}}\{(1a)^{\frac{n}{2}}(15a^2n^2+30a^2n15an^2152an`$ $`244a+116n+256)43an^2+24an+244a116n256\},`$ $`C_6`$ $`=`$ $`{\displaystyle \frac{1}{6a^2(n2)n(n+2)}}\{(1a)^{\frac{n}{2}}(a^3n^2+2a^3na^2n^220a^2n36a^2`$ $`+12an+96a48)+a^2n^216a^2n+36a^2+12an96a+48\},`$ $`C_7`$ $`=`$ $`{\displaystyle \frac{1}{6a^2(n2)n(n+2)(n+4)}}\{(1a)^{\frac{n}{2}}(a^3n^36a^3n^2+a^2n^3`$ $`8a^3n+30a^2n^2+152a^2n24an^2+192a^2240an576a+144n`$ $`+288)7a^2n^3+18a^2n^2+64a^2n48an^2192a^2+96an+576a`$ $`144n288\},`$ $`C_8`$ $`=`$ $`{\displaystyle \frac{1}{6a^2(n2)n(n+2)}}\{(1a)^{1\frac{n}{2}}(a^2n^2+2a^2n24an48a+144)`$ $`7a^2n^2+34a^2n48a^248an+192a144\},`$ $`C_9`$ $`=`$ $`{\displaystyle \frac{1}{a^2(n2)n(n+2)(n+4)}}\{(1a)^{\frac{n}{2}}(a^2n^2+6a^2n+8a^212an`$ $`48a+48)a^2n^2+6a^2n8a^212an+48a48\},`$ $`C_{10}`$ $`=`$ $`{\displaystyle \frac{1}{12a^2(n2)n(n+2)(n+4)}}\{(1a)^{\frac{n}{2}}(a^3n^36a^3n^2+a^2n^3`$ $`8a^3n+30a^2n^2+152a^2n+192a^2192an768a+576)a^2n^3`$ $`6a^2n^2+88a^2n192a^296an+768a576\},`$ $`C_{11}`$ $`=`$ $`{\displaystyle \frac{1}{3a^2(n2)n(n+2)}}\{(1a)^{1\frac{n}{2}}(9a^2n^218a^2n+76an+152a`$ $`24)+2(13a^2n^2+6a^2n+76a^232an88a+12)\},`$ $`C_{12}`$ $`=`$ $`{\displaystyle \frac{1}{3a^2(n2)n(n+2)}}\{2(1a)^{1\frac{n}{2}}(5a^2n^210a^2n+38an+76a`$ $`+12)31a^2n^2+26a^2n+152a^288an128a24\},`$ $`C_{13}`$ $`=`$ $`{\displaystyle \frac{1}{6a(n2)n(n+2)(n+4)}}\{(1a)^{\frac{n}{2}}(a^3n^3+6a^3n^2a^2n^3+8a^3n`$ $`18a^2n^280a^2n+12an^296a^2+72an+96a48n+96)+a^2n^3`$ $`18a^2n^2+8a^2n+12an^2+96a^2120an96a+48n96\},`$ $`C_{14}`$ $`=`$ $`{\displaystyle \frac{1}{2a^2(n2)n(n+2)}}\{(1a)^{1\frac{n}{2}}(a^2n^22a^2n+8an+16a16)`$ $`a^2n^22a^2n+16a^232a+16\},`$ $`C_{15}`$ $`=`$ $`{\displaystyle \frac{1}{6a(n2)n(n+2)}}\{(1a)^{\frac{n}{2}}(a^2n^22a^2n+an^2+8an+12a24)`$ $`+an^3an^212a+24\},`$ $`C_{16}`$ $`=`$ $`{\displaystyle \frac{1}{12a^2(n2)n(n+2)(n+4)}}\{(1a)^{\frac{n}{2}}(a^3n^36a^3n^2+a^2n^3`$ $`8a^3n+6a^2n^2+8a^2n+48an+192a288)+a^2n^4+3a^2n^3`$ $`+2a^2n^248a^2n+96an192a+288\},`$ $`C_{17}`$ $`=`$ $`{\displaystyle \frac{1}{24a^2(n2)n(n+2)(n+4)}}\{(1a)^{\frac{n}{2}}(a^3n^36a^3n^2+a^2n^3`$ $`8a^3n+30a^2n^2+152a^2n+192a^2192an768a+576)+a^2n^4`$ $`+3a^2n^310a^2n^2+72a^2n192a^296an+768a576\}.`$ Coefficients $`C_i`$ in the dimension $`n=2`$: $`C_1`$ $`=`$ $`{\displaystyle \frac{3\mathrm{ln}(1a)}{4a}}{\displaystyle \frac{2a}{8(1a)}},`$ $`C_2`$ $`=`$ $`{\displaystyle \frac{\mathrm{ln}(1a)}{a}}+{\displaystyle \frac{2a}{2(1a)}},`$ $`C_3`$ $`=`$ $`{\displaystyle \frac{(a4)\mathrm{ln}(1a)}{2a}}+2,`$ $`C_4`$ $`=`$ $`{\displaystyle \frac{(17a67)\mathrm{ln}(1a)}{12a}}+{\displaystyle \frac{137a134}{24(1a)}},`$ $`C_5`$ $`=`$ $`{\displaystyle \frac{(15a61)\mathrm{ln}(1a)}{12a}}+{\displaystyle \frac{119a122}{24(1a)}},`$ $`C_6`$ $`=`$ $`{\displaystyle \frac{(a^29a+6)\mathrm{ln}(1a)}{12a^2}}+{\displaystyle \frac{(3a2)(a2)}{8a(1a)}},`$ $`C_7`$ $`=`$ $`{\displaystyle \frac{(a^212a+12)\mathrm{ln}(1a)}{12a^2}}{\displaystyle \frac{2a^29a+6}{6a(1a)}},`$ $`C_8`$ $`=`$ $`{\displaystyle \frac{(a^212a+18)\mathrm{ln}(1a)}{12a^2}}+{\displaystyle \frac{a6}{4a}},`$ $`C_9`$ $`=`$ $`{\displaystyle \frac{(a2)\mathrm{ln}(1a)}{4a^2}}{\displaystyle \frac{a^212a+12}{24a(1a)}},`$ $`C_{10}`$ $`=`$ $`{\displaystyle \frac{(a^212a+12)\mathrm{ln}(1a)}{24a^2}}{\displaystyle \frac{2a^29a+6}{12a(1a)}},`$ $`C_{11}`$ $`=`$ $`{\displaystyle \frac{(9a^238a+3)\mathrm{ln}(1a)}{6a^2}}{\displaystyle \frac{73a6}{12a}},`$ $`C_{12}`$ $`=`$ $`{\displaystyle \frac{(10a^238a3)\mathrm{ln}(1a)}{6a^2}}{\displaystyle \frac{79a+6}{12a}},`$ $`C_{13}`$ $`=`$ $`{\displaystyle \frac{(a6)\mathrm{ln}(1a)}{12a}}+{\displaystyle \frac{2a3}{6(1a)}},`$ $`C_{14}`$ $`=`$ $`{\displaystyle \frac{(a^24a+2)\mathrm{ln}(1a)}{4a^2}}{\displaystyle \frac{3a2}{4a}},`$ $`C_{15}`$ $`=`$ $`{\displaystyle \frac{(a3)\mathrm{ln}(1a)}{12a}}{\displaystyle \frac{7a10}{24(1a)}},`$ $`C_{16}`$ $`=`$ $`{\displaystyle \frac{(a^26)\mathrm{ln}(1a)}{24a^2}}{\displaystyle \frac{3a^2+a6}{24a(1a)}},`$ $`C_{17}`$ $`=`$ $`{\displaystyle \frac{(a^212a+12)\mathrm{ln}(1a)}{48a^2}}{\displaystyle \frac{3a^210a+6}{24a(1a)}}.`$ ### 3.2 Lorentzian Trace $`\mathrm{tr}_\mathrm{L}E_2`$ $`=`$ $`(4\pi )^{\frac{n}{2}}\{C_1X_\alpha {}_{}{}^{\alpha }+C_2RC_3T_{\alpha \beta \gamma }T^{\alpha \beta \gamma }+C_4T_{\alpha \beta \gamma }T^{\beta \alpha \gamma }`$ $`+C_5D_\alpha T^\alpha +(C_4+C_5)T_\alpha T^\alpha \}`$ $`C_i`$ for the trace in arbitrary dimension $`n`$: $`C_1`$ $`=`$ $`{\displaystyle \frac{(1a)^{\frac{n}{2}}+n1}{n}},C_2={\displaystyle \frac{(1a)^{\frac{n}{2}}(an+n+6)+n^2n6}{6n}},`$ $`C_3`$ $`=`$ $`{\displaystyle \frac{1}{24an(n+2)}}\{(1a)^{\frac{n}{2}}(a^2n^2+2a^2nan^226an48a+96)`$ $`+an^3+an^2+22an48a+96\},`$ $`C_4`$ $`=`$ $`{\displaystyle \frac{1}{12an(n+2)}}\{(1a)^{\frac{n}{2}}(a^2n^2+2a^2nan^214an24a+48)`$ $`an^3an^210an+24a48\},`$ $`C_5`$ $`=`$ $`{\displaystyle \frac{1}{6a(n2)n}}\{(1a)^{\frac{n}{2}}(a^2n^2+4a^2nan^210an36a+48)`$ $`+an^33an^2+14an36a+48\}.`$ $`C_i`$ for the trace in the dimension $`n=2`$: $`C_1`$ $`=`$ $`{\displaystyle \frac{2a}{2(1a)}},C_2={\displaystyle \frac{2+a}{6(1a)}},C_3={\displaystyle \frac{1}{12}},C_4={\displaystyle \frac{1}{6}},`$ $`C_5`$ $`=`$ $`{\displaystyle \frac{(a4)\mathrm{ln}(1a)}{2a}}{\displaystyle \frac{11a14}{6(1a)}}.`$ ## 4 Conclusion The program computes $`E_2`$ with torsion for operator (2) rather easily (about 10 sec on a Pentium-75 PC). Unfortunately, computational complexity of the problem under consideration is very high. For example, the timings for torsionless computations of $`E_2`$ and $`E_4`$ for the same operator (2) are $`<1`$ sec and 4 h 5 min, correspondingly. It is clear that the inclusion of the torsion increases the computational efforts considerably and the computation of $`E_4`$ with torsion may take too much time. Another problem is the volume of the resulting expressions. There are two ways to handle this problem. First of all, some work is needed for developing of algorithms for further reduction of large tensor expressions. However, due to the natural complexity of the heat invariants, one can not hope to make the higher-order invariants tractable by hand. Thus, the methods for automatic usage of these invariants should be elaborated. ## Acknowledgements I would like to thank V. Gusynin for initiating this work and helpful communications. This work was supported in part by INTAS project No. 96-0842 and RFBR project No. 98-01-00101.
warning/0004/nucl-th0004045.html	ar5iv	text	# Practical approximation scheme for the pion dynamics in the three-nucleon system ## I Introduction Considerable progress has been made in understanding the coupled system of two nucleons and (at least) one pion. This enables the study of pion absorption and production processes on very light nuclei, and is a next step of including explicit pion degrees of freedom, beyond the standard nuclear picture where mesonic degrees of freedom are “frozen out”, into nucleon-nucleon potentials. We are in particular interested in a set of theories which can be classified by the acronym TRABAM (Thomas-Rinat, Afnan- Blankleider, Avishai-Mizutani), Ref. . In recent years, a fair amount of effort has been made to extend this theory to the pion-three-nucleon domain, where a richer range of phenomena is possible. In a sequence of papers , this system has been explored in the attempt to arrive at a consistent and connected theory of the coupled $`\pi `$NNN-NNN system. This effort has required the blending of the standard three- and four-body theories of Alt, Grassberger and Sandhas (AGS) , with the possibility that a pion can appear and disappear anywhere in the system, and in any of its subsystems. In a recent paper, one of the authors elucidated a connected, coupled scheme for the combined $`\pi `$NNN-NNN dynamics. Furthermore, he derived, by the use of the quasi-particle formalism at three-cluster and two-cluster levels, an equation for the coupled $`\pi `$NNN-NNN system, that has the appearance of a coupled set of Lippmann-Schwinger type equations: $$X_{ss^{}}^{(2)}=Z_{ss^{}}^{(2)}+\underset{s^{\prime \prime }}{}Z_{ss^{\prime \prime }}^{(2)}\tau _{s^{\prime \prime }}^{(2)}X_{s^{\prime \prime }s^{}}^{(2)}.$$ (1) An equation of this type was first given by Lovelace in the standard three-particle problem with separable interactions, using the original theory by Faddeev . It was also derived by AGS who applied the quasiparticle approximation to their three-body equations. Equation (1), like the Faddeev-Lovelace-AGS equation, is a coupled set of integral equations in one inter-cluster momentum variable, but here the labels $`s(s^{},s^{\prime \prime })`$ run over four values, rather than three: $`s=0`$ representing the configuration consisting of a three-nucleon cluster with the pion separate, and $`s=1,\mathrm{\hspace{0.17em}2},\mathrm{\hspace{0.17em}3}`$ representing the three possible arrangements of the three nucleons in a pair and a separated nucleon (the usual AGS-Lovelace scheme) with the pion associated with either the pair or the single nucleon, or absent. (See table II in .) The driving term in this equation is given in terms of quantities associated with the two levels of separable approximations needed when the quasiparticle method is applied to the full $`\pi `$NNN-NNN system of equations, as follows $`Z_{ss^{}}^{(2)}`$ $`=`$ $`\left(s^{(2)}\right)_{}\left\|g_0\right\|\left(s_{}^{}{}_{}{}^{(2)}\right)_{}\overline{\delta }_{ss^{}}`$ (2) $`+`$ $`{\displaystyle \underset{a^{}(s)}{}}{\displaystyle \underset{b^{}(s^{})}{}}{\displaystyle \underset{a(a^{},b^{})}{}}\left(s^{(2)}\right)_{a^{}a}\left\|\tau _a^{(3)}\right\|\left(s^{(2)}\right)_{b^{}a}(\overline{\delta }_{ss^{}}+\delta _{ss^{}}\overline{\delta }_{a^{}b^{}}).`$ (3) This equation looks rather complicated with the multiple sums and inclusion rules, but, in fact, it turns out to be rather simple, once all these rules are applied and the delta-functions invoked. We will define in detail symbols in this equation below, but let us first look at the structure of this driving term. The first term contains only components from the three-nucleon (no-pion) sector, and therefore cannot contribute at all if $`s=0`$ or $`s^{}=0`$. Therefore, the only contributions to this term can come from the case $`s0`$ and $`s^{}0`$. In this case the first term of Eq. (2) does not contribute on the diagonal (i.e. $`s=s^{}`$) because of the anti-delta function. Off the diagonal, since this contains only nucleon degrees of freedom, this term can properly be identified with the driving term of the traditional AGS-Lovelace approach to the standard three-nucleon problem. For $`s=s^{}=0`$ there is no contribution to $`Z_{ss^{}}^{(2)}`$ also from the second term, because the inclusion prescriptions in the sums cannot be satisfied if $`s=s^{}=0`$. More explicitly, since $`s=s^{}`$, only the last term in the delta-function structure survives, but this requires $`a^{}b^{}`$. However, for $`s=0`$, there is only one two-cluster state, namely $`\pi (N_1N_2N_3)`$, so the condition $`a^{}b^{}`$ cannot be satisfied. Therefore, the driving terms do not contribute at all in the case $`s=s^{}=0`$, either in the first, or in the second term. The second term does have non-zero values in the diagonal case $`s=s^{}0`$, and this provides additional diagonal terms to Eq. (1) due to the coupling to the pion degrees of freedom. Even if small, such contributions compare to zero and therefore can hardly be neglected. The second term contributes also to the off-diagonal elements, thus providing corrections to the dominant terms of the standard three-nucleon problem. In Eq. (2), $`\|\left(s^{(2)}\right)_{a^{}a},\|\left(s^{(2)}\right)_{}`$ are components of the form-factor vector coming from the separable approximation at the two-cluster level, the subscript “-” representing the no-pion sector, the subscripts $`a^{}a`$ being the Yakubovskĭ chain labels for the chains of partition of the (conserved) four-body system, $`\pi `$NNN. The quantity $`g_0`$ is the free propagator for three nucleons, in the absence of pions; the quantity $`\tau _a^{(3)}`$ denotes the intermediate propagation of the possible three-cluster structures of the $`\pi `$NNN system. In other words, the separable representation of the two-particle $`t`$-matrix, according to $$t_a(z)\|\left(a^{(3)}(z)\right)\tau _a^{(3)}(z)\left(a^{(3)}(z)\right)\|,$$ (5) describes effectively all the elastic two-body processes in the 4-body space, with $`t_a`$ representing either the $`NN`$ or the $`\pi N`$ two-body $`t`$-matrix. The “anti-delta” function is: $`\overline{\delta }_{ss^{}}=1\delta _{ss^{}}`$. The inclusions under the summation symbols are intended in the usual sense of Yakubovskĭ chain inclusions. For example, for $`s=1,a^{}`$ is one of the two partitions $`N_1(N_2N_3\pi )`$, $`(\pi N_1)(N_2N_3)`$, and $`a`$ are all the possible three-cluster partitions that can be obtained by breaking one cluster in each of the above two-cluster partitions. As an example of how the complicated-looking structure of Eq. (2) simplifies, we show two particular contributions to the driving term, where we specify the form-factor vectors by the chains of partition, and leave off the superscripts (2) and (3) that occur in Eq. (2); a typical off-diagonal element: $`Z_{12}=`$ $`N_1(N_2N_3)\left\|g_0\right\|N_2(N_3N_1)`$ $`(1.2a)`$ $`+`$ $`N_1(N_2N_3\pi );N_1N_2(N_3\pi )\left\|\tau _{(N_3\pi )}\right\|N_2(N_3N_1\pi );N_1N_2(N_3\pi )`$ $`+`$ $`N_1(N_2N_3\pi );N_1N_3(N_2\pi )\left\|\tau _{(N_2\pi )}\right\|(\pi N_2)(N_3N_1);N_1N_3(N_2\pi )`$ $`+`$ $`(\pi N_1)(N_2N_3);N_2N_3(N_1\pi )\left\|\tau _{(N_1\pi )}\right\|N_2(N_3N_1\pi );N_2N_3(N_1\pi ),`$ and a diagonal element: $`Z_{11}=`$ $`N_1\left(N_2N_3\pi \right);N_1\left(N_2N_3\right)\pi \left\|\tau _{\left(N_2N_3\right)}\right\|\left(N_1\pi \right)\left(N_2N_3\right);N_1\left(N_2N_3\right)\pi `$ $`(1.2b)`$ $`+`$ $`\left(N_1\pi \right)\left(N_2N_3\right);N_1\left(N_2N_3\right)\pi \left\|\tau _{\left(N_2N_3\right)}\right\|N_1\left(N_2N_3\pi \right);N_1\left(N_2N_3\right)\pi .`$ For $`s=1`$, as can be immediately deduced by comparing Eq. (2) with Eq. (1.2$`a`$), the state $`\|(s)_{}`$ denotes the Faddeev component $`\|N_1(N_2N_3)`$, while with $`\|(s)_{a^{}a}`$ we denote the relevant Yakubovskĭ components, such as $`\|N_1(N_2N_3\pi );N_1N_2(N_3\pi )`$ for instance. All possible 4-body components for $`s=1`$ are listed in Tab. I. To solve Eq. (1) one must first define and construct the states $`\|\left(s^{(2)}\right)_{a^{}a}`$, $`\|\left(s^{(2)}\right)_{}`$, and the two-cluster Green’s function $`\tau _s^{(2)}`$. It is precisely at this point that we propose herein a workable approximation scheme. However, before discussing the approximation, we recall for clarity the rigorous result obtained in Ref. . For complete details we refer to that work. Starting from the two-particle representation Eq. (5), one can derive the dynamical sub-amplitude containing the interactions internal to the coupled set of partitions $`(N_1)(N_2N_3)`$, $`(N_1\pi )(N_2N_3)`$, and $`(N_1)(N_2N_3\pi )`$. This is denoted by $`(𝐱)_{(s=1)}`$ where the $`s=1`$ index should be interpreted with the fact that nucleon “1” is set apart from the other two: obviously, the other cases with $`s=2,3`$ are obtained by cyclic permutations of the nucleons. The $`s=0`$ amplitude instead contains all the interactions internal to the single $`\pi (N_1N_2N_3)`$ partition, hence the pion is set apart from the three nucleons here. The dynamical equations for these four subamplitudes are given by (see Eqs. (3.21-24) of Ref. ) $`\left(x_s\right)_{a^{}a,b^{}b}=`$ $`a^{(3)}\|G_0\|b^{(3)}\overline{\delta }_{ab}\delta _{a^{}b^{}}+{\displaystyle \underset{c^{}(s)}{}}{\displaystyle \underset{c(c^{})}{}}a^{(3)}\|G_0\|c^{(3)}\overline{\delta }_{ac}\delta _{a^{}c^{}}\tau _c^{(3)}\left(x_s\right)_{c^{}c,b^{}b}`$ (6) $`+a^{(3)}\|G_0\left(f_s\right)_{a^{}a}g_0\left(x_{s}^{}{}_{}{}^{}\right)_{,b^{}b}`$ (7) $`\left(x_{s}^{}{}_{}{}^{}\right)_{,b^{}b}=`$ $`\left(f_s^{}\right)_{b^{}b}G_0\|b^{(3)}+{\displaystyle \underset{c^{}(s)}{}}{\displaystyle \underset{c(c^{})}{}}\left(f_s^{}\right)_{c^{}c}G_0\|c^{(3)}\tau _c^{(3)}\left(x_s\right)_{c^{}c,b^{}b}`$ (8) $`+𝒱_sg_0\left(x_{s}^{}{}_{}{}^{}\right)_{,b^{}b}`$ (9) $`\left(x_s\right)_{a^{}a,}=`$ $`a^{(3)}\|G_0\left(f_s\right)_{a^{}a}+{\displaystyle \underset{c^{}(s)}{}}{\displaystyle \underset{c(c^{})}{}}\overline{\delta }_{ac}\delta _{a^{}c^{}}a^{(3)}\|G_0\|c^{(3)}\tau _c^{(3)}\left(x_s\right)_{c^{}c,}`$ (10) $`+a^{(3)}\|G_0\left(f_s\right)_{a^{}a}g_0\left(x_s\right)_,`$ (11) $`\left(x_s\right)_,=`$ $`𝒱_s+𝒱_sg_0\left(x_s\right)_,+{\displaystyle \underset{c^{}(s)}{}}{\displaystyle \underset{c(c^{})}{}}\left(f_s^{}\right)_{c^{}c}G_0\|c^{(3)}\tau _c^{(3)}\left(x_s\right)_{c^{}c,}`$ (12) with $`aa^{}s`$ and $`bb^{}s`$. For each $`s`$, the subamplitudes $`\left(𝐱_s\right)`$ have components labelled by the Yakubovskĭ chain labels ($`a^{}a`$), or by the symbol “-” in the case of the no-pion sector, where the index s specifies a unique physical partition (last column of table II in ). Only when $`s=0`$, all couplings to the no-pion sector are vanishing. These coupled equations are obtained once the representation Eq. (5) has been assumed, and the amplitudes are expressed in the 4-body space in the corresponding quasiparticle representation. $`G_0`$ is the full 4-body free propagator, $`g_0`$ is the already mentioned free three-nucleon propagator in the no-pion sector, $`𝒱_s`$ is the pair potential between the two interacting nucleons in the partitions denoted by s and, in addition to the short-range part of the 2N potential, it includes explicitely the OPE diagram. Finally, $`f_s\left(f_s^{}\right)`$ is the elementary creation (annihilation) vertex for a pion into (from) the appropriate subsystem denoted by the chain-label subscript. An important aspect of these subamplitudes is that it is always possible to factor out a Dirac $`\delta `$ function in momentum space for the “spectator” nucleon (or pion, for $`s=0`$). The new aspect obtained in Ref. is that this factorization property has been maintained when there is a pion associated with either the nucleon pair or the spectator nucleon, or when there is no pion at all. The basic assumption for the $`x_s`$ subamplitudes consists in the finite-rank representation $`(x_s)_{a^{}a,b^{}b}`$ $`=`$ $`\|(s^{(2)})_{a^{}a}\tau _s^{(2)}(s^{(2)})_{b^{}b}\|`$ (13) $`(x_s^{})_{,b^{}b}`$ $`=`$ $`\|(s^{(2)})_{}\tau _s^{(2)}(s^{(2)})_{b^{}b}\|`$ (14) $`(x_s)_{a^{}a,}`$ $`=`$ $`\|(s^{(2)})_{a^{}a}\tau _s^{(2)}(s^{(2)})_{}\|`$ (15) $`(x_s)_,`$ $`=`$ $`\|(s^{(2)})_{}\tau _s^{(2)}(s^{(2)})_{}\|.`$ (16) This representation provides all the ingredients needed to construct the connected dynamical equation (1). We have limited here the discussion to the case of one separable term, but the algebraic generalization of Eq.(1) to more separable terms is straightforward. ## II Approximation scheme Although this is a very relevant issue, we will not concentrate here on the mathematical aspects and general constraints needed to obtain a mathematically converging separable expansion for the subamplitudes $`𝐱_s`$. A very clear explanation about the general methods required for considering such questions can be found in Ref. . We will instead concentrate on the development of an approximate scheme for the separable representation of such subamplitudes. The proposed approximation scheme is based on the physical assumption that in standard nuclear physics the picture of nucleons interacting via realistic nucleon-nucleon potentials, with the pion degrees of freedom “frozen out”, provides an acceptable first-order description of the low-energy/low-momenta dynamics. The effects of including explicit pionic degrees of freedom beyond that picture should then be considered only as dynamical “corrections”. If we consider Eq. (12), we observe that the last term accounts in fact for the one-pion dynamics in the 2N subsystem ($`s0`$), once the OPE diagram has been taken out (because it is already included in $`𝒱_s`$). In particular $`(x_s)_,`$ represents the complete, elastic 2N amplitude in presence of a spectator nucleon. Obviously, $`(x_s)_,`$ should be obtained from the coupled set of Eqs. (7-12), but we will identify instead $`(x_s)_,`$ with the conventional 2N $`t`$-matrix, derived by the solution of the standard 2N Lippmann-Schwinger equation (in the presence of a spectator nucleon) using as input the phenomenological NN potential. In other words, we set to zero the last term of Eq. (12) and use for $`𝒱_s`$ the conventional NN potential which includes in an effective way the contributions from the pion-2N dynamics. Note that there will be a price to pay for this; namely, all disconnected dispersive effects to the 3N dynamics, originated in the 2N subsystems by the dynamical equations (7-12), will be approximated to zero. This is a consequence of the fact that the dynamical description of the pion degrees of freedom implied by these equations has been replaced with an istantaneous, effective 2N potential. Implicitely, this approximation is assumed in all potential approaches to the 3N problem, but it has not really been tested. The exception is in Ref. , where these dispersive effects of the 2N subsystem have been sized in the extreme situation where the $`\pi `$N interaction is entirely represented by its coupling through a forward propagating $`\mathrm{\Delta }`$-isobar. Interestingly, the $`\mathrm{\Delta }`$-mediated, disconnected dispersive effects in the 3N system turned out to be not negligible. It is clear that this problem should be investigated further; nevertheless we will not do this here since our aim is to follow in this respect the standard potential approach to the 3N problem, where the 3N dispersive effects generated in the 2N subsystems are completely ignored. Once we have accepted that the 2N dynamics in the elastic channel is described by means of a phenomenological nucleon-nucleon potential, our description can be closely compared to the usual, potential-based, quantum-mechanical 3N approaches since the dynamical input of the two approaches appears to be identical. Then, if we consider Eq. (16), this reduces to the well-known, standard separable representation of the two-nucleon $`t`$-matrix, which we express in the polar form $$\left(x_s\right)_,\|\stackrel{~}{s}^{(2)}\stackrel{~}{\tau }_s^{(2)}\stackrel{~}{s}^{(2)}\|,$$ (17) and this already gives the approximations $`\|\left(s^{(2)}\right)_{}`$ $``$ $`\|\stackrel{~}{s}^{(2)}`$, and $`\tau ^{(2)}`$ $``$ $`\stackrel{~}{\tau }^{(2)}`$ which are needed for the determination of the driving term (first contribution) and for the kernel of Eq. (1). It should be noted that the form factors $`\|\stackrel{~}{s}^{(2)}`$ for NN interactions are, in this approximation, not distinct from the form factors $`\|a^{(3)}`$ of the quasiparticle approximation at the three-cluster level, Eq.(5). Both come, in fact, from the pole approximation of the elastic two-nucleon $`t`$-matrix. The only difference is that the $`\|\stackrel{~}{s}^{(2)}`$ factor refers only to the 2N $`t`$-matrices and is expressed in the (NN)+N two-cluster space, with one Jacobi coordinate removed already. The form factor $`\|a^{(3)}`$, on the other hand, refers to all the two-body $`t`$-matrices in the 4-body space, and this includes also the $`\pi `$N $`t`$-matrices in addition to the NN ones. Thus, a convenient feature emerges from the approach herein discussed: only one standard pole approximation (or expansion, in the more general case of higher ranks) has to be made for the two-nucleon interaction, to be used for the 2N $`t`$-matrices in both 4-body and 3N spaces. At this point one fact must be stressed: namely, if we do not consider further contributions, our description precisely collapses into the quantum-mechanical approach to the 3N system in terms of a 2N potential, since the resulting 3N equations have exactly the AGS form. This is because all couplings to $`\pi `$NNN state are made via the components $`\|\left(s^{(2)}\right)_{a^{}a}`$ which would be set to zero in this case. Thus, these components, together with the $`s=0`$ partition, are fundamental for an explicit treatment of the pion degrees of freedom in the 3N system, and we get an approximation of the $`\|\left(s^{(2)}\right)_{a^{}a}`$ factors, to the lowest order, by considering in particular Eq. (11). Assuming that the leading contributions to $`\left(x_s\right)_{a^{}a,}`$ are dominated by the pole structures (see Eq. (17)) of $`\left(x_s\right)_,`$ in the last term of Eq. (11), we obtain $$\left(x_s\right)_{a^{}a,}a^{(3)}\|G_0\left(f_s\right)_{a^{}a}g_0\|\stackrel{~}{s}^{(2)}\stackrel{~}{\tau }_s^{(2)}\stackrel{~}{s}^{(2)}\|,$$ and considering Eq. (14), $$\|\left(s^{(2)}\right)_{a^{}a}\|\left(\stackrel{~}{s}^{(2)}\right)_{a^{}a}a^{(3)}\left\|G_0\left(f_s\right)_{a^{}a}g_0\right\|\stackrel{~}{s}^{(2)}.$$ (18) This provides the form factors needed for the second term of Eq. (2), which in this approximation represents the contribution entirely responsible for the explicit treatment of the pion dynamics in Eq. (1), beyond that effectively contained in the static 2N potential. The $`a^{(3)}\|`$ on the left of this is needed because the quasiparticle approximation has already been made for the pair interaction in the four-body space, by Eq. (5). Eq. (2) also requires the “bra” vectors $`\left(s^{(2)}\right)_{}\|`$ , $`\left(\stackrel{~}{s}^{(2)}\right)_{a^{}a}\|`$. We simply construct the adjoints, by taking the corresponding “bra” states of the underlying quasiparticle approximation Eq. (17) of the standard 2N t-matrix. In other words $$\left(s^{(2)}\right)_{}\|\stackrel{~}{s}^{(2)}\|,$$ (19) and $$\left(s^{(2)}\right)_{a^{}a}\|\stackrel{~}{s}^{(2)}\left\|g_0\left(f_s^{}\right)_{a^{}a}G_0\right\|a^{(3)}.$$ (20) In the extended 3N equation, Eq. (1), we identify three types of new contributions which take into account, to the lowest order, the explicit pion dynamics. These contributions modify the driving term $`Z^{(2)}`$ of the standard AGS equation in a very selective way. The first contributions enlarge the number of components with respect to the three standard Faddeev components, with the addition of the fourth component ($`s=0`$) specifying the partition when the pion is set apart from the three nucleons. Such contributions provide the new couplings between $`s=0`$ and $`s^{}0`$, as well as with $`s0`$ and $`s^{}=0`$. These link the 3N Faddeev components with the partition consisting of a three-nucleon cluster and a separated pion ($`s=0`$). A second type of terms enter in the diagonal part, $`s=s^{}0`$, of $`Z_{ss^{}}^{(2)}`$ while traditionally these elements have been assumed to be vanishing in the standard 3N theory: hence, these terms should be considered important since they compare to zero in the standard AGS equation. Finally, the second term of Eq.(2) provides corrections also for $`0ss^{}0`$, that is for the off-diagonal elements of the driving term. Instances of such terms are shown in the last three contributions of Eq. (1.2$`a`$). They all represent modifications that have to be added to the standard 3N driving term, identified with the first contribution in Eq. (2). Thus, the new contributions coming from the explicit pion degrees of freedom, provide minimal, though important, modifications to the standard AGS formulation of the three-nucleon problem. ## III Discussion In the previous section we have shown that it is possible to treat approximately the pion dynamics in the 3N equations by introducing the new form factor $$\|\left(\stackrel{~}{s}^{(2)}\right)_{a^{}a}a^{(3)}\left\|G_0\left(f_s\right)_{a^{}a}g_0\right\|\stackrel{~}{s}^{(2)}.$$ (21) The operator $`\left(f_s\right)_{a^{}a}`$ has been defined in Ref. in terms of the renormalized elementary $`\pi NN`$ vertex by means of the inclusion prescription $$\left(f_s\right)_{a^{}a}=\underset{i=1}{\overset{3}{}}f_i\overline{\delta }_{ia}\delta _{i,aa^{}}\delta _{a^{}s}.$$ (22) Here, the label “$`i`$” represents the $`\pi N`$ pair interacting via the vertex $`f_i`$. In view of the key role played by this new ingredient in the extended 3N equation we provide its detailed diagrammatic interpretation. To fix the ideas, we choose $`s=1`$ and consider the various components depending on the Yakubovskĭ chain-of-partition subscript $`a^{}a`$, listed in Tab. I. By applying the inclusion prescription for the vertex operator, we obtain the set of diagrams drawn in Fig. 1. Note that there are only four diagrams in the figure while there are five components in the table. This is because the inclusion prescription automatically set to zero the contribution corresponding to the last Yakubovskĭ component, as shown in Tab I. Moreover, there is another diagram (omitted in Fig. 1) which has to be added to the first diagram shown in the figure, obtained by interchanging the two nucleons “2” and “3” within the pair. This sum is evidenced in the first row of Tab. I. All these contributions represent the proper quasiparticle generalization of the standard 2N form factor, in presence of a spectator nucleon, to the pion inelastic channel. Once we have illustrated the form factor diagrams, we can discuss the detailed structure of the new contributions to the driving term of the extended 3N equation. Indeed, we can interpret the main result obtained in Sect. III of Ref. , and reported here in Eqs.(1,2), as a simple prescription to include the pion dynamics in the AGS equation, $$Z_{ss^{}}=Z_{ss^{}}^{AGS}+Z_{ss^{}}^\pi $$ (23) and the main result of the previous section is a practical approximation scheme to get the part of the driving term which handles the pion dynamics when $`s,s^{}0`$ $$Z_{ss^{}}^\pi =\underset{a,a^{},b^{}}{}\stackrel{~}{s}^{(2)}\left\|g_0\left(f_s^{}\right)_{a^{}a}G_0\right\|a^{(3)}\tau _a^{(3)}a^{(3)}\left\|G_0\left(f_s^{}\right)_{b^{}a}g_0\right\|\stackrel{~}{s}^{(2)}(\overline{\delta }_{ss^{}}+\delta _{ss^{}}\overline{\delta }_{a^{}b^{}}).$$ (24) We begin with the contributions to Eq. (24) in the $`s=s^{}=1`$ case. Here, only the term with $`\delta _{ss^{}}`$ survives in Eq. (2), which implies $`a^{}b^{}`$, and also that only the ($`N_2N_3`$) pair can be formed in the intermediate state. Consequently, the only possible diagrams that can be constructed are obtained from only the first and last diagrams in Fig. 1 yielding the diagrams shown in Fig. 2, plus obviously those obtained by interchanging the pairing nucleons “2” and “3”, for a total of four time-ordered diagrams. Clearly, these are irreducible contributions that cannot be represented in conventional 2N potential theory. We then consider the contributions to the Z-pionic term when $`0ss^{}0`$, e.g., when $`s=1`$ and $`s^{}=2`$. In this case only the $`\overline{\delta }_{ss^{}}`$ contribution survives in the second term of Eq. (2); here the intermediate pair can only be formed with the pion and the nucleon $`N_3`$, which is uniquely defined since it does not act as spectator in both components $`s=1`$ and $`s^{}=2`$. In this case, one gets the diagram shown in Fig. 3. Note however that this is an approximated result: Once we have approximated the inelastic form factors in Eqs. (13-14) by the leading expressions (18, 20), then the last two terms in Eq. (1.2$`a`$) vanish, because the last row in Tab. I is empty in this case. However, in the more general approach of Ref. there are additional contributions (shown in Fig. 3 of Ref. ) which originate from the last two terms in Eq. (1.2$`a`$). Physically, these additional terms represent four-body multiple-rescattering contributions. So far, we have discussed the additional contributions one must include in the AGS equation to take minimally into account the pion dynamics beyond the OPE term. We observe that in both cases these diagrams represent, in fact, contributions that can be reinterpreted as irreducible 3N potentials. In particular, the diagram of Fig. 3 has the topological structure of the well-known Fujita-Miyazawa diagram. Similar forms (where the exchanged pion rescatters before being absorbed), are the basic ansatz for building the pion part of the irreducible 3N forces, such as the Tucson-Melbourne, Ruhr, Brazil, or Texas 3N interactions. However, this scheme provides at the same time also another set of irreducible 3N diagrams which must be considered as well; these are given in Fig. 2 and represent a totally different structure from that of Fig. 3. These irreducible diagrams have been proposed by Brueckner et al., as early as 1954, and have been investigated quantitatively for the triton binding energy by Pask, who finds they give a large contribution. Since that time, it has been argued (see Refs. ) that these terms cancel out against relativistic corrections to the twice-iterated one-pion exchange term. However, in a recent study , the effect of this cancellation has been questioned and in a quantitative analysis the cancellation turned out to be remarkably incomplete, with a breaking effect of 15-30%. The reason that breaks the cancellation is dynamical, and can be evidenced with an approach preserving the cluster sub-structures of the multinucleon system, while describing the pion-exchange dynamics. This is, we believe, one of the advandages of using the formulation discussed here. From the diagrams of Fig. 2 a new contribution to the 3N force has been extracted , and by means of this contribution the third nucleon affects in particular the triplet-odd waves of the 2N subsystem, with possible consequences for the N-d vector analyzing powers. Finally, we consider the contributions of the Z-pionic term when $`s=0`$ and $`s^{}0`$. Here, the pionic contributions are crucial in providing the couplings to the new, fourth component. In such a case, one has to modify Eq. (24) since the approximation discussed in the previous section concerns the structure of the $`\|s^{(2)}`$ form factors only when $`s0`$, while nothing is said otherwise. However, the $`s=0`$ case has been discussed in Ref. , see Eqs.(2.16-18) therein, and it turns out that for $`s=0`$ the only nonvanishing subamplitude is given by the equation $$\left(x_0\right)_{a^{}a,a^{}b}=a^{(3)}\|G_0\|b^{(3)}\overline{\delta }_{ab}+\underset{c(a^{})}{}a^{(3)}\|G_0\|c^{(3)}\overline{\delta }_{ac}\tau _c^{(3)}\left(x_0\right)_{a^{}c,a^{}b},$$ (25) which represents a Faddeev equation for three interacting nucleons in presence of a spectator pion, since $`a^{}=(NNN)\pi `$. In the vicinity of its poles, $`𝐱_0`$ becomes $$(x_0)_{a^{}a,a^{}b}\|(s^{(2)})_{a^{}a}\tau _s^{(2)}(s^{(2)})_{a^{}b}\|.$$ (26) The $`s=0`$ form factors are solutions of the homogeneous equation associated to Eq. (25), and represent the virtual decay of a 3N interacting cluster into a correlated 2N-pair, $`a`$, plus a nucleon, in presence of the spectator pion. For $`s=0`$, the form factor is obviously vanishing in the pure 3N sector, since there is always the presence of the spectator pion. Hence, taking $`s=0`$ and $`s^{}=1`$ for instance, the Z-pionic term becomes $$Z_{ss^{}}^\pi =\underset{a,a^{},b^{}}{}(s^{(2)})_{a^{}a}\|\tau _a^{(3)}a^{(3)}\left\|G_0\left(f_s^{}\right)_{b^{}a}g_0\right\|\stackrel{~}{s}^{(2)},$$ (27) once the complicated $`\delta `$-structure in Eq. (2) has been properly taken into account. We observe that for $`ss^{}`$ it is $`a^{}b^{}`$ always. Furthermore, since $`a^{}`$ identifies the $`s=0`$ partition, $`a`$ may represent only a $`NN`$ pair, which means that the term in Eq. (27) selects only the first and the fourth Yakubovskĭ components (in Tab. I), and this specifies completely the diagrams contributing to $`Z_{01}^\pi `$, as shown in Fig. 4. In the figure the first diagram should include an additional contribution obtained by interchanging the pairing nucleons, as usual. At this point, we have achieved the main goal of this work; indeed, starting from the more general approach of Ref. , we derived a practical, approximated scheme for the treatment of the pion dynamics in the 3N system. In particular, we obtained a set of dynamical equations where the main input is given by the 2N $`t`$-matrix generated by phenomenological NN potentials (hence, constrained by phase-shift analysis). The additional inputs required for this description are the $`\pi N`$ $`t`$-matrix, once its polar part has been subtracted, and the nonrelativistic $`\pi `$NN vertex, needed for the construction of new form factors expressed by Eq. (21). This set of equations can be considered a natural, approximated extension of 3N AGS equations for the explicit treatment of the pionic channel. We have identified the modifications implied by this treatment: they consist in additional terms, $`Z_{ss^{}}^\pi `$, representing corrections to the standard AGS driving term $`Z_{ss^{}}^{AGS}`$. We have classified these corrections by the structure of the underlying diagrams and found they all correspond to irreducible 3N-force diagrams. We have also pointed out the presence of a fourth, pionic, component representing the $`\pi `$ \+ (NNN) partition, and described how this is coupled to the other three Faddeev components. Finally, we have discussed the main limitation implied by the approach, which consists in ignoring all 3N disconnected (or reducible) dispersive effects: this is unavoidable if we use as input a potential-based description for the dynamics of the 2N subsystem; indeed, such an approximation is implicitly assumed in all conventional quantum-mechanical descriptions of the 3N system based on the potential approach. To overcome this limitation, one should consider the more general approach of Refs. , wherein the input interactions cannot be defined in a simple manner. With respect to this point, we observe that, since the 2N-subsystem dynamics is described in this work in terms of the standard 2N potential approach, the mass of the nucleon in the no-pion sector is not generated dynamically, but is static, and hence these approximated equations are not plagued by the nucleon-renormalization problem. On the other hand, above the pion threshold for the NN subsystem, the approach presented herein can not be considered unitary, since the input 2N $`t`$-matrices are not unitary either. It might be surprising to find out that the dynamical equations analyzed here as well as in Ref. lead to another, new class of irreducible 3N diagrams. Such a class of diagrams is present but still hidden in these dynamical equations since these represent formulations suited for the 3N problem above the pion threshold. If we consider the 3N system at lower energies, then it is more convenient to eliminate from these formulations the fourth Faddeev component, $`s=0`$, related to the mesonic channel. This can be accomplished by using in a very straightforward way Feshbach’s projection technique. The resulting dynamical equations have only three Faddeev components, just like the standard AGS formalism, while the mesonic channel has been projected out from Eq. (1). This recasting of the dynamics (which represents an exact result) leads to an additional modification of the driving term $`Z_{ss^{}}`$, due to the effects of the coupling to the fourth component, $$Z_{ss^{}}=Z_{ss^{}}^{AGS}+Z_{ss^{}}^\pi +Z_{ss^{}}^{^{}}.$$ (28) (Here and in the following, it must be assumed that the indices $`s`$, and $`s^{}`$ span only from 1 to 3, since the mesonic channel has been projected out.) The additional term, $`Z_{ss^{}}^{^{}}`$, is given by $$Z_{ss^{}}^{^{}}=Z_{s0}^\pi \tau _0^{(2)}Z_{0s^{}}^\pi ,$$ (29) while $`Z_{00}^\pi `$ is identically zero as has been pointed out in the introduction. We consider this new additional contribution to the driving term, $`Z_{ss^{}}^{^{}}`$, and discuss what kind of irreducible 3N diagrams are involved. This can be accomplished by considering the definition of $`Z_{0s}^\pi `$ given in Eq. (27), and taking into account the fact that $`\tau _0^{(2)}`$ represents the strength of the connected 3N correlations while the pion is “in flight”, as discussed in Ref. . A strong 3N correlation would correspond to a pole-like structure for $`\tau _0^{(2)}`$, shifted by the energy carried by the exchanged pion. Using the result expressed in Eq. (27) one can show that the term $`Z_{ss^{}}^{^{}}`$ can be written as $`Z_{ss^{}}^{^{}}=`$ $`{\displaystyle \underset{b,c,b^{},c^{}}{}}\stackrel{~}{s}^{(2)}\left\|g_0\left(f_s^{}\right)_{b^{}b}G_0\right\|b^{(3)}\tau _b^{(3)}\|(0^{(2)})_{a^{}b}\tau _0^{(2)}`$ (30) $`\times (0^{(2)})_{a^{}c}\|\tau _c^{(3)}c^{(3)}\left\|G_0\left(f_s^{}\right)_{c^{}c}g_0\right\|\stackrel{~}{s}^{(2)},`$ (31) and the diagram shown in Fig. 5 is just one contribution to this structure. Note that the pion can couple any of the three ingoing nucleon lines with each one of the three outgoing lines. Moreover this occurs with all possible recombinations in intermediate 2N pairs. This set of diagrams represent another, new class of irreducible 3N mechanisms contributing to the construction of the 3N force. Physically, these diagrams represent all possible connected correlations among the three nucleons while the exchanged pion is in flight. On the contrary, the diagrams of Fig. 2 represent all possible disconnected 2N correlations while the pion is in flight. ## IV Summary and Conclusions Recently, a new approach for the explicit treatment of the pion dynamics in the 3N system has been obtained . Herein, we have shown how to derive from this formulation a practical calculation scheme which is phenomenologically sound. The approximation is based on the assumption that the elastic 2N subamplitudes can be conveniently described with the phenomenological 2N potential approach. In other words, instead of treating explicitly the pion dynamics in full, in the 3N system we describe explicitely only those aspects of the pion dynamics which cannot be buried into the all-comprehensive, phenomenological 2N potential. Thus, the procedure could be viewed as a method to cool down (or to gradually project out) the pion dynamics from the theory of Ref. . One advantage of making such an approximation is that this approach is not plagued with the nucleon-renormalization problem, since the dynamical equation used to construct the subsystem amplitudes is represented by the standard 2N Lippmann-Schwinger equation, where the nucleon masses are static. However the approach includes also the pion degrees of freedom, via the inelastic subamplitudes. These are represented as 2N form factors defined in the Yakubovskĭ chain-labelled space of the (four-body) $`\pi `$NNN system. Here, we consider only the first-order contributions to such inelastic form factors in terms of the effective $`\pi `$NN coupling vertex, see Tab.I. (It is evident that more complex rescattering mechanisms contributing to such form factors can be implemented at a later stage.) Another advantage of this approach is that the dynamical input can be constrained by the 2N experimental data; therefore the method can be directly compared with the standard quantum-mechanical approach to the 3N problem based upon a 2N potential description. By using a separable-expansion representation of the elastic NN $`t`$-matrix, it has been possible to recast the dynamical equation into an extended AGS form, where the part of the pion dynamics not buried in the 2N potential is treated explicitly. We discussed the three modifications implied by this extended 3N equation with respect to the normal AGS one. First, the 3$`\times `$3 AGS driving term, defined in terms of the three Faddeev components, acquires additional contributions from the pion dynamics, and these new contributions act in both diagonal and off-diagonal matrix elements (while the standard AGS driving term is known to act only in the off-diagonal elements). We discussed these contributions also in terms of their diagrammatic interpretation, and found that the off-diagonal corrections correspond, to their lowest order, to irreducible 3N-force diagrams where the pion, while being exchanged between two nucleons, rescatters from the third before annihilating. This rescattering mechanism is practically the only case discussed in the construction of the pionic part of the 3N force. And the different 3NF approaches differ mainly for the model representation of the $`\pi `$N rescattering amplitude; there are, of course, additional short-range effects where there is much more ambiguity and where the various 3NF models differ considerably among each other. Conversely, the second modification refers to the diagonal part of the driving term, and represents irreducible 3N-force diagrams of different topology, where it is one of the two nucleons exchanging the pion that rescatters with the third one. These 3NF diagrams are not included in the construction of modern 3N potentials because it is usually assumed that they cancel out if one takes into account meson retardation effects in all their possible variety of time orderings. However, as has been pointed out in a recent study , this cancellation as a 100% effect is questionable. The reason for questioning the cancellation is due to the fact that a full 2N sub-amplitude enters in this 3N-force diagram, while the cancellation involves only the “instantaneous” part of the diagram. The construction of 3NF models from meson-exchange mechanisms evaluated in terms of instantaneous processes will lead inevitably to a 100% cancellation effect; however such methods do not take into account that two nucleons may clusterize while the pion is being exchanged, and this implies that an entire set of rescattering processes have to be subsummed during the pion-exchange process, thus leading to an incomplete cancellation effect. The third and last modification with respect to the normal AGS equation implies the increasing of the matrix dimension by one unit, since the 3 Faddeev components are now coupled to the additional two-cluster partition $`\pi `$+(NNN). We have discussed the structure of such couplings and the corresponding diagrams, thus providing details for all the ingredients of this new dynamical equation. Finally, we have shown that it is possible to project out the effects of the coupling to this pionic channel, thus providing a 3N dynamics with explicit treatment only for the 3N coordinates. This modification is suited for treatments of the 3N system at lower energies, below the pion threshold. Then, the extended AGS equation involves only the three standard Faddeev components, and the effect of the $`\pi `$+(NNN) channel in the intermediate states is contained into an additional, third contribution to the driving term. We have analyzed the diagrams involved and found that they correspond to a third class of irreducible 3NF diagrams, topologically different from the other, previously discussed two classes. The diagrams correspond to the variety of connected correlations amongst the three nucleons while the meson is in flight. In other words, these mechanisms take into account all possible 3N-cluster effects while the meson is being exchanged. This new correction acts in both diagonal and off-diagonal matrix elements of the driving term and may possibly influence the structure of the 3N force. ###### Acknowledgements. L.C. acknowledges funds from the Italian Murst-PRIN Project “Fisica Teorica del Nucleo e dei Sistemi a piú corpi”, and hospitality and financial support from the University of Manitoba, during two visits in 1999. J.P.S. acknowledges continuing financial support from NSERC, Canada. T.M. is grateful to the University of Manitoba for scholarship support. J.P.S. and T.M. are thankful to INFN, Padova, for hospitality and financial support during recent visits to Padova.
warning/0004/hep-th0004185.html	ar5iv	text	# 1 Introduction ## 1 Introduction The large $`N`$-limit of super-conformal gauge field theories has been proven to have dual representations in terms of the weakly coupled super-gravity or string theory via the Maldacena’s conjecture . In such a duality the radius of the space $`R_c`$ behaves like $`\mathrm{\Delta }^{1/4}`$, where $`\mathrm{\Delta }=g_{YM}^2N`$ is the ’t Hooft parameter which is fixed in the large-$`N`$ limit . Naturally for large $`\mathrm{\Delta }`$ this leads to a large radius which allows us to describe strings propagating in the background fields . This is a particular example of an idea suggested by Polyakov about the description of gauge theories in $`D`$ dimensions in terms of certain non-critical string theory in $`D+1`$ dimensions, where an extra dimension is due to the Liouville field. In this approach the properties of the geometry give us information about the properties of the gauge theory. On the other hand it should be possible to describe the weakly coupled regime of gauge theories by the strong coupling regime of string theory. One of the simplest questions that one can ask is about the derivation of the RG flow in gauge theories from the geometry, including the weakly coupled regime. Here we are not going to deal with this in the full $`10`$-dimensional theory. Instead of that we are going to study an alternative approach due to Álvarez and Gómez where the idea is to model the renormalization group equations of gauge theories. Basically their proposal implies to associate to the couplings of the gauge field theories some background fields of a closed string theory. Then one demands that the string $`\beta `$-functions for those backgrounds have to coincide with the RG equations of the gauge theories. In this approach the geometry dictates the properties of the gauge theory. In particular, as it has been shown by Álvarez and Gómez, for one-loop $`\beta `$-function in pure gluodynamics the space-time curvature is a continuously increasing function of the running-energy scale from the IR to the UV limit. The curvature behaves like an inverse power of the coupling. It implies that the theory runs continuously from the strongly coupled regime to the weakly coupled one. In this framework they have calculated the Wilson loops and shown that both confinement and over-confinement occur depending only on the area of the world-sheet of the fundamental string, at every value of the coupling. In the present paper we address the following situation. Consider a theory with a $`\beta `$-function which has a pole at some energy scale $`\mathrm{\Lambda }`$ in the infrared. The existence of the pole leads to two branches, one corresponds to the super-strongly coupled regime while the other one is the asymptotically-free regime. The point $`\mathrm{\Lambda }`$ behaves like an infrared attractive point since the RG flow of the theory goes from the UV limit to the IR one in both branches . Once the Wilson loops are computed for this kind of theories, one obtains an interesting new result: in the super-strong coupling, in the infrared, there is a range of values of the coupling in which the theory only leads to confinement and not over-confinement. We have calculated that range. Starting from the bosonic string action it is possible to derive the one loop $`\beta `$-functions which yield the equations of motion of the background fields $`\beta ^\mathrm{\Phi }`$ $`=`$ $`{\displaystyle \frac{D26}{48\pi ^2}}+{\displaystyle \frac{\alpha ^{}}{16\pi ^2}}\left(4(\mathrm{\Phi })^24^2\mathrm{\Phi }R+{\displaystyle \frac{1}{12}}H^2\right)+𝒪(\alpha _{}^{}{}_{}{}^{2}),`$ $`\beta _{\mu \nu }^G`$ $`=`$ $`R_{\mu \nu }{\displaystyle \frac{1}{4}}H_\mu ^{\lambda \sigma }H_{\nu \lambda \sigma }+2_\mu _\nu \mathrm{\Phi }+𝒪(\alpha ^{}),`$ $`\beta _{\mu \nu }^B`$ $`=`$ $`_\lambda H_{\mu \nu }^\lambda 2(_\lambda \mathrm{\Phi })H_{\mu \nu }^\lambda +𝒪(\alpha ^{}),`$ (1) where $`G_{\mu \nu }`$ and $`B_{\mu \nu }`$ are the symmetric and antisymmetric fields, respectively. $`R`$ is the scalar curvature and $`R_{\mu \nu }`$ is the Ricci tensor of the $`D`$-dimensional space-time. $`H_{\mu \nu \lambda }`$ is the antisymmetric tensor-field strength derived from $`B_{\nu \lambda }`$. These equations were derived using the world-sheet background-field perturbation theory from the non-linear sigma-model action plus the renormalizable action for the dilaton field $$S_{dilaton}=\frac{1}{4\pi }𝑑\sigma 𝑑\tau \sqrt{\gamma }R^{(2)}\mathrm{\Phi }(X).$$ (2) Here $`\gamma `$ is the determinant of the 2-dimensional world-sheet metric, $`R^{(2)}`$ is the scalar curvature of the world-sheet and $`\mathrm{\Phi }(X)`$ is the dilaton background field in the $`D`$-dimensional target space-time $`X`$. In the low-energy theory instead of strings one can deal with the background fields. Following Álvarez and Gómez we will solve the equations of motion which come from the vanishing of the $`\beta `$-functions, Eq.(1). For simplicity we will turn off all the fields but the symmetric tensor and the dilaton. We will work in the critical dimension ($`D=26`$) for which the equations reduce to those derived from the action of gravity coupled to the dilaton. Using the Liouville ansatz the equations of motion are trivially satisfied for any dilaton. In the metric we identify 4 of the 26 dimensions as the usual Euclidean (or Minkowski) space-time where the gauge theory is placed. One of the remaining $`22`$-spatial coordinates is identified as the running-energy scale ($`\mu `$) associated to the gauge field theories. In this prescription the gauge coupling is related to the dilaton by using the association<sup>4</sup><sup>4</sup>4The second part of this relation is an assumption related to the soft-dilaton theorem . $`g_s=e^\mathrm{\Phi }=g_{YM}^2`$. Since we know how the coupling constants run in gauge theories through the $`\beta `$-functions, the above correspondence determines the $`\mu `$-dependence of the dilaton. This procedure is known as the Renormalization Group approach to the string theory. See for example and references therein. We describe some features of the space-time and calculate the Wilson loops. There are two different types of $`\beta `$-functions that we will work on. The first is the one-loop $`\beta `$-function of pure gluodynamics where the coupling constant diverges in the infrared. The second case is a $`\beta `$-function which changes sign through a pole in the infrared. For instance we will consider in particular the Novikov-Shifman-Vainstein-Zakharov (NSVZ) $`\beta `$-function of $`𝒩=1`$ super-Yang-Mills theory, which has those properties. It was conjectured that because of this pole this theory has two phases . One is the super-strongly coupled phase and the other one is the asymptotically-free phase, both of which flow to the infrared attractive point at some finite scale. Finally in the discussions we mention about the theories with conformal-fixed points at the finite values of the coupling constants. In section 2 we study the properties of the background fields and the corresponding confining geometry derived from them. In section 3 the Wilson loops are studied by calculating the area of a fundamental string world-sheet in the background fields. In sections 4 and 5 we apply this framework to investigate the Yang-Mills type $`\beta `$-function and the $`𝒩=1`$ super-Yang-Mills type one, respectively. ## 2 The background fields Let us consider the RG equations of the bosonic strings . The vacuum configurations of string theory at one loop are determined by the sigma-model RG $`\beta `$-functions which, at leading order in $`\alpha ^{}`$, read as $$\beta _{\mu \nu }^G=R_{\mu \nu }+2_\mu _\nu \mathrm{\Phi },$$ (3) and $$8\pi ^2\beta ^\mathrm{\Phi }=\frac{D26}{6}\alpha ^{}^2\mathrm{\Phi }+2\alpha ^{}(\mathrm{\Phi })^2.$$ (4) Here $`_\mu \mathrm{\Phi }=_\mu \mathrm{\Phi }`$ and $`_\mu _\nu \mathrm{\Phi }=_\mu _\nu \mathrm{\Phi }\mathrm{\Gamma }_{\mu \nu }^\alpha _\alpha \mathrm{\Phi }`$. Conformal invariance implies that $`\beta _{\mu \nu }^G=\beta ^\mathrm{\Phi }=0`$. For all orders in $`\alpha ^{}`$-expansion these $`\beta `$-functions were shown to vanish . For the critical dimension these equations become $$R_{\mu \nu }+2_\mu _\nu \mathrm{\Phi }=0,$$ (5) and $$^2\mathrm{\Phi }2(\mathrm{\Phi })^2=0.$$ (6) As we have said before the four-dimensional space-time is embedded into $`26`$ dimensions and one of the remaining $`22`$-spatial coordinates is identified as the running-energy scale $`\mu `$. We will use the following form for the metric $$ds^2=a(\mu )(\pm dt^2+d\stackrel{}{x}^2)+b(\mu )d\mu ^2+c(\mu )d\stackrel{}{y}^2,$$ (7) where the $`a(\mu )`$ term is the Euclidean (Minkowski) metric in four dimensions, $`\mu `$ is the fifth-coordinate, while $`d\stackrel{}{y}^2`$ corresponds to an $`21`$-dimensional hyper-plane. The Ricci tensor is computed to be $$R_{ij}=\left(\frac{a^{}b^{}}{4b^2}\frac{21a^{}c^{}}{4bc}\frac{a^2}{2ab}\frac{a^{\prime \prime }}{2b}\right)\eta _{ij},$$ (8) for $`i,j=1,\mathrm{},\mathrm{\hspace{0.17em}\hspace{0.17em}4}`$, where we choose the mostly plus signature for the Minkowski metric, $`\eta _{ij}=(,+,+,+)`$, $`R_{55}`$ $`=`$ $`{\displaystyle \frac{a^2}{a^2}}{\displaystyle \frac{a^{}b^{}}{ab}}+{\displaystyle \frac{2a^{\prime \prime }}{a}}{\displaystyle \frac{21c^2}{4c^2}}{\displaystyle \frac{21b^{}c^{}}{4bc}}+{\displaystyle \frac{21c^{\prime \prime }}{2c}},`$ $`R_{\alpha \beta }`$ $`=`$ $`\left({\displaystyle \frac{a^{}c^{}}{ab}}{\displaystyle \frac{b^{}c^{}}{4b^2}}+{\displaystyle \frac{19c^2}{4bc}}+{\displaystyle \frac{c^{\prime \prime }}{2b}}\right)\delta _{\alpha \beta },`$ (9) for $`\alpha ,\beta =6,\mathrm{},\mathrm{\hspace{0.17em}\hspace{0.17em}26}`$, and the scalar curvature reads as $$R=\frac{2b^{}a^{}}{ab^2}\frac{21b^{}c^{}}{2b^2c}+\frac{189c^2}{2bc^2}+\frac{a^2}{a^2b}+\frac{4a^{\prime \prime }}{ab}+\frac{42a^{}c^{}}{abc}+\frac{21c^{\prime \prime }}{bc},$$ (10) where prime denotes derivative with respect to $`\mu `$. We will use the following ansatz to solve the equations of motion $`a(\mu )`$ $`=`$ $`e^{2\mathrm{\Phi }},`$ $`b(\mu )`$ $`=`$ $`4e^{4\mathrm{\Phi }}\mathrm{\Phi }^2,`$ (11) and $`c(\mu )=1`$. In this case the scalar curvature is $`R=e^{4\mathrm{\Phi }}=g_{YM}^8`$, where we have used $`e^\mathrm{\Phi }=g_{YM}^2`$. There is a naked singularity in the space-time which corresponds to the weakly-coupled regime of the gauge theory. On the other hand the strongly-coupled regime of the gauge theory corresponds to the weakly-curved spaces. This allows us to make a straightforward analysis of the geometry in terms of the quantum field theory $`\beta `$-functions. It is important to mention that we did not excite the tachyon field here. In the general case the metric will be more complicated. The geometry in this framework is universal in the sense that the metric and the scalar curvature depend only on the coupling constant. Choosing different types of theories, as long as the coupling diverges at certain point and goes to zero at another one, corresponds to essentially choosing different coordinates in this geometry. However once the gauge theory is introduced on a 4-dimensional hyper-plane and the prescription of computing the Wilson loops in terms of the minimal surfaces is given, the flow of the coupling constant becomes important and therefore different theories, in principle, can behave differently. With respect to this, and as we already mentioned before in the introduction, the important fact here is that the geometry can distinguish between one-loop $`\beta `$-function and $`\beta `$-functions with a pole at some finite value of the running-energy scale. ## 3 The Wilson loop In this section we discuss the calculation of the Wilson loops. First one should notice that the metric which we are dealing with can be trivially rewritten just by using the ansatz given in Eq.(11) $$ds^2=e^{2\mathrm{\Phi }}(\pm dt^2+dx_idx_i)+4l_c^2e^{4\mathrm{\Phi }}(d\mathrm{\Phi })^2+d\stackrel{}{y}^2,$$ (12) where $`x_i`$, $`i=1,\mathrm{\hspace{0.17em}2},\mathrm{\hspace{0.17em}3}`$, and $`l_c`$ is an arbitrary scale of dimension of length. In this framework $`l_c`$ is related to the running-energy scale of the field theory studied<sup>5</sup><sup>5</sup>5 Moreover one can also define $`\rho =e^{2\mathrm{\Phi }}`$ and rewrite the above metric in the Liouville form $`ds^2=\rho (\pm dt^2+dx_idx_i)+l_c^2(d\rho )^2+d\stackrel{}{y}^2`$. In this parameterization $`\rho =\mathrm{}`$ is the horizon and $`\rho `$=0 is the singularity. Therefore this space-time allows the description of zigzag invariant Wilson loops .. Since we want to deal with Wilson loops placed in the $`4`$-dimensional Euclidean (Minkowski) space-time, the extra 21 dimensions are irrelevant. Therefore this problem is equivalent to the problem of $`5`$-dimensional gravity coupled to the dilaton. Let us consider the Nambu-Goto action $$S_{NG}=\frac{1}{2\pi l_s^2}𝑑\sigma 𝑑\tau \sqrt{detG_{MN}_\alpha X^M_\beta X^N},$$ (13) where $`X^M`$ is a generic coordinate on the $`5`$-dimensional space-time. In the static configuration, for $`\tau =t`$ and $`\sigma =x`$ we have $$ds^2=\pm e^{2\mathrm{\Phi }}d\tau ^2+(e^{2\mathrm{\Phi }}+4l_c^2e^{4\mathrm{\Phi }}\mathrm{\Phi }_\sigma ^2)d\sigma ^2,$$ (14) where $`\mathrm{\Phi }_\sigma =\frac{\mathrm{\Phi }}{\sigma }`$. Therefore the action is $$S_{NG}=\frac{T}{2\pi l_s^2}_0^L𝑑\sigma e^{2\mathrm{\Phi }}\sqrt{1+4l_c^2e^{2\mathrm{\Phi }}\mathrm{\Phi }_\sigma ^2},$$ (15) where $`T`$ is the time. Basically we are considering the configurations as it is shown in figure 1. Figure 1: Schematic representation of the prescription of calculating the Wilson loops in terms of Liouville coordinates. In this figure we show the shape of the string world-sheet in terms of the Liouville field $`e^\mathrm{\Phi }`$. In particular $`e^{\mathrm{\Phi }_I}`$ is the limit when $`\mathrm{\Phi }\mathrm{}`$. Since the scalar curvature $`R`$ goes like $`e^{4\mathrm{\Phi }_I}`$ it implies that at this point the metric has a naked singularity (the scalar curvature blows up). The point at the origin represents the naked singularity. However in terms of the running-energy scale $`\mu `$ the location of this singularity depends on the particular field-theory $`\beta `$-function. We can put the $`4`$-dimensional hyper-plane $`\mathrm{\Sigma }`$ at any point which corresponds to choosing a particular value of the coupling constant. Here we denote a generic choice by $`e^{\mathrm{\Phi }_c}`$. Once the hyper-surface is chosen the string can fluctuate in either of the directions orthogonal to $`\mathrm{\Sigma }`$. However the minimal surfaces come from the strings that fluctuate in the direction of the naked singularity as it is depicted in the figure. In terms of the gauge theory language this means that semi-classically we allow the RG flow from strong coupling to the weak coupling regime. On the other hand, if an infrared attractive point occurs, the RG flow will be from the ultraviolet limit to the infrared one <sup>6</sup><sup>6</sup>6Observe that if one considers the Liouville coordinate to be time-like with a metric of the form $`ds^2=\rho (dt^2+dx_idx_i)l_c^2(d\rho )^2`$ the direction of RG flow is from weak coupling to strong coupling.. We will see that this is what happens for $`𝒩=1`$ super-Yang-Mills theory. Figure 2: The solid line indicates the scaled Wilson-loop size $`L/e^{\mathrm{\Phi }_c}`$ (in units of $`l_c`$) in terms of the dimensionless variable $`e^{\mathrm{\Phi }_0}/e^{\mathrm{\Phi }_c}`$. The dashed curve shows the derivative of the scaled Wilson-loop size. Taking the vertical axis as the $`\stackrel{}{x}`$-direction we will consider a Wilson-loop of size $`L`$. This means that $`L`$ is the static separation of the very heavy ”quark-anti-quark pair”. In the symmetric configuration $`\stackrel{}{x}=0`$ corresponds to a minimum of $`e^\mathrm{\Phi }`$, let us call it $`e^{\mathrm{\Phi }_0}`$, and since the dilaton is a function of $`\mu `$ we will call $`\mathrm{\Phi }_0=\mathrm{\Phi }(\mu _0)`$. For classical solutions one gets $$\frac{e^{4\mathrm{\Phi }}}{1+4l_c^2e^{2\mathrm{\Phi }}\mathrm{\Phi }_\sigma ^2}=e^{4\mathrm{\Phi }_0}.$$ (16) By inverting this expression $$\frac{L}{2}=2l_ce^{\mathrm{\Phi }_0}_1^{e^{\mathrm{\Phi }_c}/e^{\mathrm{\Phi }_0}}\frac{1}{\sqrt{v^41}}𝑑v,$$ (17) where $`v=\frac{e^\mathrm{\Phi }}{e^{\mathrm{\Phi }_0}}`$. After integration we obtain $$L=l_ce^{\mathrm{\Phi }_0}\left(\sqrt{\pi }\frac{\mathrm{\Gamma }(\frac{1}{4})}{\mathrm{\Gamma }(\frac{3}{4})}4\frac{{}_{2}{}^{}F_{1}^{}(\frac{1}{4},\frac{1}{2};\frac{5}{4};e^{4(\mathrm{\Phi }_0\mathrm{\Phi }_c})}{e^{\mathrm{\Phi }_c\mathrm{\Phi }_0}}\right),$$ (18) where $`{}_{2}{}^{}F_{1}^{}(\frac{1}{4},\frac{1}{2};\frac{5}{4};e^{4(\mathrm{\Phi }_0\mathrm{\Phi }_c)})`$ is a hypergeometric function . $`L`$ is a function of $`e^{\mathrm{\Phi }_c}`$ and $`e^{\mathrm{\Phi }_0}`$. It is convenient to study the function $`L/e^{\mathrm{\Phi }_c}`$ in terms of $`e^{\mathrm{\Phi }_0}/e^{\mathrm{\Phi }_c}`$. It has a maximum in the interval between $`0`$ and $`1`$ which indicates that there are two regions (I and II) to consider when one calculates the Wilson loops . For the interval below the maximum in the $`e^{\mathrm{\Phi }_0}/e^{\mathrm{\Phi }_c}`$-axis, region I, one considers large world-sheets. In figure 2 we plot $`L/e^{\mathrm{\Phi }_c}`$ in terms of $`e^{\mathrm{\Phi }_0}/e^{\mathrm{\Phi }_c}`$. There is a maximum at $`e^{\mathrm{\Phi }_0^M}/e^{\mathrm{\Phi }_c}0.62`$ and at this point $`L/e^{\mathrm{\Phi }_c}\|_M=0.42l_c`$. Let us calculate the energy for the static configuration. The Nambu-Goto action and the corresponding energy are related by $$E=\frac{S_{NG}}{T}=𝑑\sigma =\left(\frac{dv}{d\sigma }\right)^1𝑑v,$$ (19) where the lagrangian is given by $$=\frac{e^{2\mathrm{\Phi }}}{2\pi l_s^2}\sqrt{1+4l_c^2e^{2\mathrm{\Phi }}\mathrm{\Phi }_\sigma ^2}=\frac{1}{2\pi l_s^2}e^{2\mathrm{\Phi }_0}v^4.$$ (20) Then the integral becomes $$E=\frac{2l_c}{\pi l_s^2}e^{3\mathrm{\Phi }_0}_1^{\frac{e^{\mathrm{\Phi }_c}}{e^{\mathrm{\Phi }_0}}}\frac{v^4}{\sqrt{v^41}}𝑑v,$$ (21) which can be integrated as $$E=\frac{2l_c}{\pi l_s^2}e^{3\mathrm{\Phi }_0}\left(\frac{\sqrt{\pi }}{12}\frac{\mathrm{\Gamma }(\frac{1}{4})}{\mathrm{\Gamma }(\frac{3}{4})}+\frac{1}{3}e_2^{3(\mathrm{\Phi }_c\mathrm{\Phi }_0)}F_1(\frac{1}{2},\frac{3}{4};\frac{1}{4};e^{4(\mathrm{\Phi }_0\mathrm{\Phi }_c)})\right).$$ (22) Expanding Eqs.(18) and (22) in series of powers of $`e^{\mathrm{\Phi }_0}/e^{\mathrm{\Phi }_c}`$ we obtain $$\frac{L}{l_c}=\sqrt{\pi }e^{\mathrm{\Phi }_0}\frac{\mathrm{\Gamma }(\frac{1}{4})}{\mathrm{\Gamma }(\frac{3}{4})}4e^{2\mathrm{\Phi }_0\mathrm{\Phi }_c}+𝒪(e^{6\mathrm{\Phi }_05\mathrm{\Phi }_c}),$$ (23) and $$\frac{\pi l_s^2E}{2l_c}=\frac{e^{3\mathrm{\Phi }_c}}{3}+\frac{\sqrt{\pi }}{12}\frac{\mathrm{\Gamma }(\frac{1}{4})}{\mathrm{\Gamma }(\frac{3}{4})}e^{3\mathrm{\Phi }_0}\frac{1}{2}e^{4\mathrm{\Phi }_0\mathrm{\Phi }_c}+𝒪(e^{8\mathrm{\Phi }_05\mathrm{\Phi }_c}).$$ (24) Replacing $`L`$ in terms of $`e^{\mathrm{\Phi }_0}`$ in Eq.(24) we get an expression which shows that there is over-confining. $$E=\frac{1}{6\pi ^2l_s^2l_c^2}\left(\frac{\mathrm{\Gamma }(\frac{3}{4})}{\mathrm{\Gamma }(\frac{1}{4})}\right)^2L^3+\frac{2}{3\pi }\frac{l_c}{l_s^2}e^{3\mathrm{\Phi }_c}.$$ (25) In the limit of extremely strong coupling, which corresponds to $`e^{\mathrm{\Phi }_c}\mathrm{}`$, the above expression becomes divergent. However the divergent part is independent of the size of the Wilson loop. In the context of AdS/CFT duality this kind of divergence was regularized by a mass renormalization or by considering the Legendre transform of the minimal area . In our case observing that zero-size Wilson loops diverge we will drop the $`L`$-independent divergent term and measure the energy with respect to the zero-size Wilson loops. In the extremely strong coupling limit $`L`$ becomes $$L_{\mathrm{}}=\sqrt{\pi }l_ce^{\mathrm{\Phi }_0}\frac{\mathrm{\Gamma }(\frac{1}{4})}{\mathrm{\Gamma }(\frac{3}{4})},$$ (26) while the energy for the static configuration reads as follows $$E_{\mathrm{}}=\frac{l_ce^{3\mathrm{\Phi }_0}}{6\sqrt{\pi }l_s^2}\frac{\mathrm{\Gamma }(\frac{1}{4})}{\mathrm{\Gamma }(\frac{3}{4})},$$ (27) so that the over-confining is shown by the relation $$E_{\mathrm{}}=\frac{1}{6\pi ^2l_s^2l_c^2}\frac{\mathrm{\Gamma }(\frac{3}{4})^2}{\mathrm{\Gamma }(\frac{1}{4})^2}L_{\mathrm{}}^3.$$ (28) Figure 3: Energy (in units of $`\frac{l_c}{l_s^2}`$) for the static configuration as a function of $`\frac{e^{\mathrm{\Phi }_0}}{e^{\mathrm{\Phi }_c}}`$, in a small interval around $`0`$. In figure 3 we plot the energy versus $`e^{\mathrm{\Phi }_0}/e^{\mathrm{\Phi }_c}`$ in a small interval around zero (region I). It shows the cubic behaviour of the potential. A quite different situation arises when the region II is analyzed. In this region $`e^{\mathrm{\Phi }_0}`$ is close to $`e^{\mathrm{\Phi }_c}`$. In such a case we have to do the integration between $`1`$ and $`1+ϵ`$ so that we have $$\frac{L_ϵ}{2}=2l_ce^{\mathrm{\Phi }_0}\left(\sqrt{\pi }\frac{\mathrm{\Gamma }(\frac{5}{4})}{\mathrm{\Gamma }(\frac{3}{4})}(1+ϵ)_2F_1(\frac{1}{4},\frac{1}{2};\frac{5}{4};\frac{1}{(1+ϵ)^4})\right)=2l_ce^{\mathrm{\Phi }_0}\left(\sqrt{ϵ}+𝒪(ϵ^{\frac{3}{2}})\right),$$ (29) and $`E_ϵ`$ $`=`$ $`{\displaystyle \frac{2l_c}{\pi l_s^2}}e^{3\mathrm{\Phi }_0}\left({\displaystyle \frac{\sqrt{\pi }}{12}}{\displaystyle \frac{\mathrm{\Gamma }(\frac{1}{4})}{\mathrm{\Gamma }(\frac{3}{4})}}+{\displaystyle \frac{(1+ϵ)^3}{3}}_2F_1({\displaystyle \frac{1}{2}},{\displaystyle \frac{3}{4}};{\displaystyle \frac{1}{4}};{\displaystyle \frac{1}{(1+ϵ)^4}})\right)`$ (30) $`=`$ $`{\displaystyle \frac{2l_c}{\pi l_s^2}}e^{3\mathrm{\Phi }_0}\left(\sqrt{ϵ}+𝒪(ϵ^{\frac{3}{2}})\right).`$ At order $`\sqrt{ϵ}`$ it gives $$E_ϵ=\frac{e^{2\mathrm{\Phi }_0}}{2\pi l_s^2}L_ϵ,$$ (31) indicating confinement at region II. Figure 4 shows the whole picture of $`\frac{E}{e^{3\mathrm{\Phi }_c}}`$ (solid line) and $`\frac{L}{e^{\mathrm{\Phi }_c}}`$ (dashed line), as functions of $`e^{\mathrm{\Phi }_0}/e^{\mathrm{\Phi }_c}`$. In the region II we observe the area law ($`SLT`$). pt Figure 4: $`\frac{E}{e^{3\mathrm{\Phi }_c}}`$ (solid line) and $`\frac{L}{e^{\mathrm{\Phi }_c}}`$ (dashed line), as functions of $`e^{\mathrm{\Phi }_0}/e^{\mathrm{\Phi }_c}`$. ## 4 One-loop $`\beta `$-function In this section we will study the one-loop $`\beta `$-function of pure gluodynamics and compute the Wilson loops in this theory. This case was studied in . The $`\beta `$-function is $$\mu \frac{dg}{d\mu }=\frac{11N}{24\pi ^2}g^3.$$ (32) The dilaton field is $$\mathrm{\Phi }(\mu )=\mathrm{log}\mathrm{log}\left(\frac{\mu }{\mathrm{\Lambda }}\right)\mathrm{log}\left(\frac{11N}{24\pi ^2}\right),$$ (33) where $`\mathrm{\Lambda }`$ is a renormalization scale. The solutions for the metric are $$a(\mu )=\frac{576\pi ^4}{121N^2}\frac{1}{\mathrm{log}^2(\mu /\mathrm{\Lambda })},$$ (34) and $$b(\mu )=\frac{3^4\mathrm{\hspace{0.17em}2}^{14}\pi ^8}{(11N)^4}\frac{1}{\mu ^2\mathrm{log}^6(\frac{\mu }{\mathrm{\Lambda }})}.$$ (35) The curvature reads as $$R(\mu )=\left(\frac{11}{3}\right)^4\frac{N^4}{2^{12}\pi ^8}\mathrm{log}^4(\mu /\mathrm{\Lambda }).$$ (36) Figure 5 shows a picture for the Yang-Mills coupling constant in terms of the running energy scale $`\mu `$. Figure 5: Coupling constant for pure gluodynamics as a function of the running-energy scale $`\mu `$. The picture on the right shows the qualitative behaviour of the four-dimensional space-time. The picture on the right shows the behavior of the space-time radius (inverse of the curvature) as $`\mu `$ approaches to the IR limit. At $`\mu =\mathrm{\Lambda }`$ the space-time becomes flat. In the weak coupling there is a naked singularity. We will use the previous analysis for the Wilson loop calculation. In figure 6 we show the shape of the Wilson loop. The $`4`$-dimensional hyper-plane $`\mathrm{\Sigma }`$ where the Wilson loop is drawn can be placed at any value of $`\mu `$, let us say $`\mu _c`$, while the minimum $`e^{\mathrm{\Phi }_0}`$ corresponds to $`\mu _0`$. The naked singularity, represented by a dot, is the point at infinity and it is labeled as $`\mu _I`$. The corresponding point to the maximum in figure 4, i.e. $`\frac{e^{\mathrm{\Phi }_0^M}}{e^{\mathrm{\Phi }_c}}`$, is for this case $`\mu _0^M\mathrm{\Lambda }\left(\frac{\mu _c}{\mathrm{\Lambda }}\right)^{1.61}`$. For larger world-sheets it leads to over-confinement $$E=\frac{1}{6\pi ^2l_s^2l_c^2}\left(\frac{\mathrm{\Gamma }(\frac{3}{4})}{\mathrm{\Gamma }(\frac{1}{4})}\right)^2L^3.$$ Note that we have dropped the divergent part. While for smaller world-sheets we get $$E_ϵ=\frac{288\pi ^3}{121N^2l_s^2\mathrm{log}^2(\frac{\mu _0}{\mathrm{\Lambda }})}L_ϵ,$$ (37) indicating the confinement in region II. The factor in front of $`L_ϵ`$ can be interpreted as an effective tension. Figure 6: Schematic representation of the Wilson loop calculation for pure gluodynamics. ## 5 $`\beta `$-function with a pole On the other hand, for a $`\beta `$-function with a pole at some finite value of the coupling the situation is quite different. For instance let us consider the $`𝒩=1`$ super-Yang-Mills theory. In this case we will study NSVZ $`\beta `$-function $$\beta (g)=\mu \frac{dg}{d\mu }=\frac{g^3}{16\pi ^2}\frac{(𝒩4)N}{\left(1\frac{(2𝒩)Ng^2}{8\pi ^2}\right)},$$ (38) where $`N`$ is the corresponding label to the gauge group $`SU(N)`$ while $`𝒩`$ labels the number of super-symmetries. By integrating this equation one gets a transcendental equation $$\frac{\mu }{\mathrm{\Lambda }}=e^{\frac{8\pi ^2}{(4𝒩)Ng^2}}\left(\frac{g^2}{4\pi }\right)^{\frac{2𝒩}{4𝒩}}.$$ (39) For $`𝒩=1`$ super-Yang-Mills theory the $`\beta `$-function is exact both perturbatively and non-perturbatively and it reads $$\beta (g)=\mu \frac{dg}{d\mu }=\frac{3}{16\pi ^2}\frac{g^3N}{\left(1\frac{Ng^2}{8\pi ^2}\right)}.$$ (40) This $`\beta `$-function has a pole at $`g^2=8\pi ^2/N`$ and it changes sign through the pole. As it is shown in figure 7 the coupling constant is a double-valued function of $`\mu `$. The pole is an infrared attractive point. Figure 7: Coupling constant for $`𝒩=1`$ super-Yang-Mills theory as a function of the running-energy scale $`\mu `$. The picture on the right shows the qualitative behaviour of the four-dimensional space-time. The theory can flow, both from the asymptotically-free phase where the coupling is small at large $`\mu `$ (lower branch ($``$)), and from the super-strongly coupled phase where the coupling is large (upper branch ($`+`$)), to the infrared attractive point. Since we can not invert the transcendental equation (39) we will analyze the behaviour around the infrared point $`\mathrm{\Lambda }`$. By expanding the Eq.(39) around $`\mathrm{\Lambda }`$ one gets $$g_\pm ^2=\frac{8\pi ^2}{N}\left(1\pm \sqrt{3\frac{\mu \mathrm{\Lambda }}{\mathrm{\Lambda }}}\right).$$ (41) Since $`\mathrm{\Phi }=\mathrm{log}(g^2)`$ it follows that $$\mathrm{\Phi }_\pm =\mathrm{log}(8\pi ^2/N)+\mathrm{log}(1\pm \zeta ),$$ (42) where $`\zeta =\sqrt{3(\mu \mathrm{\Lambda })/\mathrm{\Lambda }}`$. The solutions are $`a_\pm (\mu )`$ $`=`$ $`{\displaystyle \frac{64\pi ^4}{N^2}}(1\pm \zeta )^2,`$ $`b_\pm (\mu )`$ $`=`$ $`9{\displaystyle \frac{2^{12}\pi ^8}{N^4\mathrm{\Lambda }^2}}{\displaystyle \frac{(1\pm \zeta )^2}{\zeta ^2}}.`$ (43) We set $`c_\pm (\mu )=1`$. The scalar curvature is given by $`R_\pm `$ $`=`$ $`{\displaystyle \frac{N^4}{2^{12}\pi ^8(1\pm \zeta )^4}}\zeta <<1.`$ (44) In figure 8 we show the picture of the Wilson loop for the present case. Notice that we use the same coordinates as in figure 1. In these coordinates the picture becomes more clear. The vertical line at the point $`e^{\mathrm{\Phi }_\mathrm{\Lambda }}`$ represents the infrared attractive point. Here it acts like an effective horizon. This is because in the super-strongly coupled phase the direction of the RG flow is from strong coupling to the weak coupling, however since we have an infrared attractive point we can not continue the flow to the smaller couplings. On the other hand in the asymptotically-free phase the flow is from weak coupling to the strong coupling. The singularity placed at the origin corresponds to the limit of weak coupling shown in the lower-branch of figure 7. Again confinement and over-confinement have to be understood as properties of the world-sheet size. Close to effective horizon, for both branches we have $$E_{ϵ\pm }=\frac{32\pi ^3\left(1\pm \sqrt{3(\mu _0\mathrm{\Lambda })/\mathrm{\Lambda }}\right)^2}{N^2l_s^2}L_ϵ.$$ (45) In figure 8 the left side of the vertical line represents the weakly-coupled phase and the right side represents the super-strongly coupled phase. Figure 8: Schematic representation of the Wilson loop for $`𝒩=1`$ super Yang-Mills theory. The dashed ellipses represent the behaviour of the $`4`$-dimensional space-time as it has seen in the right picture of figure 8. When one studies the theory in the super-strongly coupled phase and sufficiently close to the attractive point one always gets the area law for the Wilson loops. In fact there is a critical coupling $`g_{critical}=1.27g_\mathrm{\Lambda }`$. For all values of the coupling between $`g_{critical}`$ and $`g_\mathrm{\Lambda }`$ we have the area law. ## 6 Discussion and conclusions In this paper, using the RG approach to string theory, we studied two types of confining theories. As it was shown by Álvarez and Gómez a theory which has a $`\beta `$-function with a zero at UV-limit is confining or over-confining, depending on the minimal surfaces we use. We have considered the theories which have $`\beta `$-functions with a pole at some finite value of the coupling. For example $`𝒩=1`$ super-Yang-Mills theory is of this form. It has two phases, the super-strongly coupled phase and the asymptotically-free phase which flow to an infrared attractive point. In this theory one has also both confinement and over-confinement. We showed that in the super-strongly coupled regime of the theory there is a critical value below which one has confinement only. In this regime large world-sheets are not allowed and one does not have over-confinement. It is interesting to stress that this fact is due to the existence of a pole at some finite scale in the $`\beta `$-function. This leads to a new scale $`g_{critical}=1.27g_\mathrm{\Lambda }`$ below which these theories satisfy the area law for their Wilson loops. Concerning to the flow of the space-time for the one-loop $`\beta `$-function we have seen that it goes from the strong coupling to the weak coupling regime. In the case of $`\beta `$-functions with a pole the situation changes in the sense that now the RG flow goes from the UV limit to the IR limit, in both branches. The geometry we have considered is universal in the sense that it leads to confinement. So presumably this geometry is not suitable for conformal theories and Abelian theories. For the conformal case we know that the space-time should be anti-de Sitter. However let us assume that, by using the metric in this paper, we can describe theories with conformal-fixed points. The $`\beta `$-function changes sign through zero instead of a pole but the geometry and the analysis of the Wilson loops will be similar to the case of $`\beta `$-functions with a pole which was discussed in section 5. Although this is surprising, an argument given by Damgaard and Haangensen may help one to understand this. They showed a theory with a self-dual point is either conformal at the self-dual point or its $`\beta `$-function has a pole at this point. Here we reproduce their argument. Let us consider a theory with a coupling $`g`$. If this theory is self-dual there is a relation between $`g`$ and the coupling of its dual theory $`g^{}`$ $$g^{}=f(g).$$ (46) The only essential requirement is that the interaction part in the action and its dual have to be of the same form. The map in Eq.(46) is assumed to be one-to-one and $`ff=1`$, which implies that its derivative $`\frac{f}{g}`$ is negative. We also assumed that $`g`$ and $`g^{}`$ are positive-valued functions. At the self-dual point, let us call it $`g_{self}`$ we have $$g_{self}=f(g_{self}).$$ (47) On the other hand we know that the RG flow is dictated by the $`\beta `$-function $$\beta (g)=\mu \frac{g}{\mu },$$ by applying the operator $`\mu \frac{}{\mu }`$ to the map of Eq.(46) we obtain a consistency relation $$\mu \frac{g^{}}{\mu }=\mu \frac{g}{\mu }\frac{g^{}}{g}=\beta (g)\frac{f}{g},$$ (48) if one identifies the left hand side of the above equation with the $`\beta `$-function for the dual theory one obtains $$\beta (g^{})=\beta (g)\frac{f}{g}.$$ (49) From the condition $`ff=1`$ at the self-dual point results that $`\frac{f}{g}=1`$, and therefore there are two solutions of Eq.(48), i.e. the $`\beta `$-function can be zero at this point or it is discontinuous. The most natural way to realize the discontinuity is through a pole in the $`\beta `$-function. Then one has at the self-dual point $`\beta (g_{self}+ϵ)=\beta (g_{self}ϵ)`$. Finally we would like to speculate that when one deforms $`𝒩=4`$ theory down to $`𝒩=1`$ to obtain a dual description in terms of string theory or gravity, a prescription which works only in the strong coupling regime of the gauge theory, one might be dealing with the super-strongly coupled regime, the upper branch, of $`𝒩=1`$ theory. ## Acknowledgments We benefited from useful discussions with Alex Kovner. M.S. also acknowledges discussions with Gastón Giribet, Esteban Moro and Carlos Núñez, and the organizers of the Spring Workshop on Superstrings and Related Matters (2000) at the Abdus Salam ICTP for kind hospitality where a part of this work was carried out. The work of I.K. and B.T. was supported by PPARC Grant PPA/G/O/1998/00567. The work of M.S. was supported by the CONICET of Argentina, the Fundación Antorchas of Argentina and The British Council.
warning/0004/astro-ph0004040.html	ar5iv	text	# Optical imaging and spectroscopy of BL Lac objects. ## 1 Introduction In the past decade BL Lac objects have been actively investigated in direct imaging and spectroscopy using ground based telescopes and HST.The imaging effort has been directed towards detecting the host galaxy, and when possible towards measuring its absolute luminosity and colors and determining its morphological properties. The aim of the spectroscopy has been to measure the redshift of the host or to measure the redshift of companions galaxies in order to assess a possible group or cluster membership. Apart from studies on individual objects a number of papers have presented optical images for samples or lists of objects. Twenty three objects have been imaged with the William Herschel Telescope in the R filter and 14 are resolved (Abraham et al. 1991). However due to either unknown redshift or poorly detected nebulosity only for 6 sources absolute quantities are derived. Some cases of disc dominated host galaxies are proposed. Sixteen objects in the southern sky have been studied by Falomo (1996) using sub-arcsec images obtained at the ESO 3.5m New Technology Telescope (NTT). Eleven sources were resolved and the hosts found to be luminous ellipticals (M$`{}_{R}{}^{}`$ –23.5). For a number of objects close companion galaxies are detected. Due to their small projected distance it is likely that they are associated with the BL Lac but spectroscopy is needed to assess this point. A larger sample but with poorer average resolution was investigated using the 3.6m CFHT (Wurtz et al. 1996). Fifty objects have been observed and 36 well resolved. For another ten objects the host galaxy has been only marginally detected. No difference of host properties is found between objects discovered in radio surveys (i.e. 1Jy sample) and those derived from X-ray surveys (i.e. EMSS). With very few exceptions all the BL Lac objects investigated are classified as ellipticals based on the surface brightness profiles. More recently a study of the host galaxies in a large sample of X-ray selected (high frequency peaked) BL Lacs have been presented by (Falomo & Kotilainen 1999). They used high resolution images in the R filter at the Nordic Optical Telescope (NOT) to image 52 targets from EMSS and Einstein Slew samples. All the 45 objects resolved are well represented by elliptical models. On average the hosts are found 1 magnitude more luminous than M (M$`{}_{R}{}^{}`$ -22.5; Mobasher et al. 1993; assuming R-K = 2.7). In addition to ground based studies several 0.1 arcsec resolution short exposure images have been obtained with WFPC2 camera on board of HST during a snapshot survey (Scarpa et al. 2000; Urry et al. 2000). Objects from various samples, and in the redshift interval 0.05 $`<z<`$ 1.3, were observed and 69 out of 110 observed are resolved. The highest redshift host galaxy detected is $`z=0.664`$ for 1823+568. For 80% of the resolved host galaxies an elliptical model is clearly preferred over a disc galaxy. The median absolute magnitude of these host galaxies ($`M_R23.7`$) is at least one magnitude brighter than $`M^{}`$. The nuclei are always well centered over the body of the galaxy and have luminosity similar to that of its host galaxy. From the point of view of the optical morphology the hosts of BL Lacs appear indistinguishable from “normal” (non active) ellipticals. The main aim of all these observations outlined above was to detect the host galaxies and to determine their structural and photometric properties. The knowledge of the kind of galaxies that host a BL Lac phenomenon in the nucleus is of importance not only for understanding/studying the nuclear activity vs galaxy connection (see e.g. Lawrence 1999) but also as a probe to test unified models of radio loud AGN. In particular if BL Lacs are FR I radio galaxies whose jet is aligned along the line of sight (e.g. Urry & Padovani 1995; Ulrich 1989) their host galaxies should exhibit exactly the same photometrical and morphological properties as the hosts of FR I. The properties of the BL Lacs hosts can also be compared with those of related beamed objects such as FSRQ and HPQ (see e.g. Kotilainen et al. 1998a). The aim of this work is to complement the existing data on BL Lac host galaxies and close environment with new imaging and spectroscopy for a dozen of (previously not well studied) objects. A general discussion and comparison of the properties of BL Lacs and radio galaxies will be presented elsewhere. In this paper we therefore present results from optical images of BL Lac objects collected at the NTT with mostly sub-arcsec resolution. Most of the objects presented here were not previously investigated with adequate capabilities. These observations therefore complement the existing data on BL Lac host galaxies. We also present spectroscopic observations for some of the objects performed with the aim of deriving the redshift of the host galaxies and of some nearby companion galaxies. When no spectroscopic redshift is available we give an estimate of the photometric redshift derived by assuming that the host has M<sub>R</sub> = –23.85 and R<sub>e</sub> = 9 kpc (the typical median values found in previous studies of BL Lacs hosts; e.g. Falomo & Kotilainen 1999). In Sect. 2 we describe the observations and data analysis. Section 3 reports the results obtained for each individual objects. Section 4 gives a summary of the results and discussion. ## 2 Observations and data analysis Optical observations were obtained using the 3.5m New Technology Telescope (NTT) at the European Southern Observatory (ESO), operated via remote control from the ESO headquarters in Garching (Germany). We acquired images using the Superb Seeing Imager (SUSI; Melnick et al. 1992) which is installed at one of the Nasmyth foci of the NTT. Configuration used was R-band filter and a CCD (TK 1024) with 24$`\mu `$m pixel size corresponding to 0.13<sup>′′</sup> on the sky. Conditions were photometric and seeing was ranging from 0.55 to 1.2 arcsec (FWHM), and in most cases $`<`$ 1<sup>′′</sup> . Observations of standard stars (Landolt 1992) were used to set the photometric zero point. We obtained images centered on the BL Lac object with exposure times ranging from 10 to 30 min (see Table 1). For many objects we also secured one short (2 minutes) exposure in order to be sure to get unsaturated images of the nucleus of the targets and to enable us to use bright stars in the field to study the PSF. The images were processed in the standard way (bias subtracted, trimmed, flat fielded, and cleaned of cosmic rays) using the Image Reduction and Analysis Facility (IRAF) procedures. A journal of the observations is given in Table 1. Spectroscopy of the the objects and/or of galaxies in the field were obtained for some targets in order to determine the redshift of BL Lacs and/or nearby companion galaxies. For this purpose the ESO multi mode instrument EMMI (Melnick et al. 1992) was used with red arm and grism elements. In general the slit has been oriented in order to obtain in a single observation both the BL Lac object and one or more galaxies around the source. All the images have been analyzed following the methods and procedure described in Falomo (1996). In particular surface photometry analysis was performed down to the surface brightness magnitude $`\mu _R`$ 26 mag./arcsec<sup>2</sup> in order to derive the properties of the host galaxies. A fit of the radial brightness profile was performed assuming a simple two model components: a point source plus a elliptical galaxy described by a de Vaucouleurs law $$I(r)=I_0exp\{7.67[(r/r_e)^{1/4}1]\}$$ where $`I(r)`$ is the surface brightness and $`r_e`$ the effective radius. Also disc galaxies models were attempted but in no cases they gave a better fit than the elliptical model. This is consistent with what was found in previous studies on a larger number of sources ( Falomo & Kotilainen 1999; Urry et al. 1999; Scarpa et al. 2000). To obtain absolute quantities we applied correction for Galactic extinction and redshift (K-correction). The former was determined using the Bell Lab Survey of neutral hydrogen N<sub>H</sub> converted to E<sub>B-V</sub> (Stark et al. 1992; Shull & Van Steenberg 1985), while the latter was computed from the model of Coleman et al. (1980) for elliptical galaxies. Throughout this paper, H<sub>0</sub> = 50 km s<sup>-1</sup> Mpc<sup>-1</sup> and q$`{}_{0}{}^{}=0`$ are adopted. ## 3 Results In Fig. 1 we report the observed radial brightness profile of the objects together with the best fit with the two components (point source plus elliptical galaxy) for the objects resolved. Parameters of the fit and absolute quantities for host galaxies and the nuclei are given in Table 2. In this Table columns 4–8 we give the results from this paper. The redshift in column 2 is drawn from literature except that for 0301-24 and two cases where a photometric redshift (given in parenthesis) is derived from the observed host properties. In the following discussion absolute quantities are given including corrections for galactic extinction and redshift (K-correction). Optical spectra of the BL Lacs or companion galaxies are reported in Fig. 2 together with the main identifications of observed spectral features. ### 3.1 Comments for individual objects PKS 0138-097 This object was observed under 1.2<sup>′′</sup> seeing and it looks unresolved. Heidt et al. 1996 have presented deep sub-arcsec images of this source that indicate the presence of close companion objects. These could be responsible for the intervening absorption system at z = 0.501 (Stickel et al. 1993) seen in the spectrum of the BL Lac object. Our image was taken under relatively poor seeing but nevertheless some evidence of the southern feature at $``$1.5<sup>′′</sup> from the center of the source is present in our image. This object has also been imaged by HST and found to be unresolved (Scarpa et al. 2000) but the presence of a companion galaxy at 1.5<sup>′′</sup> South from the nucleus is clearly apparent. Recent spectroscopy (Stocke & Rector 1997) detects for the first time the emission-line redshift of z=0.733 based upon weak Mg II and \[O II\] emission features. At this relatively high redshift our image result is consistent with this object being in a luminous (not detected) host galaxy at z = 0.733. 0301-243 We took a 20 minute image under good seeing ( 0.8<sup>′′</sup> ) of this BL Lac object that clearly shows an extended nebulosity (ellipticity $`\epsilon `$ = 0.3; $`ϵ`$ = 1- $`b/a`$ ) with a complex close environment (see Fig. 3). The immediate region around the object is rich with faint galaxies and there is a marked enhancement of the galaxy density within $``$60<sup>′′</sup> from the BL Lac object. The spectra of two galaxies (G1 and G2; see Fig. 3) at $``$ 6<sup>′′</sup> and 20<sup>′′</sup> from 0301-243 indicate that they are at $`z=0.263`$ suggesting a cluster of galaxies of Abell richness class 0 might be associated with the BL Lac source at this redshift (Pesce et al. 1995). The radial profile is adequately well represented by a point source plus the elliptical model while the fit with an exponential disk is not acceptable. Fig. 3 (right panel) shows the field after subtraction of the BL Lac model (nucleus plus host galaxy) revealing the faint galaxy $``$ 3.5<sup>′′</sup> South of the nucleus. After masking out the companion from the image we find that the surrounding nebulosity is very well centered on the nucleus within an accuracy of 0.2<sup>′′</sup> . We took three optical spectra of the nebulosity with the slit off the nucleus by 2<sup>′′</sup> . They are still dominated by the signal from the non-thermal source but all three show one weak emission line at $`\lambda `$ = 6303 Å (see Figure 2). The most plausible identification for this emission is \[ O III\] 5007 Å that yields a redshift of 0.26. Fainter emissions like \[ O III\] 4959 Å or H<sub>β</sub> could be present at this $`z`$ but not detectable in our spectrum because the features are lost in the noise. Other possible identifications like MgII 2800 (at z = 2.25) are not acceptable because the host galaxy would be too luminous (M $`<`$ –30). The redshift of 0.26 is very similar to the redshift of the companion galaxies G1 and G2 (respectively of M<sub>R</sub> = -20.7 and M<sub>R</sub> = -22.3) and supports the idea that the host of the BL Lac is the dominant member of a cluster of galaxies. We note that few other examples have been reported in the literature of BL Lacs in clusters whose membership has been proved spectroscopically. H0414+00 (Falomo et al. 1993a) is in a cluster of Abell class 0; PKS 0548-32 (Falomo et al. 1995) is in a cluster of Abell class 1-2. At this redshift (z = 0.26) the absolute magnitude of the host galaxy of 0301-24 is M<sub>R</sub> = -24.1 0338-21 Our image obtained with 0.6<sup>′′</sup> seeing shows the object to be resolved. This is the first detection of the nebulosity for the source. However its magnitude is not consistent with the redshift of the object published twenty years ago ( z = 0.048 ; Wright et al. 1977). The strongest absorption line identified in Wright et al. is indeed a telluric band at 6870 Å. Subsequent spectroscopy of the source has failed to confirm this redshift and a pure featureless optical spectrum has been observed (Falomo et al. 1994). In fact at this redshift of z = 0.048 the nebulosity would correspond to an unreasonably faint and small host galaxy (M$`{}_{R}{}^{}`$ -18.5). Assuming a typical host galaxy (see Sect. 1) we can well fit the radial brightness profile with a nucleus plus host galaxy obtaining a photometric redshift z $``$ 0.45. REX 0353-36 The source was identified as a BL Lac object in the REX survey of AGN (Wolter et al. 1997). Its optical spectrum is featureless (Wolter et al. 1998). We obtained an image under very good seeing (0.6 arcsec) and are able to detect the surrounding nebulosity and measure its luminosity and R<sub>e</sub>. This is the first detection of its host. There is no spectroscopic redshift but we can estimate a photometric redshift from the image decomposition assuming the host galaxy has average properties for BL Lacs hosts. The value of the photometric redshift so obtained is z $``$ 0.4. PKS 0548-322 This is a very well know BL Lac object at z = 0.068 (Fosbury & Disney, 1976) with a very large host galaxy in a rich environment (Falomo et al 1995). We took one relatively short exposure but good signal-to-noise spectrum centered in the nucleus to search for possible emission lines as have been reported in a number of nearby BL Lacs (e.g. BL Lac itself, Vermeulen et al. 1995 ).The spectrum, shown in Fig. 2, exhibits a substantial contribution from the stellar population of the host galaxy. The MgI 5175 Å and Na blend 5892 Å are well detected with equivalent widths of 12 Å and 6.5 Å, respectively. We could not find any emission down to a limit of equivalent with of 2 Å. This limit corresponds to H<sub>α</sub> line luminosity L(H<sub>α</sub>) $``$ 5 $`\times `$ 10<sup>40</sup> erg s<sup>-1</sup> which is about a factor 10 lower than the line detected in BL Lac (Vermeulen et al. 1995). PKS 0735+178 This BL Lac object is bright and strongly variable. It has been extensively studied in the radio range and several moving components have been detected in VLBI. The optical spectrum shows the absorption line due to an intervening system at 3980 Å, which if identified with Mg II gives z $`>`$ 0.424 (Carswell et al. 1974) Our images were obtained under seeing of 0.8<sup>′′</sup> but the source remains unresolved. Previous images were presented by Hutchings et al. (1988) who also found this source unresolved. There is no sign in our image (see Fig. 4) that the galaxy 7<sup>′′</sup> NW is distorted by interaction with 0735+178 as suggested by previous lower resolution images (Hutchings et al. 1988). The object was also imaged by Stickel et al. (1993) who are not able to detect the surrounding nebulosity. They obtained a spectrum of the galaxy 7<sup>′′</sup> NW and found z = 0.645. This BL Lac object is unresolved also in a short exposure image obtained with HST (Scarpa et al. 2000). In addition to the two well resolved companion galaxies we detect a faint emission at $``$ 3.5<sup>′′</sup> East from the BL Lac (see Fig. 4). Given its projected distance from the BL Lac (25 kpc at z = 0.424) it could be related to the intervening absorption at z = 0.424 but we cannot exclude that it is just a faint background source. From our image we can set a lower limit to the redshift (again assuming the typical properties for the host) of z $`>`$ 0.5, consistent with the limit derived from intervening absorption. PKS 0736+017 The excellent (seeing 0.55 arcsec) image (see Fig. 1) shows the flat spectrum radio quasar PKS 0736+01 (z = 0.191) as well as two close resolved faint companions that are embedded in the nebulosity of the object. The radial luminosity profile (see Fig. 1 ) is very well represented by an elliptical galaxy with a bright point source in the nucleus. It is found that the galaxy has M<sub>R</sub> = -24.3 and effective radius of $``$ 12 kpc. This host galaxy was previously detected in the optical with lower resolution by Wright et al. (1998). They derive M<sub>R</sub> = -22.0, which is substantially fainter than our value. We note that this discrepancy could be due to a problem in the Wright et al. image calibration as their surface brightness goes unbelievably faint. At 5<sup>′′</sup> from the nucleus their surface brightness is about $`\mu _R`$ =28 while our value at the same radius is $`\mu _R`$ = 24. The object has been also resolved in the NIR by Taylor et al. (1996) who found M<sub>K</sub> = -26.3, and by Kotilainen et al. (1998a) who found M<sub>H</sub> = -26.2. The R-H color turns out to be $``$ 2.0, consistent with the range of values reported by Kotilainen et al. 1998b for a number of BL Lacs. PKS 0754+10 We took two spectra of this BL Lac object for which no firm value of the redshift is available but whose host galaxy had already been detected (Abraham et al. 1991; Falomo 1996). The tentative redshift ($`z=0.66`$) proposed by Persic & Salucci (1986) based on inspection of the photographic spectrum reported by Wilkes et al. (1983) is unlikely as the host galaxy would be extremely luminous (($`M_R26`$ mag). Our spectra were obtained positioning the slit 2<sup>′′</sup> from the nucleus in order to reduce the contamination of light from the bright nucleus. Therefore the spectrum (see Fig. 2) is noisy and it is still dominated by the nuclear non thermal emission. We are not able to unambiguously identify spectral lines but some hint of the CaII break signature from the host galaxy is possibly apparent at $`\lambda `$ = 5045 Å which corresponds to z = 0.28. At this redshift the detected surrounding nebulosity would be M<sub>R</sub> $``$ -23. We note that this is consistent with the value of the redshift of the companion galaxy (see Fig. 5) 13.6<sup>′′</sup> north-east of the BL Lac object ($`z`$=0.27; Pesce et al. 1995) and could be another case of a companion galaxy physically associated with a BL Lac object. A definitive redshift determination is however still needed for this BL Lac object. PKS 0818-128 There is no redshift for this object and its optical spectrum is featureless (Falomo et al. 1994). Our optical images, obtained with seeing of 0.7<sup>′′</sup> , are not able to detect the host galaxy. The radial brightness profile is well matched by that of a scaled PSF (see Fig. 1). We can set a lower limit to the redshift assuming its nucleus is hosted by a standard luminous ( M<sub>R</sub> $``$ -23.8) elliptical. The limit of redshift we found for such a galaxy to be undetected in our image is z $`>`$ 0.5. In order to search for emission or intervening absorption line we gathered spectra in a wide wavelength range. Our spectrum ( see Fig. 2) is still dominated by the non-thermal featureless emission. The only feature (in addition to telluric bands) we can detect is an absorption at 6284 Å (e.w. 0.7 Å). The most likely identification of this feature is with an interstellar diffuse absorption band at the same wavelength. This is consistent with the low galactic latitude (b$`{}_{II}{}^{}`$ 13<sup>o</sup>) of the source. Alternatively the absorption line could be identified with MgII 2800 Å and this would yield approximately z $`>`$ 1.2 and, consequentially, the object would be extremely luminous (M$`{}_{R}{}^{}<29`$). PKS 0829+046 Previous images obtained at sub-arcsec resolution showed that the host galaxy (z = 0.18) has M<sub>R</sub> $``$ -23 (Falomo 1996). There is also an excess density of galaxies around this object (Pesce et al. 1994). But our spectroscopy shows that only some of them may be physically associated with the BL Lac object. Pesce et al. 1994 obtained the redshifts of galaxies G1 and G2 (see Fig. 6) at respectively z =0.24 and z = 0.204. We took additional spectra of two other galaxies (G3 and G4, see Fig.s 2 and 6). We found that G4 is at significantly higher redshift (z = 0.29) while G3 is at z = 0.175, consistent with being associated with PKS 0829+04 at projected distance of $``$ 120 kpc. In fact G3 is the only galaxy which is at the same redshift as the BL Lac. On one hand this is another case of similar redshift of a companion galaxy and its BL Lac. On the other hand the environment of 0829+04 must be less rich than what can be estimated from galaxy counts. H 1101-23 This is a BL Lac discovered from X-ray survey and is surrounded by a conspicuous rather elongated nebulosity (see Fig. 7 ) at the proposed $`z=0.186`$ (Remillard et al. 1989) confirmed by Falomo et al. (1994). The radial brightness profile extends to 15 arcsec along the major axis. We found that the luminosity profile is well fitted by an elliptical galaxy model plus a point source. The luminosity of the host galaxy is very high. The absolute magnitude, M<sub>R</sub> = –24.45, sets this galaxy among the brightest hosts of BL Lac objects (Falomo & Kotilainen 1999). For this object (see Fig. 8) we performed detailed surface photometry analysis using the AIAP package (Fasano 1990) in order to study the structural properties of the galaxy. From this analysis we derived photometric and structural parameters (surface brightness, ellipticity, position angle and Fourier coefficient C<sub>4</sub> describing the deviation of isophotes from the ellipse) as a function of the equivalent radius $`r=a\times (1ϵ)^{1/2}`$ where $`a`$ is the semi-major axis and $`ϵ`$ is the ellipticity of the ellipse fitting a given isophote. We found the ellipticity profile is increasing from the center outwards up to $`ϵ`$ = 0.45. The profile of the C<sub>4</sub> (see Fig. 9) shows disky (positive C<sub>4</sub>) trend in the inner region while the external isophotes are substantially boxy (negative C<sub>4</sub>), possibly due to merging processes ( e.g. Bender et al. 1988). This is the only clear evidence of significantly boxy isophotes ever found in a BL Lac host. Another example of a very luminous host galaxy (M<sub>R</sub>= –24.8, or -24.45 if de Vaucouleurs law is fitted) was reported by Heidt et al. (1999) for 1ES 1741+196. Also in this case the host galaxy isophotes have high ellipticity ( $``$ 0.35). There is no information, however, about the detailed shape of the isophotes and the amount of possible boxiness. MS 1312.1-422 This source, drawn from the EMSS of BL Lacs (Maccacaro et al. . 1994) was observed during bad seeing conditions (seeing of 1.4<sup>′′</sup> ) but since it is at relatively low redshift ( z = 0.108; Morris et al. 1991) it is rather well resolved. The host galaxy is indeed dominant with respect to the nuclear source ( ratio nucleus/host = 0.1). Our fit of the brightness profile yields M<sub>R</sub> = -23.4. No other detection of this host galaxy can be found in the literature. Note that in the calculation of $`\alpha _{OX}`$ it is usually the luminosity of the whole object (nucleus + host) which is used in the calculation. Such a procedure if applied to 1312-42 would overestimate the optical flux by a factor $``$ 10. MS 1332.6-293 The target belongs to the EMSS sample of BL Lacs although its classification is uncertain. Optical spectra showed either emission lines at z=0.256 or strong CaII break (Stocke et al. 1991 ) due to a substantial contribution from stellar emission. Our image shows this object is only marginally resolved. This is in part due to the bad seeing ($``$ 1.5<sup>′′</sup> ) and also because the host galaxy is substantially under-luminous (M<sub>R</sub> = -21.3) with respect to the average of the host galaxies of BL Lacs (M<sub>R</sub> = -23.8 Falomo & Kotilainen 1999). We note that the same object (1ES1322-297) is listed in the Einstein Slew sample of BL Lacs (Perlman et al. 1996) and has a redshift z = 0.512 quite different from the previous finding. Since in neither cases there are spectra published we are not able to make our own judgment of the validity of the redshift values. However, the latter value seems confirmed by another optical spectrum (albeit noisy) reported by Rector (1998). At z = 0.512 the host galaxy and point source would be much more luminous (M<sub>nuc</sub> = -23.5 and M<sub>host</sub> = –23.3) and well within the averages of these types of objects. ## 4 Summary and conclusions We have presented optical images of a number of BL Lacs that were not previously well studied. For several of these objects the first detection of the host galaxy is presented here. The properties of the hosts are consistent with them being luminous ellipticals as found in previous similar studies. For two of the resolved objects that have not a spectroscopic redshift we derive a photometric redshift based on the observed properties of the surrounding nebulosity. A bright and boxy elliptical We find that the external isophotes of the luminous host galaxy of 1101-23 are significantly boxy while the inner disky region suggests the presence of a small disc component. This is the first clear example of a boxy galaxy hosting a BL Lac object. Boxy isophotes are observed in a fraction of luminous ellipticals (Bender 1988) and could be ascribed to merging events from equal mass galaxies (e.g. Naab et al. 1999). It would be interesting to know what fraction of hosts of BL Lacs exhibit boxy isophotes as compared with non active ellipticals. Very little data are, however, available on isophote shapes of BL Lacs host galaxies because the presence of the bright nucleus and the quality of data often hinder a reliable estimation of this parameter especially at high redshift. For relatively low redshift objects with high resolution images it should be possible to investigate the isophote shape in a systematic way. The immediate environment of BL Lacs Our spectroscopy has allowed us to derive a redshift for 0301-24 (z = 0.26) and possibly for 0754+10 (z=0.28). Both objects have companion galaxies at redshifts very similar to that of the BL Lacs. The companions and the BL Lacs are thus very likely to be gravitationally bound. A third case is PKS 0829+04 for which we took the spectra of two galaxies in the immediate environment and found that one is at the same redshift as the BL Lac object. These spectroscopic results improve the scanty data on redshifts of companion galaxies of BL Lacs. Together with previous findings ( Falomo et al. 1993a,b; Pesce et al 1994,1995; Heidt et al. 1999) our new results yield convincing evidence that galaxies around BL Lacs are (often) gravitationally bound with the BL Lacs. On the other hand only in very few cases do these interactions lead to significantly (observable) disturbed morphology (see e.g. Falomo et al. 1995, Heidt et al. 1999). ###### Acknowledgements. This work was partly supported by the Italian Ministry for University and Research (MURST) under grant Cofin98-02-32. This research has made use of the NASA/IPAC Extragalactic Database (NED) which is operated by the Jet Propulsion Laboratory, California Institute of Technology, under contract with the National Aeronautics and Space Administration. We thank A. Wolter for providing coordinates and finding charts of REX 0353-36 before publication. RF thanks the ESO visitor program for hospitality during several visits at the ESO headquarter. ## 5 References Abraham R.G., McHardy I.M., Crawford C.S. 1991, MNRAS 252, 482 Bender R. 1988, A&A 193, L7 Bender R., Dobereiner S., Mollenhoff C. 1988, A&AS 74,385. Carswell R.F., Strittmatter P.A., Williams R.E., Kinman T.D., Serkowski K. 1974, ApJ 190, L101 Coleman G.D., Wu C.C., Weedman D.W. 1980, ApJS 43, 393 Falomo R. 1996, MNRAS 283, 241 Falomo R., Kotilainen J. 1999 A&A 321, 374 Falomo R., Pesce J. 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warning/0004/math0004030.html	ar5iv	text	# Characterization of Cyclic and Separating Vectors and Application to an Inverse Problem in Modular Theory II. Semifinite Factors ## 1 The Inverse Problem in Modular Theory Let $`_0`$ be a von Neumann algebra on a separable Hilbert space $`_0`$ with a cyclic and separating vector $`u_0`$. Then modular theory shows the existence of a modular operator $`\mathrm{\Delta }_0`$ and a modular conjugation $`\mathrm{J}_0`$ (the modular objects $`(\mathrm{\Delta }_0,\mathrm{J}_0)`$) belonging to the vector $`u_0`$. In this paper we examine the inverse problem of constructing algebras $``$ having the same cyclic and separating vector and modular objects as $`_0`$: The Inverse Problem Let $`(\mathrm{\Delta }_0,\mathrm{J}_0)`$ be the modular objects for the von Neumann algebra $`_0`$ with cyclic and separating vector $`u_0`$. Characterize all von Neumann algebras $``$ isomorphic to $`_0`$ with the following properties: 1. $`u_0`$ is also cyclic and separating for $``$, 2. $`(\mathrm{\Delta }_0,\mathrm{J}_0)`$ are the modular objects for $`(,u_0)`$. Let $`NF__0(\mathrm{\Delta }_0,\mathrm{J}_0,u_0)`$ denote all solutions $``$ of the inverse problem. In \[Bol\] the following theorems were shown: ###### Theorem 1.1. Let ($`_0,_0`$) be a finite von Neumann factor. Let further $`u_0`$. Then there is exactly one operator $`\mathrm{T}_u\eta _0`$ associated with the vector $`u`$, s.t. $`u=\mathrm{T}_uu_{\mathrm{tr}}`$ where $`u_{\mathrm{tr}}_0`$ is a cyclic trace vector. This operator has the following properties: 1. $`\mathrm{tr}(\mathrm{T}_u\mathrm{T}_u^{})=\mathrm{tr}(\mathrm{T}_u^{}\mathrm{T}_u)<\mathrm{}`$. 2. $`u`$ is cyclic, iff $`\mathrm{T}_u`$ is injective. 3. $`u`$ is separating, iff $`\mathrm{T}_u`$ has dense range. 4. $`u`$ is cyclic and separating iff $`\mathrm{T}_u`$ is injective and has dense range, i.e. iff $`\mathrm{T}_u`$ is invertible. ###### Theorem 1.2. Let $`\mathrm{T}\eta _0`$. Then there is a vector $`u_0`$ s.t $`\mathrm{T}=\mathrm{T}_u`$ in the sense of Theorem 1.1 iff $`\mathrm{tr}\mathrm{TT}^{}=\mathrm{tr}\mathrm{T}^{}\mathrm{T}<\mathrm{}`$. In the second section of this paper these theorems were generalized to infinite semifinite factors. For this purpose we first consider a special case of such factors, a matrix algebra of finite factors, where the trace vector is replaced by a sequence of vectors, constructed from the trace vectors of the constituting factors. With the help of this result we show the analogue of the following result, also obtained in \[Bol\] for finite factors: ###### Theorem 1.3. Let $`_0`$ be a finite von Neumann factor with cyclic and separating vector $`u_0_0`$ and cyclic trace vector $`u_{\mathrm{tr}}_0`$. Let further $`\mathrm{T}_{u_0}\eta _0`$ be the invertible operator corresponding to $`u_0`$ and $`\mathrm{T}_{u_0}=\mathrm{HV}=(\mathrm{T}_{u_0}\mathrm{T}_{u_0}^{})^{1/2}\mathrm{V}`$ the polar decomposition of $`\mathrm{T}_{u_0}`$. Then we can calculate the modular objects $`(\mathrm{\Delta }_0,\mathrm{J}_0)`$ of $`(_0,u_0)`$ as follows: $$\mathrm{J}_0=\mathrm{JV}^{}\mathrm{JVJ}=\mathrm{VJV}^{},$$ where $`\mathrm{J}`$ is the conjugation corresponding to $`u_{\mathrm{tr}}`$, and $$\mathrm{\Delta }_0=\mathrm{J}_0\mathrm{H}_0^1\mathrm{J}_0\mathrm{H}_0,$$ where $`\mathrm{H}_0=\mathrm{H}^2=\mathrm{T}_{u_0}\mathrm{T}_{u_0}^{}`$. Then we will be in exactly the same situation as in the finite case, and can examine the inverse problem as in \[Bol\]. In contrast to that case the second simple class of solutions will never exists in this case (s. §4), but the classification of the solutions in the pure point spectrum case will be the same. Notice that in this paper all Hilbert spaces are separable, i.e. the von Neumann algebras are countably decomposable. ## 2 Characterization of Vectors by Affiliated Operators In this section we consider infinite but still semifinite factors, i.e. $`_0`$ is of type $`I_{\mathrm{}}`$ or $`II_{\mathrm{}}`$. In this case we have no trace vectors left. But nevertheless we can make a similar construction as in the finite case by considering the infinite factor as an infinite matrix of finite factors and using the results presented in \[Bol\]. As a model for such a matrix of finite factors we examine now the semifinite factor $`=𝒯L(_{\mathrm{}})`$ on $`𝒦=_{\mathrm{}}_{\mathrm{}}`$, where $`(𝒯,)`$ is a finite factor possessing a cyclic and separating vector and $`_{\mathrm{}}`$ is a infinite dimensional separable Hilbert space which we can identify with $`l_2()`$. Now $``$ is an (infinite) type $`I`$ ($`II`$) factor, if $`𝒯`$ is type $`I`$ ($`II`$). Further, since $`(𝒯,)`$ is finite and possesses a cyclic and separating vector, it possesses a cyclic trace vector $`u_{\mathrm{tr}}`$ (cf. \[KR86, Th. 8.2.8, Lem. 7.2.8\]). In the following we consider the elements of $`𝒦`$ as infinite dimensional matrices $`u=(u_i^k)_{i,k}`$ with entries $`u_i^k`$ s.t. $`_{i,k}u_i^k^2<\mathrm{}`$, where the lower index corresponds to the second component of the tensor product and the upper to the third, resp. Then we can write the elements of $``$ as matrices $`\mathrm{T}=(\mathrm{T}_{li})_{l,i}`$ with entries $`\mathrm{T}_{li}𝒯`$, where $$\mathrm{T}u=(\underset{i}{}\mathrm{T}_{li}u_i^k)_l^k.$$ Then the commutant $`^{^{}}`$ of $``$ is $`𝒯^{^{}}L(_{\mathrm{}})`$, where we can write an element in $`^{^{}}`$ as $`\mathrm{T}^{^{}}=(\mathrm{T}_{}^{^{}}{}_{}{}^{lk})_{l,k}`$ with entries $`\mathrm{T}_{}^{^{}}{}_{}{}^{lk}𝒯^{^{}}`$, where $$\mathrm{T}^{^{}}u=(\underset{k}{}\mathrm{T}_{}^{^{}}{}_{}{}^{lk}u_i^k)_i^l.$$ For the proofs in this section the following subalgebras of $``$ and $`^{^{}}`$ are important: $$\begin{array}{cc}\hfill _0& :=\{\mathrm{M}=(\mathrm{M}_{ij})_{ij}\|\mathrm{M}_{ij}0\text{ for only finitely many }i,j\}\hfill \\ \hfill _0^{^{}}& :=\{\mathrm{M}^{^{}}=(\mathrm{M}_{}^{^{}}{}_{}{}^{ij})^{ij}^{^{}}\|\mathrm{M}_{}^{^{}}{}_{}{}^{ij}0\text{ for only finitely many }i,j\}.\hfill \end{array}$$ Now we define the following sequence of vectors in $`𝒦`$, which is the analogue to the trace vector: $$v_k:=(\delta _i^j\delta _{ik}u_{\mathrm{tr}})_i^j.$$ Further we define $$D_0:=\mathrm{lin}\{\mathrm{M}v_k\|\mathrm{M},k\}𝒦$$ and $$D_0^{^{}}:=\mathrm{lin}\{\mathrm{M}^{^{}}v_k\|\mathrm{M}^{^{}}^{^{}},k\}𝒦.$$ Now the trace $`\mathrm{tr}`$ of $``$, which is a n.s.f. tracial weight, is $`\mathrm{tr}=\mathrm{tr}_𝒯\mathrm{tr}_{L(_{\mathrm{}})}`$, where $`\mathrm{tr}_𝒯`$ is the trace on $`𝒯`$ and $`\mathrm{tr}_{L(_{\mathrm{}})}`$ the standard trace on $`_{\mathrm{}}`$. It can be written with the help of the vectors $`(v_k)`$: $$\begin{array}{cc}\hfill \mathrm{tr}(\mathrm{M})& =\underset{k}{}\mathrm{tr}_𝒯(\mathrm{M}_{kk})=\underset{k}{}\mathrm{M}_{kk}u_{\mathrm{tr}}\|u_{\mathrm{tr}}\hfill \\ & =\underset{k}{}\mathrm{M}v_k\|v_k\hfill \\ & =\underset{k}{}\lambda d\mathrm{E}_\lambda ^\mathrm{M}v_k^2\mathrm{M}=(\mathrm{M}_{ij})^+,\hfill \end{array}$$ where $`\mathrm{M}=\lambda 𝑑\mathrm{E}_\lambda ^\mathrm{M}`$ is the spetral measure of $`\mathrm{M}`$. As in \[Bol\] we can continue the trace to all the positive closed operators $`\mathrm{A}`$ affiliated with $``$ by $$\mathrm{tr}(\mathrm{A}):=\underset{k}{}\lambda d\mathrm{E}_\lambda ^\mathrm{A}v_k^2,$$ (2.1) where $`\mathrm{E}_\lambda ^\mathrm{A}`$ is the spectral measure of $`\mathrm{A}`$. Now we can associate an operator $`\mathrm{T}_{ij}\eta 𝒯`$ with every component $`u_i^j`$ of a vector $`u=(u_i^j)_i^j`$, s.t. $`u_{\mathrm{tr}}𝒟(\mathrm{T}_{\mathrm{ij}})`$ and $`\mathrm{T}_{ij}u_{\mathrm{tr}}=u_i^j`$ (cf. \[Bol\]). These operators give rise to a linear operator $`\stackrel{~}{\mathrm{T}}_u`$ defined by $$\begin{array}{cc}\hfill \stackrel{~}{\mathrm{T}}_u:𝒟(\stackrel{~}{\mathrm{T}}_\mathrm{u}):=D_0^{^{}}𝒦& 𝒦\hfill \\ \hfill \mathrm{M}^{^{}}v_k& \stackrel{~}{\mathrm{T}}_u\mathrm{M}^{^{}}v_k:=(\mathrm{T}_{ik}\mathrm{M}_{}^{^{}}{}_{}{}^{jk}u_{\mathrm{tr}})_i^j.\hfill \end{array}$$ (2.2) Now we can prove ###### Lemma 2.1. Let $`\stackrel{~}{\mathrm{T}}_u`$ defined by (2.2). Then $`\stackrel{~}{\mathrm{T}}_u`$ is densely defined and closable. Let $`\mathrm{T}_u`$ be its closure. Then $`D_0^{^{}}𝒟(\mathrm{T}_\mathrm{u}^{})`$, $`\mathrm{T}_u`$ is affiliated with $``$, and $$\underset{k}{}\mathrm{T}_uv_k=u,$$ where the convergence is absolute. ###### Proof. 1. First we must show that $`\stackrel{~}{\mathrm{T}}_u`$ is well defined. Observe first that $$\begin{array}{cc}\hfill \underset{i,j}{}\mathrm{T}_{ik}\delta ^{jk}u_{\mathrm{tr}}^2& =\underset{i}{}\mathrm{T}_{ik}u_{\mathrm{tr}}^2\hfill \\ & =\underset{i}{}u_i^k^2<\mathrm{},\hfill \end{array}$$ i.e. $`\stackrel{~}{\mathrm{T}}_uv_k𝒦`$ for every $`k`$. Let now $`\mathrm{M}^{^{}}^{^{}}`$ be arbitrary. Then $`\mathrm{M}^{^{}}\stackrel{~}{\mathrm{T}}_uv_k𝒦`$ and $$\begin{array}{cc}\hfill \mathrm{}& >\mathrm{M}^{^{}}\stackrel{~}{\mathrm{T}}_uv_k^2\hfill \\ & =\underset{i,j}{}\mathrm{M}_{}^{^{}}{}_{}{}^{jk}\mathrm{T}_{ik}u_{\mathrm{tr}}^2\hfill \\ & =\underset{i,j}{}\mathrm{T}_{ik}\mathrm{M}_{}^{^{}}{}_{}{}^{jk}u_{\mathrm{tr}}^2\text{ (cf. }\text{[Bol, Prop.2.1.]}\text{)}\hfill \\ & =\stackrel{~}{\mathrm{T}}_u\mathrm{M}^{^{}}v_k^2\hfill \end{array}$$ (2.3) for every $`k`$, hence $`\stackrel{~}{\mathrm{T}}_u`$ is well defined. 2. Now we show that $`𝒟(\stackrel{~}{\mathrm{T}}_\mathrm{u})`$ is dense in $`𝒦`$. First the elements with only finitely many entries not $`0`$ are dense in $`𝒦`$. Further every such element is a linear combination of elements of the type $`(u_i^j\delta _{ik})_i^j`$, again all but a finite number equal $`0`$. Since $`u_{\mathrm{tr}}`$ is cyclic for $`𝒯^{^{}}`$, we can approximate these elements by elements of the form $`(\mathrm{M}_{}^{^{}}{}_{}{}^{jk}\delta _{ik}u_{\mathrm{tr}})_i^j=:\mathrm{M}^{^{}}v_k`$ with $`\mathrm{M}^{^{}}=(\mathrm{M}_{}^{^{}}{}_{}{}^{jk}\delta ^{ik})^{ji}_0^{^{}}^{^{}}`$, hence $`𝒟(\stackrel{~}{\mathrm{T}}_\mathrm{u})=D_0^{^{}}`$ is dense in $`𝒦`$. 3. In this step we want to show that $`\stackrel{~}{\mathrm{T}}_u`$ is closable. Let $`x=\mathrm{M}^{^{}}v_k𝒟(\stackrel{~}{\mathrm{T}}_\mathrm{u})`$, $`y=\mathrm{N}^{^{}}v_j𝒟(\mathrm{S}):=D_0^{^{}}`$ ($`k,j`$), where $`\mathrm{S}:=(\mathrm{T}_{li}^{})_{i,l}`$ is defined analogously to $`\stackrel{~}{\mathrm{T}}_u`$, hence it is a densely defined operator, too (All $`\mathrm{T}_{li}^{}`$ are closed operators affiliated with $`𝒯`$ and $`u_{\mathrm{tr}}𝒟(\mathrm{T}_{\mathrm{li}}^{})`$, cf. \[Bol, Prop. 2.1.\], and $`\mathrm{T}_{li}^{}u_{\mathrm{tr}}^2=\mathrm{T}_{li}u_{\mathrm{tr}}^2`$). Now $$\begin{array}{cc}\hfill \stackrel{~}{\mathrm{T}}_ux\|y& =\stackrel{~}{\mathrm{T}}_u\mathrm{M}^{^{}}v_k\|\mathrm{N}^{^{}}v_j\hfill \\ & =\underset{i,l}{}\mathrm{T}_{ik}\mathrm{M}_{}^{^{}}{}_{}{}^{lk}u_{\mathrm{tr}}\|\mathrm{N}_{}^{^{}}{}_{}{}^{lj}\delta _{ij}u_{\mathrm{tr}}\hfill \\ & =\underset{l}{}\mathrm{M}_{}^{^{}}{}_{}{}^{lk}u_{\mathrm{tr}}\|\mathrm{T}_{jk}^{}\mathrm{N}_{}^{^{}}{}_{}{}^{lj}u_{\mathrm{tr}}\hfill \\ & =\underset{i,l}{}\delta _{ik}\mathrm{M}_{}^{^{}}{}_{}{}^{lk}u_{\mathrm{tr}}\|\mathrm{T}_{ji}^{}\mathrm{N}_{}^{^{}}{}_{}{}^{lj}u_{\mathrm{tr}}\hfill \\ & =x\|\mathrm{S}y.\hfill \end{array}$$ This shows $`y𝒟(\stackrel{~}{\mathrm{T}}_{\mathrm{u}}^{}{}_{}{}^{})`$, $`\stackrel{~}{\mathrm{T}}_u^{}y=\mathrm{S}y`$, and $`\mathrm{S}(\stackrel{~}{\mathrm{T}}_u)^{}`$, hence $`\stackrel{~}{\mathrm{T}}_u`$ is closable. This shows also, that $`D_0^{^{}}𝒟((\stackrel{~}{\mathrm{T}}_\mathrm{u})^{})=𝒟(\mathrm{T}_\mathrm{u}^{})`$. 4. To show that $`\mathrm{T}_u`$ is affiliated with $``$, let $`\mathrm{U}^{^{}}=(\mathrm{U}_{}^{^{}}{}_{}{}^{ij})_{i,j}^{^{}}`$ be a unitary. Then $`\mathrm{U}^{^{}}D_0^{^{}}=D_0^{^{}}`$. Let now $`x=\mathrm{M}^{^{}}v_kD_0^{^{}}=𝒟(\stackrel{~}{\mathrm{T}}_\mathrm{u})`$. Then $$\begin{array}{cc}\hfill \mathrm{U}^{^{}}\stackrel{~}{\mathrm{T}}_ux& =\mathrm{U}^{^{}}(\mathrm{T}_{ik}\mathrm{M}_{}^{^{}}{}_{}{}^{jk}u_{\mathrm{tr}})_i^j\hfill \\ & =(\underset{j}{}\mathrm{U}_{}^{^{}}{}_{}{}^{lj}\mathrm{T}_{ik}\mathrm{M}_{}^{^{}}{}_{}{}^{jk}u_{\mathrm{tr}})_i^l\hfill \\ & =(\underset{j}{}\mathrm{T}_{ik}\mathrm{U}_{}^{^{}}{}_{}{}^{lj}\mathrm{M}_{}^{^{}}{}_{}{}^{jk}u_{\mathrm{tr}})_i^l\text{ (cf. }\text{[Bol, Prop.2.1.]}\text{)}\hfill \\ & =(\mathrm{T}_{ik}\underset{j}{}\mathrm{U}_{}^{^{}}{}_{}{}^{lj}\mathrm{M}_{}^{^{}}{}_{}{}^{jk}u_{\mathrm{tr}})_i^l\hfill \\ & =\stackrel{~}{\mathrm{T}}_u\mathrm{U}^{^{}}x\hfill \end{array}$$ This shows $`\mathrm{U}^{^{}}\stackrel{~}{\mathrm{T}}_u=\stackrel{~}{\mathrm{T}}_u\mathrm{U}^{^{}}`$ for every unitary $`\mathrm{U}^{^{}}^{^{}}`$, hence, since $`D_0^{^{}}`$ is a core for $`\mathrm{T}_u`$, $$\mathrm{U}^{^{}}\mathrm{T}_u=\mathrm{T}_u\mathrm{U}^{^{}}\mathrm{U}^{^{}}𝒰(^{^{}}).$$ 5. In the last step we calculate $$\begin{array}{cc}\hfill \underset{k}{}\mathrm{T}_uv_k& =\underset{k}{}(\mathrm{T}_{ik}\delta ^{jk}u_{tr})_i^j\hfill \\ & =(\mathrm{T}_{ij}u_{tr})_i^j=u.\hfill \end{array}$$ Now we can give the following definition: ###### Definition 2.1. For every vector $`u=(u_i^j)𝒦`$ we denote by $`\mathrm{T}_u=(\mathrm{T}_{ij})`$ an operator affiliated with $``$ s.t. $`u_{\mathrm{tr}}𝒟(\mathrm{T}_{\mathrm{ij}})`$ for all $`i,j`$, $`\mathrm{T}_{ij}u_{\mathrm{tr}}=u_i^j`$, and $`_k\mathrm{T}_uv_k=u`$, which exists according to Lemma 2.1. The next proposition shows some usefull properties of the operators occuring in Definition 2.1 ###### Proposition 2.2. Let $`\mathrm{T}\eta `$, $`v_k𝒟(\mathrm{T})`$ ($`k`$), $`_k\mathrm{T}v_k^2<\mathrm{}`$. Then 1. $`v_k𝒟(\mathrm{T})`$, $`v_k𝒟(\mathrm{T}^{})`$, and $`v_k𝒟((\mathrm{T}^{}\mathrm{T})^{1/2})`$ for all $`k`$. 2. $`D_0`$ is a core for $`\mathrm{T}`$, $`\mathrm{T}^{}`$, and $`(\mathrm{T}^{}\mathrm{T})^{1/2}`$. 3. $`^{^{}}v_k𝒟(\mathrm{T})`$, $`^{^{}}v_k𝒟(\mathrm{T}^{})`$, and $`^{^{}}v_k𝒟((\mathrm{T}^{}\mathrm{T})^{1/2})`$ for all $`k`$. 4. $`D_0^{^{}}`$ is a core for $`\mathrm{T}`$, $`\mathrm{T}^{}`$, and $`(\mathrm{T}^{}\mathrm{T})^{1/2}`$. ###### Proof. 1. Let $`\mathrm{T}=\mathrm{VH}`$ the polar decomposition of $`\mathrm{T}`$, and $`\mathrm{E}_\lambda `$ the spectral resolution of $`\mathrm{H}`$. Then $$\begin{array}{cc}\hfill \lambda ^2d\mathrm{E}_\lambda \mathrm{U}v_k^2& \underset{l}{}\lambda ^2d\mathrm{E}_\lambda \mathrm{U}v_l^2\hfill \\ & =\underset{l}{}\lambda ^2d\mathrm{U}^{}\mathrm{E}_\lambda v_l^2\text{(s. (}\text{2.1}\text{))}\hfill \\ & =\underset{l}{}\lambda ^2d\mathrm{E}_\lambda v_l^2\hfill \\ & =\underset{l}{}\mathrm{H}v_l^2=\underset{l}{}\mathrm{T}v_l^2<\mathrm{}\hfill \end{array}$$ for every unitary $`\mathrm{U}`$ and every $`k`$, i.e. $`v_k𝒟(\mathrm{H})=𝒟(\mathrm{T})`$ for every $`k`$. Now $`\mathrm{T}^{}=\mathrm{HV}^{}`$, and, since $`\mathrm{V}^{}`$, also $`v_k𝒟(\mathrm{T}^{})`$. 2. 1) shows that $`D_0𝒟(\mathrm{T})`$, further $`D_0`$ is dense in $`𝒦`$. Now $`D_0`$ is invariant under the unitary group $`e^{it\mathrm{H}}`$, i.e. $`D_0`$ is a core for $`\mathrm{H}`$ and also for $`\mathrm{T}`$. The assertion for $`\mathrm{T}^{}`$ follows analogous. 3. This follows from 1) and \[Bol, Prop. 2.1.\]. 4. Now for every $`\mathrm{M}=(\mathrm{M}_{ij})_{ij}`$ there exists exactly one $`\mathrm{M}^{^{}}=(\mathrm{M}_{}^{^{}}{}_{}{}^{ij})^{^{}}`$ s.t. $`\mathrm{M}_{}^{^{}}{}_{}{}^{ij}u_{\mathrm{tr}}=\mathrm{M}_{ji}u_{\mathrm{tr}}`$ ($`\mathrm{M}_{}^{^{}}{}_{}{}^{ij}:=\mathrm{JM}_{ji}\mathrm{J}`$, where $`\mathrm{J}`$ is the conjugation w.r.t. $`u_{\mathrm{tr}}`$). Now define $$\mathrm{M}_{(k,l)}^{^{}}:=\mathrm{E}_k^{^{}}\mathrm{M}^{^{}}\mathrm{E}_l^{^{}},$$ where $`\mathrm{E}_{(k)}^{^{}}:=(\delta _{ik}\delta ^{ij})^{ij}`$. Then $$\mathrm{M}v_k=\underset{l}{}\mathrm{M}_{(k,l)}^{^{}}v_l$$ (2.4) and $$\begin{array}{cc}\hfill \underset{l}{}\mathrm{TM}_{(k,l)}^{^{}}v_l^2& =\underset{l}{}\mathrm{M}_{(k,l)}^{^{}}\mathrm{T}v_l^2\hfill \\ & \underset{l}{}\mathrm{M}^{^{}}^2\mathrm{T}v_l^2\hfill \\ & \mathrm{M}^2\underset{l}{}\mathrm{T}v_l^2<\mathrm{}.\hfill \end{array}$$ Hence $`_l\mathrm{TM}_{(k,l)}^{^{}}v_l`$ converges and therefore $`_l\mathrm{M}_{(k,l)}^{^{}}v_l`$ converges in the graph norm of $`\mathrm{T}`$ to $`\mathrm{M}v_k`$, i.e. also $`D_0^{^{}}`$ is a core, since $`D_0`$ is it. ###### Lemma 2.3. Let $`\mathrm{T}\eta `$ be as in Proposition 2.2. Then there are $`\mathrm{T}_{ij}\eta 𝒯`$ with $`u_{\mathrm{tr}}𝒟(\mathrm{T}_{\mathrm{ij}})`$ s.t. $`\mathrm{T}=(\mathrm{T}_{ij})_{i,j}`$ and $`\mathrm{T}=\mathrm{T}_u`$ with $`u:=_k\mathrm{T}v_k`$ in the sense of Definition 2.1. ###### Proof. Set $`\mathrm{E}_k:=[^{^{}}v_k]`$ ($`k`$). Then matrix calculation shows that $`(\mathrm{E}_k)`$ is a family of orthogonal, equivalent, finite projections, s.t. $$\mathrm{E}_k\mathrm{E}_k=(\delta _{ki}\delta _{ij}𝒯)_{i,j}$$ and $`_k\mathrm{E}_k=\mathrm{Id}`$. Now $`\mathrm{E}_{ij}:\mathrm{M}^{^{}}v_j\mathrm{M}^{^{}}v_i`$ defines a selfadjoint system of matrix units $`(\mathrm{E}_{ij})`$ s.t. $`\mathrm{E}_{kk}=\mathrm{E}_k`$ for every $`k`$. Now define operators $$\begin{array}{cc}\hfill \mathrm{S}_{ij}:𝒟(\mathrm{S}_{\mathrm{ij}}):=D_0^{^{}}𝒦& 𝒦\hfill \\ \hfill \underset{k}{}\mathrm{M}_k^{^{}}v_k& \underset{k}{}\mathrm{E}_{ki}\mathrm{TE}_{jk}\mathrm{M}_k^{^{}}v_k\hfill \end{array}$$ and $$\begin{array}{cc}\hfill \stackrel{~}{\mathrm{S}}_{ji}:𝒟(\stackrel{~}{\mathrm{S}}_{\mathrm{ij}}):=D_0^{^{}}𝒦& 𝒦\hfill \\ \hfill \underset{k}{}\mathrm{M}_k^{^{}}v_k& \underset{k}{}\mathrm{E}_{kj}\mathrm{T}^{}\mathrm{E}_{ik}\mathrm{M}_k^{^{}}v_k.\hfill \end{array}$$ Since $`D_0^{^{}}`$ is dense in $`𝒦`$ (cf. proof of Lemma 2.1) and a core both for $`\mathrm{T}`$ and for $`\mathrm{T}^{}`$ they are well defined and densely defined. Let now $`x𝒟(\mathrm{S}_{\mathrm{ij}})`$ and $`y𝒟(\stackrel{~}{\mathrm{S}}_{\mathrm{ji}})`$. Then $$\begin{array}{cc}\hfill \mathrm{S}_{ij}x\|y& =\underset{k}{}\mathrm{E}_{ki}\mathrm{TE}_{jk}x\|y\hfill \\ & =\underset{k}{}x\|\mathrm{E}_{kj}\mathrm{T}^{}\mathrm{E}_{ik}y\hfill \\ & =x\|\stackrel{~}{\mathrm{S}}_{ji}y.\hfill \end{array}$$ This means that $`y𝒟(\mathrm{S}_{\mathrm{ij}}^{})`$ and $`\mathrm{S}_{ij}^{}y=\stackrel{~}{\mathrm{S}}_{ji}y`$, i.e. $`\stackrel{~}{\mathrm{S}}_{ji}\mathrm{S}_{ij}^{}`$, hence $`\mathrm{S}_{ij}`$ is closable since $`\stackrel{~}{\mathrm{S}}_{ji}`$ is densely defined. Let now $`\stackrel{~}{\mathrm{T}}_{ij}`$ be the closure of $`\mathrm{S}_{ij}`$. Then $`D_0^{^{}}=𝒟(\mathrm{S}_{\mathrm{ij}})`$ is a core for $`\stackrel{~}{\mathrm{T}}_{ij}`$. Since $`\mathrm{U}^{^{}}D_0^{^{}}=D_0^{^{}}`$ and $$\begin{array}{cc}\hfill \mathrm{U}^{^{}}\mathrm{S}_{ij}(\underset{k}{}\mathrm{M}_k^{^{}}v_k)& =\underset{k}{}\mathrm{U}^{^{}}\mathrm{E}_{ki}\mathrm{TE}_{jk}\mathrm{M}_k^{^{}}v_k\hfill \\ & =\underset{k}{}\mathrm{E}_{ki}\mathrm{TE}_{jk}\mathrm{U}^{^{}}\mathrm{M}_k^{^{}}v_k\hfill \\ & =\mathrm{S}_{ij}\mathrm{U}^{^{}}(\underset{k}{}\mathrm{M}_k^{^{}}v_k)\hfill \end{array}$$ for every unitary $`\mathrm{U}^{^{}}^{^{}}`$ and every element $`(\mathrm{M}_k^{^{}})v_kD_0^{^{}}`$, it follows that $`\mathrm{U}^{^{}}\stackrel{~}{\mathrm{T}}_{ij}=\stackrel{~}{\mathrm{T}}_{ij}\mathrm{U}^{^{}}`$ and $`\stackrel{~}{\mathrm{T}}_{ij}`$ is affiliated with $``$. Further $$\begin{array}{cc}\hfill \mathrm{E}_{mn}\mathrm{S}_{ij}(\underset{k}{}\mathrm{M}_k^{^{}}v_k)& =\underset{k}{}\mathrm{E}_{mn}\mathrm{E}_{ki}\mathrm{TE}_{jk}\mathrm{M}_k^{^{}}v_k\hfill \\ & =\mathrm{E}_{mi}\mathrm{TE}_{jm}\mathrm{E}_{mn}\mathrm{M}_n^{^{}}v_n\hfill \\ & =\underset{k}{}\mathrm{E}_{ki}\mathrm{TE}_{jk}\mathrm{E}_{mn}\mathrm{M}_n^{^{}}v_n\hfill \\ & =\mathrm{S}_{ij}\mathrm{E}_{mn}(\underset{k}{}\mathrm{M}_k^{^{}}v_k),\hfill \end{array}$$ hence $`\stackrel{~}{\mathrm{T}}_{ij}`$ is affiliated with $`𝒯=\{\mathrm{E}_{mn}\|m,n\}^{^{}}`$. Now set $`\mathrm{T}_{ij}:=\mathrm{V}^{}\stackrel{~}{\mathrm{T}}_{ij}\mathrm{V}`$, where $$\begin{array}{cc}\hfill \mathrm{V}:& 𝒦\hfill \\ \hfill v& (\delta _{1i}\delta _i^jv)_i^j\hfill \end{array}$$ is the canonical partial isometry from $``$ to $`𝒦=_{\mathrm{}}_{\mathrm{}}`$. Now $`u_{\mathrm{tr}}𝒟(\mathrm{T}_{\mathrm{ij}})`$ since $`\mathrm{V}u_{\mathrm{tr}}=v_1𝒟(\mathrm{T}_{\mathrm{ij}})`$, and with $`u_i^j:=\mathrm{T}_{ij}u_{\mathrm{tr}}`$ $$\begin{array}{cc}\hfill \underset{i,j}{}u_i^j^2& =\underset{i,j}{}\mathrm{T}_{ij}u_{\mathrm{tr}}^2\hfill \\ & =\underset{i,j}{}\mathrm{V}^{}\stackrel{~}{\mathrm{T}}_{ij}\mathrm{V}u_{\mathrm{tr}}^2\hfill \\ & =\underset{i,j}{}\mathrm{E}_{1i}\mathrm{TE}_{j1}v_1^2\hfill \\ & =\underset{i,j}{}\mathrm{E}_i\mathrm{T}v_j^2\hfill \\ & =\underset{j}{}\mathrm{T}v_j^2<\mathrm{}\hfill \end{array}$$ s.t. $`u:=_k\mathrm{T}v_k=(u_i^j)_i^j=(\mathrm{T}_{ij}u_{\mathrm{tr}})_i^j𝒦`$. This means that we can construct the operator $`\mathrm{T}_u=(\mathrm{T}_{ij})_{ij}`$ according to Lemma 2.1. Now $`\mathrm{T}_u`$ and $`\mathrm{T}`$ coincide on the core $`D_0^{^{}}`$, and hence they are equal. ∎ ###### Corollary 2.4. The operator $`\mathrm{T}_u`$ defined in Definition 2.1 is unique. ###### Corollary 2.5. Let $`\mathrm{T}_u`$ be the operator defined in Definition 2.1. Then $`v_k𝒟(\mathrm{T}_\mathrm{u})`$ for every $`k`$ and $$\mathrm{T}_u\mathrm{M}v_k=(\underset{l}{}\mathrm{T}_{il}\mathrm{M}_{lk}\delta _k^ju_{\mathrm{tr}})_i^j.$$ ###### Proof. Proposition 2.2 shows that $`v_k𝒟(\mathrm{T}_\mathrm{u})`$ for every $`k`$ and $`\mathrm{M}v_k=_l\mathrm{M}_{(k,l)}^{^{}}v_l`$ (cf. (2.4)). Now $$\begin{array}{cc}\hfill \mathrm{T}_u\mathrm{M}v_k& =\mathrm{T}_u\underset{l}{}\mathrm{M}_{(k,l)}^{^{}}v_l\hfill \\ & =\underset{l}{}\mathrm{T}_u\mathrm{M}_{(k,l)}^{^{}}v_l\hfill \\ & =\underset{l}{}(\mathrm{T}_{il}\mathrm{M}_{}^{^{}}{}_{}{}^{k,l}\delta _{jk}u_{\mathrm{tr}})_i^j\hfill \\ & =\underset{l}{}(\mathrm{T}_{il}\mathrm{M}_{lk}\delta _k^ju_{\mathrm{tr}})_i^j.\hfill \end{array}$$ Now we can formulate the following lemma: ###### Lemma 2.6. Let $`\mathrm{T}_u`$ be the operator defined in Definition 2.1. Then: 1. $`\mathrm{tr}(\mathrm{T}_u^{}\mathrm{T}_u)=\mathrm{tr}(\mathrm{T}_u\mathrm{T}_u^{})<\mathrm{}`$. 2. $`u`$ is cyclic, iff $`\mathrm{T}_u`$ is injective. 3. $`u`$ is separating, iff $`\mathrm{T}_u`$ has dense range. 4. $`u`$ is cyclic and separating iff $`\mathrm{T}_u`$ is injective and has dense range, i.e. iff $`\mathrm{T}_u`$ is invertible. For the proof we need: ###### Proposition 2.7. Let $`𝒯`$ be a (finite) von Neumann algebra with cyclic trace vector $`u_{\mathrm{tr}}`$. Let further $`\mathrm{S},\mathrm{T}\eta 𝒯`$ with $`u_{\mathrm{tr}}𝒟(\mathrm{S})𝒟(\mathrm{T})`$ and $`\mathrm{M},\mathrm{N}𝒯`$. Then $$\mathrm{MT}u_{\mathrm{tr}}\|\mathrm{NS}u_{\mathrm{tr}}=\mathrm{S}^{}\mathrm{N}^{}u_{\mathrm{tr}}\|\mathrm{T}^{}\mathrm{M}^{}u_{\mathrm{tr}}.$$ (2.5) ###### Proof. Let $`(\mathrm{E}_n)`$ and $`(\mathrm{F}_n)`$ be bounding sequences for $`\mathrm{T}`$ and $`\mathrm{S}`$, resp. (cf. \[KR83, Lem. 5.6.14\]). Then: $$\begin{array}{cc}\hfill \mathrm{MT}u_{\mathrm{tr}}\|\mathrm{NS}u_{\mathrm{tr}}& =\underset{n\mathrm{}}{lim}\mathrm{MTE}_nu_{\mathrm{tr}}\|\mathrm{NSF}_nu_{\mathrm{tr}}\hfill \\ & =\underset{n\mathrm{}}{lim}(\mathrm{SF}_n)^{}\mathrm{N}^{}u_{\mathrm{tr}}\|(\mathrm{TE}_n)^{}\mathrm{M}^{}u_{\mathrm{tr}}\hfill \\ & =\underset{n\mathrm{}}{lim}\mathrm{F}_n\mathrm{S}^{}\mathrm{N}^{}u_{\mathrm{tr}}\|\mathrm{E}_n\mathrm{T}^{}\mathrm{M}^{}u_{\mathrm{tr}}\hfill \\ & =\mathrm{S}^{}\mathrm{N}^{}u_{\mathrm{tr}}\|\mathrm{T}^{}\mathrm{M}^{}u_{\mathrm{tr}},\hfill \end{array}$$ since $`\mathrm{N}^{}u_{\mathrm{tr}}𝒟(\mathrm{S}^{})`$ and $`\mathrm{M}^{}u_{\mathrm{tr}}𝒟(\mathrm{T}^{})`$ (cf. \[Bol, Prop2.1\]). ∎ ###### Proof of Lemma 2.6. 1. Since $`v_k𝒟(\mathrm{T}_\mathrm{u})=𝒟(\mathrm{H})`$ for all $`k`$ , where $`\mathrm{T}_u=\mathrm{VH}`$ is the polar decomposition of $`\mathrm{T}_u`$, we can write the trace, defined in (2.1), as follows ($`\mathrm{E}_\lambda `$ is the spectral measure of $`\mathrm{H}`$): $$\mathrm{tr}(\mathrm{T}_u^{}\mathrm{T}_u)=\mathrm{tr}(\mathrm{H}^2)=\underset{k}{}\lambda ^2d\mathrm{E}_\lambda v_k^2=\underset{k}{}\mathrm{H}v_k^2=\underset{k}{}\mathrm{T}_uv_k^2.$$ Since the $`[v_k]`$ are mutually orthogonal, we have $$\begin{array}{cc}\hfill \mathrm{tr}(\mathrm{T}_u^{}\mathrm{T}_u)& =\underset{k}{}\mathrm{T}_uv_k\|\mathrm{T}_uv_k\hfill \\ & =\underset{k,j}{}\mathrm{T}_uv_k\|\mathrm{T}_uv_j\hfill \\ & =\underset{k}{}\mathrm{T}_uv_k^2=u^2<\mathrm{}.\hfill \end{array}$$ Further $$\begin{array}{cc}\hfill \mathrm{tr}(\mathrm{T}_u\mathrm{T}_u^{})& =\underset{j}{}\mathrm{T}_u^{}v_j\|\mathrm{T}_u^{}v_j\hfill \\ & =\underset{j}{}\underset{k,i}{}\mathrm{T}_{jk}^{}\delta ^{ji}u_{\mathrm{tr}}\|\mathrm{T}_{jk}^{}\delta ^{ji}u_{\mathrm{tr}}\hfill \\ & =\underset{k,i}{}\mathrm{T}_{ik}^{}u_{\mathrm{tr}}\|\mathrm{T}_{ik}^{}u_{\mathrm{tr}}\hfill \\ & =\underset{k}{}\underset{i}{}\mathrm{T}_{ik}u_{\mathrm{tr}}\|\mathrm{T}_{ik}u_{\mathrm{tr}}\hfill \\ & =\underset{k}{}\mathrm{T}_uv_k\|\mathrm{T}_uv_k\hfill \\ & =\mathrm{tr}(\mathrm{T}_u^{}\mathrm{T}_u).\hfill \end{array}$$ 2. Let $`u`$ be cyclic. Then there are $`\mathrm{M}^{(n)}=(\mathrm{M}_{ik}^{(n)})`$ with $$\underset{n\mathrm{}}{lim}\mathrm{M}^{(n)}u=v$$ for every $`v=(\mathrm{S}_{ij}u_{\mathrm{tr}})_i^j𝒦`$, where $`\mathrm{S}=(\mathrm{S}_{ij})_0`$. This means, using Proposition 2.7 and Corollary 2.5, $$\begin{array}{cc}\hfill 0\stackrel{\mathrm{}n}{}& \underset{i,j}{}\underset{k}{}\mathrm{M}_{ik}^{(n)}\mathrm{T}_{kj}u_{\mathrm{tr}}\mathrm{S}_{ij}u_{\mathrm{tr}}^2\hfill \\ \hfill =& \underset{i,j}{}(\underset{k}{}\mathrm{M}_{ik}^{(n)}\mathrm{T}_{kj}u_{\mathrm{tr}}^22\underset{k}{}\mathrm{Re}\mathrm{M}_{ik}^{(n)}\mathrm{T}_{kj}u_{\mathrm{tr}}\|\mathrm{S}_{ij}u_{\mathrm{tr}}+\mathrm{S}_{ij}u_{\mathrm{tr}}^2)\hfill \\ \hfill =& \underset{i,j}{}(\underset{k}{}\mathrm{T}_{kj}^{}(\mathrm{M}_{ik}^{(n)})^{}u_{\mathrm{tr}}^22\underset{k}{}\mathrm{Re}\mathrm{T}_{kj}^{}(\mathrm{M}_{ik}^{(n)})^{}u_{\mathrm{tr}}\|\mathrm{S}_{ij}^{}u_{\mathrm{tr}}+\mathrm{S}_{ij}^{}u_{\mathrm{tr}}^2)\hfill \\ \hfill =& \underset{i,j}{}\underset{k}{}\mathrm{T}_{kj}^{}(\mathrm{M}_{ik}^{(n)})^{}u_{\mathrm{tr}}\mathrm{S}_{ij}^{}u_{\mathrm{tr}}^2,\hfill \end{array}$$ i.e., since $`(\mathrm{T}_{ki}^{})_{i,k}\mathrm{T}_u^{}`$ and $`_0u_{\mathrm{tr}}`$ is dense in $`𝒦`$, $`\mathrm{T}_u^{}`$ has dense range, i.e. $`\mathrm{T}_u`$ is injective. Let now $`\mathrm{T}_u`$ be injective and $`\mathrm{M}^{^{}}=(\mathrm{M}_{}^{^{}}{}_{}{}^{ij})^{^{}}`$ with $`\mathrm{M}^{^{}}u=0`$. Now $$\mathrm{M}^{^{}}u=(\underset{j}{}\mathrm{M}_{}^{^{}}{}_{}{}^{ij}\mathrm{T}_{kj}u_{\mathrm{tr}})_{ik}=0,$$ and $$\begin{array}{cc}\hfill 0=\mathrm{M}^{^{}}u^2& =\underset{i,k}{}\underset{j}{}\mathrm{M}_{}^{^{}}{}_{}{}^{ij}\mathrm{T}_{kj}u_{\mathrm{tr}}\hfill \\ & =\underset{i,k}{}\underset{j}{}\mathrm{T}_{kj}\mathrm{M}_{}^{^{}}{}_{}{}^{ij}u_{\mathrm{tr}}=\mathrm{T}_uv,\hfill \end{array}$$ where $`v:=(\mathrm{M}_{}^{^{}}{}_{}{}^{ij}u_{\mathrm{tr}})_i^j=_k\mathrm{M}^{^{}}v_k𝒟(\mathrm{T}_\mathrm{u})`$ ($`\mathrm{T}_u`$ is closed), hence $`\mathrm{T}_uv=0`$, and, since $`\mathrm{T}_u`$ is injective, $`v_i^j=\mathrm{M}_{}^{^{}}{}_{}{}^{ij}u_{\mathrm{tr}}=0`$ for all $`i,j`$. Because $`u_{\mathrm{tr}}`$ is cyclic for $`𝒯`$ hence separating for $`𝒯^{^{}}`$, $`\mathrm{M}_{}^{^{}}{}_{}{}^{ij}=0`$ for all $`i,j`$, s.t. $`\mathrm{M}^{^{}}=0`$. 3. Let $`u`$ be separating. This means that $`u`$ is cyclic for $`^{^{}}`$. Then there are $`\mathrm{M}_{(n)}=(\mathrm{M}_{(n)}^{ik})^{^{}}`$ and $$\underset{n\mathrm{}}{lim}\mathrm{M}_{(n)}u=v$$ for every $`v=(\mathrm{S}_{ij}u_{\mathrm{tr}})_i^j𝒦`$, where $`(\mathrm{S}_{ij})_0^{^{}}`$ ($`i,j`$). This means $$\begin{array}{cc}\hfill 0\stackrel{\mathrm{}n}{}& \underset{i,j}{}\underset{k}{}\mathrm{M}_{(n)}^{jk}\mathrm{T}_{ik}u_{\mathrm{tr}}\mathrm{S}_{ij}u_{\mathrm{tr}}^2\hfill \\ \hfill =& \underset{i,j}{}\underset{k}{}\mathrm{T}_{ik}\mathrm{M}_{(n)}^{jk}u_{\mathrm{tr}}\mathrm{S}_{ij}u_{\mathrm{tr}}^2.\hfill \end{array}$$ Since $`_0^{^{}}u_{\mathrm{tr}}`$ is dense in $`𝒦`$ we have proven that $`\mathrm{T}_u`$ has dense range. For the converse read the argument backwards. 4. This follows from 2. and 3. ###### Remark 2.1. Also here, as in the finite case, the finite trace condition of Lemma 2.6 is not only necessary but also sufficient for an operator being the operator associated with a vector in the sense of Definition 2.1. Suppose that $`\mathrm{tr}(\mathrm{T}^{}\mathrm{T})<\mathrm{}`$ with $`\mathrm{T}\eta `$. Then $$\begin{array}{cc}\hfill \mathrm{}>\mathrm{tr}(\mathrm{T}^{}\mathrm{T})& =\mathrm{tr}(\mathrm{H}^2)\hfill \\ & =\underset{k}{}\lambda ^2d\mathrm{E}_\lambda v_k\hfill \end{array}$$ hence $$\lambda ^2d\mathrm{E}_\lambda v_k<\mathrm{}k,$$ i.e. $`v_k𝒟(\mathrm{H})=𝒟(\mathrm{T})`$, and $$\underset{k}{}\mathrm{T}v_k^2=\underset{k}{}\mathrm{H}v_k^2=\underset{k}{}\lambda ^2d\mathrm{E}_\lambda v_k<\mathrm{}.$$ This shows that the assumptions of Lemma 2.3 are fulfilled. ###### Corollary 2.8. $``$ possesses a cyclic and separating vector $`u_0𝒦`$. ###### Proof. Set $`\mathrm{T}:=(\delta _{ij}j^2\mathrm{Id})_{i,j}`$ or $`u_0:=_jj^2v_j`$. Then $`\mathrm{T}`$ fulfills the conditions of Lemma 2.3 and is invertible, s.t. from Lemma 2.6 follows that $`u_0`$ is cyclic and separating. ∎ In the last step of this subsection we show that the model we have just treated is really representative for the general situation, in the sense that all infinite type $`I`$ or type $`II`$ factors can be considered as a matrix algebra of finite type $`I`$ or type $`II`$ factors, resp. This is shown by the next ###### Lemma 2.9. Every infinite but semifinite von Neumann factor $`(_0,_0)`$ with cyclic and separating vector $`u_0_0`$ is unitarily equivalent to $`𝒯L(_{\mathrm{}})=:,_{\mathrm{}}_{\mathrm{}}=:𝒦)`$, where $`𝒯`$ is a finite von Neumann factor acting on the Hilbert space $``$ with cyclic and separating vector and $`_{\mathrm{}}`$ is a separable infinite dimensional Hilbert space. ###### Proof. Since $`_0`$ is infinite but semifinite there is a countable orthogonal family of finite equivalent projections $`(\mathrm{E}_n)_n`$ in $`_0`$, s.t. $`\mathrm{E}_n=\mathrm{Id}`$. Now there is a selfadjoint system of matrix units $`(\mathrm{E}_{ab})_{a,b}`$ with $`\mathrm{E}_{aa}=\mathrm{E}_a`$ (cf. \[KR86, 6.6.4\]). This shows that $`_0`$ is isomorphic to $`\stackrel{~}{𝒯}L(_{\mathrm{}})`$ where $`\stackrel{~}{𝒯}:=\{\mathrm{E}_{ab}\}^{^{}}_0`$ and $`\stackrel{~}{𝒯}`$ is isomorphic to every $`\mathrm{E}_n_0\mathrm{E}_n`$ ($`n`$). Since the projections $`\mathrm{E}_n`$ are finite also $`\stackrel{~}{𝒯}`$ is a finite factor. Since $`_0`$ possesses the separating vector $`u_0`$ we can represent the algebras $`\mathrm{E}_n_0\mathrm{E}_n`$ by the GNS representation for the faithful state $`\omega _n`$ induced by the separating vector $`\mathrm{E}_nu_0`$ on a Hilbert space $`_n`$, s.t. the vector $`u_n_n`$ implementing the state $`\omega _n`$ is a cyclic and separating vector for $`\mathrm{E}_n_0\mathrm{E}_n`$. Since all the $`\mathrm{E}_n_0\mathrm{E}_n`$ are isomorphic and they possess in this representation a cyclic and separating vector, they are all unitarily equivalent. This means that we can choose as $`𝒯`$ one of the $`\mathrm{E}_n_0\mathrm{E}_n`$ acting on the representation space $`_n`$. Since the factor $`(𝒯L(_{\mathrm{}})=:,_{\mathrm{}}_{\mathrm{}}=:𝒦)`$ possesses a cyclic and separating vector if $`(𝒯,)`$ does (see Corollary 2.8) and it is isomorphic to $`_0`$ it is unitarily equivalent to $`_0`$. ∎ The results of this section (and the analogues in \[Bol\]) can be subsumed in the next two theorems: ###### Theorem 2.10. Let ($`_0,_0`$) be a semifinite von Neumann factor. Let further $`u_0`$. Then there is exactly one operator $`\mathrm{T}_u\eta _0`$ associated with the vector $`u`$ in the sense of \[Bol, Def 2.1.\] in the finite case and in the sense of Definition 2.1 in the infinite case, resp., having the following properties: 1. $`\mathrm{tr}(\mathrm{T}_u\mathrm{T}_u^{})=\mathrm{tr}(\mathrm{T}_u^{}\mathrm{T}_u)<\mathrm{}`$. 2. $`u`$ is cyclic, iff $`\mathrm{T}_u`$ is injective. 3. $`u`$ is separating, iff $`\mathrm{T}_u`$ has dense range. 4. $`u`$ is cyclic and separating iff $`\mathrm{T}_u`$ is injective and has dense range, i.e. iff $`\mathrm{T}_u`$ is invertible. ###### Proof. The finite case is just Theorem 1.1. In the infinite case the existence and the asserted properties follow from Lemma 2.9 and Lemma 2.6 infinite case, the uniqueness from Corollary 2.4. ∎ ###### Theorem 2.11. Let $`\mathrm{T}\eta _0`$. Then there is a vector $`u_0`$ s.t $`\mathrm{T}=\mathrm{T}_u`$ iff $`\mathrm{tr}(\mathrm{TT}^{})=\mathrm{tr}(\mathrm{T}^{}\mathrm{T})<\mathrm{}`$. ###### Proof of Theorem 2.11. Again the finite case is just Theorem 1.2. In the infinite case the necessarity of the trace condition follows from Theorem 2.10 and the sufficiency from Remark 2.1, resp. ∎ ## 3 Generation of Modular Objects In this section we show how the modular objects of a cyclic and separating vector $`u_0`$ for a semifinite von Neumann factor $`(_0,_0)`$ are related to the operator $`\mathrm{T}_{u_0}`$ constructed in the last section. As in §2 we consider as a model for the infinite but semifinite factor the factor $`𝒯L(_{\mathrm{}})=:,_{\mathrm{}}_{\mathrm{}})=:𝒦)`$, where $`𝒯`$ is a finite factor with cyclic trace vector $`u_{\mathrm{tr}}`$. If $`u_0𝒦`$ is a cyclic and separating vector for $``$, according to Lemma 2.6, there is an invertible operator $`\mathrm{T}_{u_0}\eta `$, s.t. $`u_0=_k\mathrm{T}_{u_0}v_k`$, where $`v_k=(\delta _i^j\delta _{ik}u_{\mathrm{tr}})_i^j`$. Using this operator we can formulate the following analogue to Theorem 1.3: ###### Theorem 3.1. Use the notations from above. Let further $$\mathrm{T}_{u_0}=\mathrm{HV}=(\mathrm{H}_{\mathrm{ij}})_{ij}(\mathrm{V}_{ij})_{ij}$$ be the polar decomposition of $`\mathrm{T}_{u_0}`$. With the conjugation $`\stackrel{~}{\mathrm{J}}`$ defined as $$\stackrel{~}{\mathrm{J}}(\mathrm{M}_{ij}u_{\mathrm{tr}})_i^j:=(\mathrm{M}_{ji}^{}u_{\mathrm{tr}})_i^j:=(\mathrm{JM}_{ji}u_{\mathrm{tr}})_i^j\mathrm{M}=(\mathrm{M}_{ij})_{ij},$$ where $`\mathrm{J}`$ is the conjugation corresponding to the trace vector $`u_{\mathrm{tr}}`$, we can calculate the modular objects $`(\mathrm{\Delta }_0,\mathrm{J}_0)`$ of $`(_0,u_0)`$ as follows: $$\mathrm{J}_0=\stackrel{~}{\mathrm{J}}\mathrm{V}^{}\stackrel{~}{\mathrm{J}}\mathrm{V}\stackrel{~}{\mathrm{J}}=\mathrm{V}\stackrel{~}{\mathrm{J}}\mathrm{V}^{},$$ and $$\mathrm{\Delta }_0=\mathrm{J}_0\mathrm{H}_0^1\mathrm{J}_0\mathrm{H}_0,$$ where $`\mathrm{H}_0=\mathrm{H}^2=\mathrm{T}_{u_0}\mathrm{T}_{u_0}^{}`$. ###### Proof. 1. First we observe that $`\stackrel{~}{\mathrm{J}}\mathrm{R}\stackrel{~}{\mathrm{J}}^{^{}}`$ for every $`\mathrm{R}`$. For let $`\mathrm{R}=(\mathrm{R}_{ij})`$ and $`v=(u_i^j)_i^j=(\mathrm{M}_{ij}u_{\mathrm{tr}})_i^j𝒦`$, $`(\mathrm{M}_{ij})_0`$, then $$\begin{array}{cc}\hfill \stackrel{~}{\mathrm{J}}\mathrm{R}\stackrel{~}{\mathrm{J}}v& =\stackrel{~}{\mathrm{J}}\mathrm{R}(\mathrm{JM}_{ji}u_{\mathrm{tr}})_i^j\hfill \\ & =\stackrel{~}{\mathrm{J}}(\underset{i}{}\mathrm{R}_{ki}\mathrm{JM}_{ji}u_{\mathrm{tr}})_k^j\hfill \\ & =(\mathrm{J}\underset{i}{}\mathrm{R}_{ji}\mathrm{JM}_{ki}u_{\mathrm{tr}})_k^j\hfill \\ & =\underset{:=\mathrm{R}^{^{}}^{^{}}}{\underset{}{(\mathrm{JR}_{ji}\mathrm{J})^{ji}}}(\mathrm{M}_{ki}u_{\mathrm{tr}})_k^i\hfill \\ & =\mathrm{R}^{^{}}v.\hfill \end{array}$$ Further $$\stackrel{~}{\mathrm{J}}\stackrel{~}{\mathrm{J}}v=\stackrel{~}{\mathrm{J}}(\mathrm{JM}_{ji}u_{\mathrm{tr}})_{ij}=(\mathrm{M}_{ij}u_{\mathrm{tr}})_{ij}=v,$$ s.t. $`\stackrel{~}{\mathrm{J}}`$ is an (algebraic) conjugation for $``$. 2. Let now $`\mathrm{T}_{u_0}`$ be bounded ($``$ all the $`\mathrm{T}_{\mathrm{ij}}`$ and $`\mathrm{H}_{ij}`$, resp. are bounded). Then we show that the Tomita operator $`\mathrm{S}`$ defined by $$\mathrm{SA}u_0=\mathrm{A}^{}u_0\mathrm{A}$$ can be written as $$\mathrm{S}=\mathrm{H}^1\mathrm{V}\stackrel{~}{\mathrm{J}}\mathrm{V}^{}\mathrm{H}.$$ (3.1) For this let $`\mathrm{A}=(\mathrm{A}_{ij})_{ij}`$ and $`u_0=(_k\mathrm{H}_{jk}\mathrm{V}_{kl}u_{\mathrm{tr}})_j^l`$. Then $$\mathrm{A}u_0=(\underset{j,k}{}\mathrm{A}_{ij}\mathrm{H}_{jk}\mathrm{V}_{kl}u_{\mathrm{tr}})_i^l$$ and $$\mathrm{A}^{}u_0=(\underset{j,k}{}\mathrm{A}_{ji}^{}\mathrm{H}_{jk}\mathrm{V}_{kl}u_{\mathrm{tr}})_i^l.$$ Now $$\begin{array}{cc}\hfill (\mathrm{H}^1\mathrm{V}\stackrel{~}{\mathrm{J}}\mathrm{V}^{}\mathrm{H})\mathrm{A}u_0& =\mathrm{H}^1\mathrm{V}\stackrel{~}{\mathrm{J}}(\underset{i,j,k,m}{}\mathrm{V}_{mn}^{}\mathrm{H}_{mi}\mathrm{A}_{ij}\mathrm{H}_{jk}\mathrm{V}_{kl}u_{\mathrm{tr}})_n^l\hfill \\ & =\mathrm{H}^1\mathrm{V}(\underset{i,j,k,m}{}\mathrm{JV}_{ml}^{}\mathrm{H}_{mi}\mathrm{A}_{ij}\mathrm{H}_{jk}\mathrm{V}_{kn}u_{\mathrm{tr}})_n^l\hfill \\ & =\mathrm{H}^1\mathrm{V}(\underset{i,j,k,m}{}\mathrm{V}_{kn}^{}\mathrm{H}_{jk}^{}\mathrm{A}_{ij}^{}\mathrm{H}_{mi}^{}\mathrm{V}_{ml}u_{\mathrm{tr}})_n^l\hfill \\ & =\mathrm{H}^1\mathrm{V}(\underset{i,j,k,m}{}\mathrm{V}_{kn}^{}\mathrm{H}_{kj}\mathrm{A}_{ij}^{}\mathrm{H}_{im}\mathrm{V}_{ml}u_{\mathrm{tr}})_n^l\hfill \\ & =\mathrm{H}^1(\underset{i,j,m}{}\mathrm{H}_{nj}\mathrm{A}_{ij}^{}\mathrm{H}_{im}\mathrm{V}_{ml}u_{\mathrm{tr}})_n^l\hfill \\ & =(\underset{i,m}{}\mathrm{A}_{in}^{}\mathrm{H}_{im}\mathrm{V}_{ml}u_{\mathrm{tr}})_n^l=\mathrm{A}^{}u_0,\hfill \end{array}$$ which proves (3.1). Now $`\mathrm{S}^{}=\mathrm{HV}\stackrel{~}{\mathrm{J}}\mathrm{V}^{}\mathrm{H}^1`$ and $$\begin{array}{cc}\hfill \mathrm{\Delta }_0& =\mathrm{S}^{}\mathrm{S}\hfill \\ & =\mathrm{HV}\stackrel{~}{\mathrm{J}}\mathrm{V}^{}\mathrm{H}^1\mathrm{H}^1\mathrm{V}\stackrel{~}{\mathrm{J}}\mathrm{V}^{}\mathrm{H}\hfill \\ & =\mathrm{V}\stackrel{~}{\mathrm{J}}\mathrm{V}^{}\mathrm{H}^2\mathrm{V}\stackrel{~}{\mathrm{J}}\mathrm{V}^{}\mathrm{H}^2\hfill \\ & =\mathrm{J}_0\mathrm{H}_0^1\mathrm{J}_0\mathrm{H}_0.\hfill \end{array}$$ Further $$\mathrm{J}_0\mathrm{\Delta }_0^{1/2}=\mathrm{H}^1\mathrm{J}_0\mathrm{H}=\mathrm{S},$$ and all the assertions are proven in the bounded case. 3. In the last step we approximate the (unbounded) operator $`\mathrm{T}_{u_0}`$ by bounded operators $`\mathrm{T}_n`$ in exactly the same way as in the proof of Theorem 3.1. in \[Bol\] and show the assertions like there also in the unbounded case. ## 4 The Second Simple Class of Solutions of the Inverse Problem In this section we want to use the results of the last two sections to examine the second simple classes of solutions of the inverse problem introduced by Wollenberg in \[Wolb\] for type $`I`$ factors, and considered in \[Bol\] also for type $`II_1`$ factors. For the construction of this class it is crucial that the inverse $`\mathrm{\Delta }_0^1`$ of the modular operator is again a modular operator. To this scope there was shown the following ###### Lemma 4.1. Let $`\mathrm{\Delta }_0=\mathrm{J}_0\mathrm{H}_0^1\mathrm{J}_0\mathrm{H}_0`$ be the decomposition of the modular operator $`\mathrm{\Delta }_0`$, where $`\mathrm{J}_0=\mathrm{JV}^{}\mathrm{JVJ}=\mathrm{VJV}^{}`$ and $`\mathrm{T}_{u_0}=\mathrm{H}_0^{1/2}\mathrm{V}`$ is the operator corresponding to $`u_0`$ (cf. Theorem 3.1). Then $`\mathrm{\Delta }_0^1=\mathrm{J}_0\mathrm{H}_0\mathrm{J}_0\mathrm{H}_0^1`$ and the following is equivalent: 1. ($`\mathrm{\Delta }_0^1,\mathrm{J}_0)`$ are the modular objects w.r.t. a cyclic and separating vector $`u_1_0`$. 2. $$\mathrm{tr}(\mathrm{H}_0^1)<\mathrm{}.$$ (4.1) This lemma can be proven with the same techniques as in \[Bol\] also for the infinite case taking into account Theorem 2.10, Theorem 2.11, and Theorem 3.1. Now we must examine, whether or not the second condition in Lemma 4.1 is fulfilled: ###### Lemma 4.2. For type $`I_{\mathrm{}}`$ and type $`II_{\mathrm{}}`$ factors the condition (4.1) is never true. ###### Proof. Let $`_0`$ now be a type $`I_{\mathrm{}}`$ or $`II_{\mathrm{}}`$ factor and $`\mathrm{T}_{u_0}=\mathrm{H}_0^{1/2}\mathrm{V}`$ the operator corresponding to the cyclic and separating vector $`u_0`$. Let further $`\mathrm{E}_\lambda _0`$ the spectral resolution of $`\mathrm{H}_0`$. Then we can define a positive measure $`\mu _{\mathrm{tr}}`$ on the $`\sigma `$-algebra of Borel sets in $``$, s.t. $$\mathrm{tr}(\mathrm{H}_0)=\lambda 𝑑\mu _{\mathrm{tr}}(\lambda ),$$ where $$\mu _{\mathrm{tr}}(B):=\mathrm{tr}\mathrm{E}(B)$$ for all Borel sets $`B`$. Now $`c:=\mathrm{tr}(\mathrm{H}_0)<\mathrm{}`$. Assume w.l.o.g. $`c=1`$. Then $$1=\lambda 𝑑\mu _{\mathrm{tr}}(\lambda )_{[0,1]}\lambda 𝑑\mu _{\mathrm{tr}}(\lambda )+_{(1,\mathrm{})}𝑑\mu _{\mathrm{tr}}(\lambda ),$$ i.e. $$_{(1,\mathrm{})}𝑑\mu _{\mathrm{tr}}(\lambda )<\mathrm{}.$$ Since $`_0`$ is infinite $`\mathrm{}=\mathrm{tr}(\mathrm{Id})=\mu _{\mathrm{tr}}()`$, i.e. $$\mathrm{}=_\lambda 𝑑\mu _{\mathrm{tr}}(\lambda )=_{[0,1]}𝑑\mu _{\mathrm{tr}}(\lambda )+\underset{<\mathrm{}}{\underset{}{_{(1,\mathrm{})}𝑑\mu _{\mathrm{tr}}(\lambda )}},$$ hence $$_{[0,1]}𝑑\mu _{\mathrm{tr}}(\lambda )=\mathrm{}.$$ Suppose now that also $`\mathrm{tr}(\mathrm{H}_0^1)<\mathrm{}`$, then $$\mathrm{}>_\lambda \lambda ^1𝑑\mu _{\mathrm{tr}}(\lambda )\underset{=\mathrm{}}{\underset{}{_{[0,1]}𝑑\mu _{\mathrm{tr}}(\lambda )}}+_{(1,\mathrm{})}\lambda ^1𝑑\mu _{\mathrm{tr}}(\lambda ),$$ which is a contradiction. ∎ Hence the last lemma shows that for infinite semifinite factors the second class of solutions of the inverse problem can never be constructed. This result was yet obtained by Wollenberg in \[Wolb\] for the type $`I_{\mathrm{}}`$ case. ## 5 The Classification of Solutions in the Pure Point Spectrum Case In this section we want to show the modifications of classification results obtained in \[Bol\]. The definition of the equivalence relation does not use any special properties of the finite factors, and can just be repeated here: ###### Definition 5.1. Two semifinite von Neumann factors $`,𝒩NF__0(\mathrm{\Delta }_0,\mathrm{J}_0,u_0)`$ are called equivalent, $`𝒩`$, if $`NF_𝒩^1(\mathrm{\Delta }_0,\mathrm{J}_0,u_0)`$, i.e. if there exists a unitary operator $`\mathrm{U}`$ on $`_0`$, s.t. $`=\mathrm{U}𝒩\mathrm{U}^{}`$, $`\mathrm{U}`$ commutes with $`\mathrm{\Delta }_0`$ and $`\mathrm{J}_0`$ and $`\mathrm{U}^{}u_0=\pm u_0`$ (For the definition of the class $`NF_𝒩^1(\mathrm{\Delta }_0,\mathrm{J}_0,u_0)`$ see \[Bol\]). Also the next lemmas can be formulated and proved in exactly the same way as in the finite case. Assume in the following that $`\mathrm{H}_0`$ has pure point spectrum, i.e. $`\mathrm{H}_0=_{kK}\mu _k\mathrm{E}_k`$ where the $`\mu _k`$ ($`kK`$) are the eigenvalues of $`\mathrm{H}_0`$ and $`\mathrm{E}_k_0`$ are the corresponding (orthogonal) eigenprojections with $`m_k:=\mathrm{tr}\mathrm{E}_k=:D__0(\mathrm{E}_k)`$ their von Neumann dimension. Then we have for $`\mathrm{\Delta }_0`$ the following decomposition $$\begin{array}{cc}\hfill \mathrm{\Delta }_0& =\mathrm{H}_0\mathrm{J}_0\mathrm{H}_0^1\mathrm{J}_0\hfill \\ & =\underset{k,lK}{}\mu _k\mu _l^1\mathrm{E}_k\mathrm{J}_0\mathrm{E}_l\mathrm{J}_0\hfill \\ & =\underset{jJ}{}\lambda _j\mathrm{F}_j,\hfill \end{array}$$ (5.1) where the $`\lambda _j`$ ($`jJ`$) are the eigenvalues of $`\mathrm{\Delta }_0`$ and $`\mathrm{F}_j`$ are the corresponding eigenprojections. Now ###### Lemma 5.1. With the notations introduced above we can compute the spectrum of $`\mathrm{\Delta }_0`$ in the following way: $$\{\lambda _j\|jJ\}=\{\mu _k\mu _l^1\|k,lK\}jJ$$ (5.2) and $$n_j=\underset{\mu _k\mu _l^1=\lambda _j}{}m_km_ljJ\text{ if }_0\text{ is type }I\text{,}$$ (5.3a) $$n_j=\mathrm{}jJ\text{ if }_0\text{ is type }II\text{,}$$ (5.3b) where $`n_j:=D_{L(_0)}(\mathrm{F}_j)`$ with $`D_{L(_0)}(\mathrm{F}_j)`$ the dimension function in the type $`I_{\mathrm{}}`$ factor $`L(_0)`$, which corresponds to the normalized Hilbert space dimension. ###### Lemma 5.2. If there are two solutions of the inverse problem $`_1`$, $`_2`$ s.t. the corresponding selfadjoint operators $`\mathrm{H}_1`$ and $`\mathrm{H}_2`$ have the same eigenvalues modulo a positive constant $`c>0`$ and same (von Neumann) multiplicities, then $`_1_2`$. ###### Lemma 5.3. If there are two equivalent solutions $`_1`$, $`_2`$ of the inverse problem with the corresponding positive operators $`\mathrm{H}_1`$ and $`\mathrm{H}_2`$, resp., (having pure point spectrum) then $`\mathrm{H}_1`$ and $`\mathrm{H}_2`$ have the same eigenvalues (up to a positive constant) and von Neumann multiplicities, i.e. they are unitarily equivalent in $`_0`$. The only difference to the finite case is shown by the next ###### Lemma 5.4. Let $`(\mu _k,m_k)_{kK}`$ be a sequence of pairs of positive reals $`\mu _k>0`$ and $`m_k>0`$, s.t. $$m_k\text{ if }_0\text{ is type }I_{\mathrm{}},$$ (5.4a) $$m_k_{>0}\text{ if }_0\text{ is type }II_{\mathrm{}},$$ (5.4b) and $$\underset{kK}{}m_k=\mathrm{}$$ (5.4c) and $$\underset{kK}{}m_k\mu _k=1$$ (5.4d) and the relations (5.2) and (5.3) are fulfilled. Then there exists a solution $`=\mathrm{U}_0\mathrm{U}^{}NF__0(\mathrm{\Delta }_0,\mathrm{J}_0,u_0)`$, s.t. $`\mathrm{U}^{}\mathrm{\Delta }_0\mathrm{U}=\mathrm{HJ}_0\mathrm{H}^1\mathrm{J}_0`$ and $`\mathrm{H}`$ has the eigenvalues and multiplicities $`(\mu _k,m_k)_{kK}`$ (cf. \[Wolb, prop.4.1\]). For the proof we need the following auxiliary results: ###### Proposition 5.5. If $`(m_k)`$ is countable family of positive reals with $`m_k=\mathrm{}`$, then there exists in a type $`II_{\mathrm{}}`$ von Neumann factor $``$ a family of pairwise orthogonal projections $`(\mathrm{E}_k)`$, s.t. $`D(\mathrm{E}_k)=m_k`$ for every $`k`$. ###### Proof. We construct the $`\mathrm{E}_k`$ inductively: Since the range of $`D_{}`$ is all of $`_0`$ (cf. \[KR86, 8.4.4\]) there is a projection in $``$, s.t. $`D(\mathrm{E}_1)=m_1`$. Suppose now that for $`N`$ the $`\mathrm{E}_k`$ are pairwise orthogonal with $`D__0(\mathrm{E}_k)=m_k`$ ($`1k<N`$). Setting $`\mathrm{F}_N:=\mathrm{Id}_{k=1}^N\mathrm{E}_k`$ the restricted algebra $`\mathrm{F}_N\mathrm{F}_N`$ is again a type $`II`$ factor, finite, if $`\mathrm{F}_N`$ is finite, and infinite, if $`\mathrm{F}_N`$ is infinite (cf. \[KR86, Ex. 6.9.16\]) with the dimension function $$D_N(\mathrm{F}_n\mathrm{EF}_N):=D__0(\mathrm{F}_n\mathrm{EF}_N)/D(\mathrm{F}_N)\mathrm{F}_n\mathrm{EF}_N\mathrm{F}_N\mathrm{F}_N,$$ if $`\mathrm{F}_N`$ is finite, and $`D_N=D__0`$ else, where $$D__0(\mathrm{F}_N)=D__0(\mathrm{Id}\underset{k=1}{\overset{N}{}}\mathrm{E}_k)=1\underset{k=1}{\overset{N}{}}D__0(\mathrm{E}_k)m_N.$$ With the same argument as above there is again a projection $`\mathrm{E}_N\mathrm{F}_N\mathrm{F}_N`$, s.t. $`D_N(\mathrm{E}_N)=D(\mathrm{F}_N)^1m_N1`$, if $`\mathrm{F}_N`$ is finite, and $`D_N(\mathrm{E}_N)=m_N`$ else. In both cases $`D__0(\mathrm{E}_N)`$ and $`\mathrm{E}_N<\mathrm{F}_N\mathrm{E}_k`$ ($`1k<N`$). ∎ Now the proof of Lemma 5.4 is the same as in \[Bol\]. ###### Remark 5.1. (5.4) show that in the infinite case we have always an infinite set of eigenvalues with $`0`$ as cummulation point, i.e. $`K=`$ and $`0`$ is in the spectrum $`\sigma (\mathrm{H})`$ of $`\mathrm{H}`$. Now we can summarize the lemmas of this section in the following ###### Theorem 5.6. Let $`_0`$ be a semifinite von Neumann factor with cyclic and separating vector $`u_0`$ and $`\mathrm{T}_{u_0}=\mathrm{H}_0^{1/2}\mathrm{V}`$ the operator corresponding to $`u_0`$. If $`\mathrm{H}_0`$ has pure point spectrum, also $`\mathrm{\Delta }_0`$ has it. In this case let $`(\lambda _j)`$ $`(jJ)`$ be the eigenvalues of $`\mathrm{\Delta }_0`$. Then 1. Two solutions $`_1,M_2NF__0(\mathrm{\Delta }_0,\mathrm{J}_0,u_0)`$ of the inverse problem with corresponding invertible operators $`\mathrm{H}_i\eta _0`$ $`(i=1,2)`$ having pure point spectrum are equivalent iff $`\mathrm{H}_1`$ and $`\mathrm{H}_2`$ have the same eigenvalues and (von Neumann) multiplicities. 2. A positive invertible operator $`\mathrm{H}\eta _0`$ with pure point spectrum gives rise to a solution of the inverse problem iff its eigenvalues and multiplicities satisfy (5.2), (5.3), and (5.4). 3. When the corresponding operators $`\mathrm{H}`$ has pure point spectrum the equivalence classes of $``$ are completely classified by the spectrum of the corresponding operators, i.e. by sequences of pairs of positive reals $`(\mu _k,m_k)`$ satisfying (5.2), (5.3), and (5.4). ###### Example 5.1. Here we want to give some examples to illustrate Theorem 5.6. 1. In \[Wolb\] you can find some examples for the type $`I`$ case. 2. Let $$(\mathrm{},10^3,10^2,10^1,1,10,10^2,10^3,\mathrm{})$$ be the eigenvalues of a modular operator for a type $`II_{\mathrm{}}`$ factor. Then $$((c_11,1),(c_110^1,1),(c_110^2,1),(c_110^3,1),\mathrm{})$$ and $$((c_21,1),(c_210^1,1),(c_210^3,1),(c_210^5,1),\mathrm{})$$ characterize two different classes of solutions of the inverse problem, i.e. they both satisfy (5.2), (5.3), and (5.4), where $`c_i`$ ($`i=1,2`$) are appropriate chosen constants. This shows that in this case there are more than the simple classes of solutions of the inverse problem. 3. Let $`(\mu _k,m_k)_k`$ characterize a class of solutions of the inverse problem in the type $`II_{\mathrm{}}`$ case, where $`m_lm_k`$ for at least one pair $`k,l`$, then for every finite permutation $`\sigma `$ of $``$ interchanging $`k`$ and $`l`$ also $`(c\mu _k,m_{\sigma (k)})`$ characterize another class of solutions of the inverse problem ($`c>0`$ a norming constant) which is really a new one. 4. Let again $`(\mu _k,m_k)_k`$ be a solution of the inverse problem in the type $`II_{\mathrm{}}`$ case, and let $`k,l`$ be a pair of indices and $`ϵ>0`$. Then we get another class by adding $`ϵ`$ to $`m_k`$ and subtracting it from $`m_l`$ where again we have really a new class. ###### Remark 5.2. 1. Example 5.1.3 and Example 5.1.4 shows that in the type $`II_{\mathrm{}}`$ case, when $`\mathrm{H}_0`$ has pure point spectrum, we can always construct a second class of solutions, different from the simple class discussed in §4, i.e. $`NF__0NF__0^1`$, in contrast to the type $`I`$ case, where for modular operators with generic spectrum we have $`NF__0=NF__0^1`$ (cf. \[Wolb\]). 2. Unfortunately the classification result presented here applies only to operators with pure point spectrum. Whereas in general there are also operators with more complicated spectrum (cf. \[Bol, Remark 4.1\]), for type $`I`$ factors this is no restriction, since all operators generating modular operators are trace class operators, hence have pure point spectrum. Acknowledgements The author thanks professor M. Wollenberg for discussing and his usefull hints and the DFG and the Graduiertenkolleg for financial support.
warning/0004/cond-mat0004431.html	ar5iv	text	# Comment on “Phase ordering in chaotic map lattices with conserved dynamics” ## I Introduction The dynamical phenomenon of domain growth occurs in many different physical contexts. Once a fairly well established subject , it has been recently the center of renewed interest for several distinct reasons: one is the advent of new quantifiers of the associated dynamical scaling regimes, such as first-passage or persistence exponents, i.e. the rate of algebraic decay of the probability for a given point in space to have remained in the same phase since some initial time . Another reason is the natural question of the extent of known “universality classes” to new types of systems, e.g. spatio-temporally chaotic ones. In this context, the recent study of the ordering properties of chaotic coupled map lattices (CMLs) possessing several symmetric phases in competition brought up some intriguing results : in the simple case of two competing phases and a non-conserved order parameter, the “normal” growth law $`L(t)t^{1/z}`$ with $`z=2`$, where $`L(t)`$ is the single lengthscale characterizing the coarsening pattern, was observed but with some exponent $`z2`$ continuously varying with parameters. However, this was later shown to be only a (slow) transient behavior due to the non-trivial effect of space-discretization in these deterministic systems. For larger lattices and longer times than those considered in , $`L`$ was shown to behave normally when plotted against $`t^{1/2}`$. In a recent Rapid Communication, Angelini, Pellicoro, and Stramaglia (APS) , motivated by the above study of non-conserved order parameter CMLs, presented a class of sequentially-updated lattices of chaotic maps designed to investigate the case where the order parameter is locally conserved. In this case, $`L`$ is also expected to grow algebraically with time, but with $`z=3`$ . APS claim, however, that larger exponents are commonly found. Here we show that APS were misled by their treatment of data and that in fact the normal ($`z=3`$) growth law is observed in all cases. We argue, moreover, that fully deterministic, synchronously-updated, coupled map lattices which conserve the order parameter can be easily constructed following the ideas of Oono and Puri , and we show that these systems behave very smoothly, enabling the precise measurement of persistence exponents in this context. ## II Revisiting APS results ### A The hybrid map lattices of APS The lattices of maps introduced by APS are hybrid in several ways: a given local map $`f`$ is first applied to all sites $`x_i`$ of the lattice (a deterministic and synchronous operation), then pairs of nearest neighbors are sequentially and regularly visited and swapped probabilistically. (The regularity of the sweeps of the lattice is at the origin of the anisotropy of the domains in Fig. 1 of .) Furthermore, such systems are not “coupled” maps as in usual CMLs, since the values taken by the sites are not influenced by the swaps (they are always taken according to the invariant measure of the local map). These systems are designed to mimick Ising systems (with fluctuating couplings) corresponding to the “spins” $`\sigma _i=\mathrm{sgn}(x_i)`$. The local map $`f`$ has not much importance, and it is convenient to choose, following , an odd map of the $`[1,1]`$ interval with two symmetric attractors.. The energy of one configuration is given by $`E=_{i,j}x_ix_j`$ where the sum is over nearest-neighbor pairs. The exchange probability reads $`P_{\mathrm{swap}}=1/(1+\mathrm{exp}(\beta \mathrm{\Delta }E)`$ where $`\beta `$ is the inverse temperature and $`\mathrm{\Delta }E`$ is the energy change of the swap. The zero-temperature limit is deterministic: swaps are effective if and only if they decrease the energy. ### B Domain growth is normal We have performed numerical simulations of the APS system at zero temperature with the piecewise linear local map used in : $$f(x)=\{\begin{array}{ccc}\mu X\hfill & \mathrm{if}\hfill & X[1/3,1/3]\hfill \\ 2\mu /3\mu X\hfill & \mathrm{if}\hfill & X[1/3,1]\hfill \\ 2\mu /3\mu X\hfill & \mathrm{if}\hfill & X[1,1/3]\hfill \end{array}$$ (1) with $`\mu =1.9`$. Coarsening occurs, with, again, a strong anisotropy due to the mode of update. The growth of $`L`$, defined as the width at mid-height of the two-point autocorrelation function, is slow at short times, but then reaches the expected $`t^{1/3}`$ behavior (Fig. 1a), contrary to the claims of APS. The short-time behavior may be mistaken for anomalously slow algebraic growth (with an exponent close to the value $`1/z=0.07`$ reported by APS) when logarithmic scales are used, but a closer inspection shows a systematic increase of the local exponent (Fig. 1b). As a matter of fact, the system is so anisotropic that domains elongate in time (Fig. 1a). This is due to our choice of updating $`x`$-wise pairs before $`y`$-wise pairs. Alternating this order would presumably suppress this effect. Runs of the same system at finite temperatures indicate that domain growth is faster and that the crossover to the $`z=3`$ behavior occurs sooner. We are confident that similar results hold for the complicated local map mostly used by APS. As a conclusion, Fig. 3 of has to be replaced by the variation of the prefactor of the $`Lt^{1/3}`$ law, similarly to the final conclusions of for the non-conserved order-parameter case. ### C Persistence scaling is hard to measure The persistence probability $`p(t)=\mathrm{Prob}\{\sigma _i(t^{})=\sigma _i(t_0),t^{}[t_0,t]\}`$ is usually observed to decay algebraically with time ($`p(t/t_0)^\theta `$) in systems with algebraic growth laws. But persistence scaling for conserved order parameter systems is notoriously difficult to observe . An additional difficulty lies in the fact that the available models only show coarsening at finite temperatures, so that one has to resort to block-scaling of persistence. This is also the case of APS systems, even at zero temperature, since the chaotic fluctuations of the “couplings” amount to a finite temperature. This, by the way, is the reason why APS systems coarsen in this case. Needless to say, the estimates of persistence exponents presented in are then highly unreliable, if only because of the slow crossover for the growth law of $`L`$. Ideally, since persistence is a complex quantity involving all times since the reference time $`t_0`$, one should in principle choose $`t_0`$ in the asymptotic scaling regime and simulate the system up to times $`tt_0`$. Given the typical values of the crossover times (Fig. 1) this is hardly possible. Another difficulty for the APS systems is the possible influence of their strong anisotropy on the persistence exponent $`\theta `$ . Rather than trying to measure properly persistence scaling in APS systems, a possible but difficult task, we now turn ourselves to truly deterministic models, i.e. regular coupled map lattices, which are devoid of the drawbacks underlined above for APS systems. ## III Well-behaved deterministic models ### A Oono-Puri style CMLs Deterministic models for phase ordering of conserved systems were introduced by Oono and Puri . Dynamical scaling was observed to hold for these CMLs, with a growth law compatible with the expected $`z=3`$. We now present similar models which, in addition, can be constructed for any local map. A usual CML, such as that studied in for the non-conserved order parameter case, can be written: $$x_i^{t+1}=(x_i^t)(1𝒩g)f(x_i^t)+g\underset{ji}{}f(x_j^t)$$ (2) where $`g`$ is the coupling strength, $`𝒩`$ is the number of neighbors in the choosen coupling range, and the sum is over these neighbors. Following , a CML conserving exactly the continuous local field can easily be constructed as: $$x_i^{t+1}=(x_i^t)\frac{1}{𝒩}\underset{ji}{}\left((x_j^t)x_j^t\right)$$ (3) where the last term corresponds to the extra Laplacian in the Cahn-Hilliard equation. Conservation of the continuous field $`x`$ is obvious. On the other hand, strictly speaking, the discrete field $`\sigma =\mathrm{sgn}(x)`$ is not exactly conserved and fluctuates slightly, because the last term in (3) may change the sign of sites situated in domain walls. The synchronous mode of update prevents excessive (i.e. other than lattice-derivated) anisotropy. The above structure insures that “true” zero-temperature regimes are observed if the local map possesses two disjoint attractors. Changing the nature of these attractors (fixed points, limit cycles, chaotic sets), one can study competition between phases of different nature. ### B An example We now present results obtained on a particular case of the models defined above. More comprehensive results will be reported elsewhere . For simplicity reasons, we again choose the map given by Eq. (1). For “extra smoothness”, the Moore neighborhood on the square lattice ($`𝒩=8`$ neighbors of equal weight) was used. As in , domain walls are strictly pinned for small $`g`$. For too strong coupling, on the other hand, antiferromagnetic-like phases appear. There is, however, an intermediate range of $`g`$ values for which domain growth proceeds forever between two weakly-chaotic phases. The expected $`z=3`$ law is then easily observed even at short times and in log-log scales (Fig. 2a). The above CMLs reveal their strongest advantage when persistence scaling is considered. As already noticed, they show normal coarsening with true zero temperature. This allows to avoid studying block persistence scaling, a somewhat tedious task. Figures 2b shows persistence decay for different reference times $`t_0`$. Nice scaling is easily observed. This constitutes, to our knowledge, the first clean evidence of algebraic decay of persistence in a two-dimensional conserved order parameter system. Our results give an exponent $`\theta 0.25(2)`$, i.e. a value larger than that observed for the non-conserved order parameter case (for which $`\theta 0.200.22`$ ). The above CMLs constitute an excellent base for a reliable study of persistence scaling in conserved order parameter systems. Ongoing work is probing the degree of universality of the persistence exponent measured above, alongside a similar study for the non-conserved order parameter case for which this issue is still unresolved . ## IV Conclusion In this Comment, we showed how APS were misled in their interpretation of simulation data and that their conclusions about possible non-trivial values of the dynamical exponent $`z`$ in chaotic systems with conserved order parameter dynamics do not hold. To a large extent, these systems can be seen as too close to the Ising model with Kawasaki dynamics, which is well-known to be difficult to study numerically (although $`z=3`$ scaling is now well documented ). We introduced a class of CMLs which are devoid of all the problems encountered in APS systems and which show normal scaling already at early times. Moreover, these systems also present nice scaling behavior for the persistence probability, whereas similar investigations in APS systems are riddled with problems. We showed unambiguously that two-dimensional systems with conserved order parameter domain growth show an algebraic decay of persistence. Future work will try to assess the universality of both the so-called Fisher-Huse exponent and the persistence exponent $`\theta `$ in such systems.
warning/0004/nucl-th0004053.html	ar5iv	text	# nucl-th/0004053 LA-UR-00-2054 HYBRID BARYON SIGNATURES ## 1 What is a hybrid baryon? Historically a low–lying hybrid baryon was defined as a three quark – gluon composite. However, from the viewpoint of the Lagrangian of Quantum Chromodymanics (QCD) this definition is non–sensical. This is because gluons are massless, and hence there is no reason not to define a hybrid baryon, for example, as a three quark – two gluon composite. Neither is one possibility distinguishable from the other, since strong interactions mix the possibilities. The place where this definition of a hybrid baryon is most useful is in large $`Q^2`$ deep inelastic scattering, where a Fock state expansion of a state can rigorously be defined, and one can at least talk about the three quark – gluon component of such a state.$`^\mathrm{?}`$ However, in other situations the definition becomes perilous. A case in point is recent work on large $`N_c`$ hybrid baryons, where their properties depend critically on the fact that the gluon is in colour octet, and hence the three quarks in colour octet, so that the entire state is colour singlet.$`^\mathrm{?}`$ The bag model circumvents the objections raised against this definition, since gluons become massive due to their confinement inside the bag.$`^\mathrm{?}`$ More recently, a low–lying hybrid baryon was defined as three quarks moving on the low–lying excited adiabatic potential.$`^\mathrm{?}`$ From the viewpoint of QCD this can be a perfectly sensible definition. One can always evaluate the energy of a system of three fixed quarks as a function of the three quark positions, called the adiabatic potential. A calculation along these lines has been performed in flux–tube models$`^\mathrm{?}`$ and a first attempt has been made in lattice QCD.$`^\mathrm{?}`$ The three quarks are then allowed to move in a three–body equation, typically a non–relativistic Schrödinger equation. Treating a three quark system via this two step process is called the adiabatic or Born–Oppenheimer approximation. Note that it is not appropriate to talk about the case where the quarks are actually infinitely heavy, because of the lack of kinetic energy. The adiabatic approximation is expected to become exact if the quark masses are much greater than the scale of the strong interactions $`\mathrm{\Lambda }_{QCD}`$. The criterion for the validity of the adiabatic approximation is that the slow degrees of freedom (the quarks) should move much slower than the fast degrees of freedom (the gluons). It is possible to argue that for conventional baryons the relative velocities of quarks behave like the strong coupling constant $`\alpha _S`$ as the quarks become heavier.$`^\mathrm{?}`$ Because of the asymptotic freedom of QCD, the quark relative velocities go to zero, ensuring the validity of the adiabatic approximation. In fact, this is the basis for the NRQCD expansion. However, since $`\alpha _S`$ goes to zero only logarithmically, one may need quarks heavier than the bottom quark for the adiabatic approximation to be valid. Depending on the shape of the adiabatic potential the possibility of an NRQCD expansion for hybrid baryons,$`^\mathrm{?}`$ and hence the validity of an adiabatic approximation, may be in jeopardy. Here we point out for the first time that for specific dynamics, the adiabatic approximation can be exact, even for light quarks, if one redefines the adiabatic potential suitably. We call this the “redefined adiabatic approximation”, which employs a “redefined adiabatic potential”. It was noted in a flux–tube model that “For light quarks almost all corrections may be incorporated into a redefinition of the potentials. Mixing between \[new\] potentials is of the order of 1%”.$`^\mathrm{?}`$ We develop the following technique for obtaining the redefined potential. The quark positions are still fixed relative to each other, but the quark and gluon positions are not defined relative to the quark positions, but relative to the centre of mass of the quarks and the gluons, as recently pioneered.$`^\mathrm{?}`$ For dynamics where the redefined adiabatic approximation is exact one can rigorously define a low–lying hybrid baryon as three quarks moving in the low–lying redefined excited adiabatic potential. ## 2 Redefined adiabatic approximation Consider the non–relativistic hamiltonian for three quarks at positions $`𝐫_1,𝐫_2,𝐫_3`$ and a junction at position $`𝐫_4`$. The junction represents the gluons. The classical hamiltonian with a simple harmonic oscillator potential is $$H=\frac{1}{2}M(\dot{𝐫}_1^2+\dot{𝐫}_2^2+\dot{𝐫}_3^2)+\frac{1}{2}m\dot{𝐫}_4^2+\frac{1}{2}k((𝐫_1𝐫_4)^2+(𝐫_2𝐫_4)^2+(𝐫_3𝐫_4)^2)$$ (1) where $`M,m`$ and $`k`$ are the mass of the quarks, the mass of the junction and coupling constant respectively. The first two and the last terms are the kinetic and potential energy terms respectively. ### 2.1 Exact solution Making the variable transformations $$𝝆=\frac{1}{\sqrt{2}}(𝐫_1𝐫_2)𝝀=\frac{1}{\sqrt{6}}(𝐫_1+𝐫_22𝐫_3)$$ $$𝝈=\frac{2}{\sqrt{3}}\frac{m}{m+3M}(𝐫_1+𝐫_2+𝐫_33𝐫_3)$$ (2) one obtains the diagonal hamiltonian $$H=\frac{1}{2}M(\dot{𝝆}^2+\dot{𝝀}^2+\frac{1}{4}(1+3\frac{M}{m})\dot{𝝈}^2)+\frac{1}{2}k(𝝆^2+𝝀^2+\frac{1}{4}(1+3\frac{M}{m})^2𝝈^2)$$ (3) when working in the overall centre of mass frame. The hamiltonian describes three independent simple harmonic oscillators so that the energies solving the (quantized) Schrödinger equation, are $$E=(n_\rho +\frac{3}{2})\omega _\rho +(n_\lambda +\frac{3}{2})\omega _\lambda +(n_\sigma +\frac{3}{2})\omega _\sigma $$ (4) where e.g. $`n_\rho =n_\rho ^x+n_\rho ^y+n_\rho ^z`$, with $`n_\rho ^x,n_\rho ^y`$ and $`n_\rho ^z`$ the degrees of excitation in the three space directions. The vibrational frequencies are $$\omega _{\rho }^{}{}_{}{}^{2}=\omega _{\lambda }^{}{}_{}{}^{2}=\frac{k}{M}\omega _{\sigma }^{}{}_{}{}^{2}=(1+3\frac{M}{m})\frac{k}{M}$$ (5) The wave functions solving the Schrödinger equation can be denoted by $$\psi _{n_\rho }(𝝆)\psi _{n_\lambda }(𝝀)\psi _{n_\sigma }(𝝈)$$ (6) where e.g. $`n_\rho `$ denotes the set $`\{n_\rho ^x,n_\rho ^y,n_\rho ^z\}`$. One can verify that the quark relative velocities $`(\frac{k}{Mm^2})^{\frac{1}{4}}`$, so that they go to zero if $`M\frac{k}{m^2}`$. The criterion for the validity of the adiabatic approximation is hence satisfied for large quark masses. ### 2.2 Adiabatic solution Fix the quark positions so that $`\dot{𝐫}_1=\dot{𝐫}_2=\dot{𝐫}_3=0`$. The potential energy depends only on the differences between positions, and is hence unaffected by fixing the quark positions. The kinetic energy in Eq. 1 is $`\frac{1}{2}m\dot{𝐫}_4^2`$, and using $`2\sqrt{3}\dot{𝐫}_4=(1+3\frac{M}{m})\dot{𝝈}`$ from Eq. 2, the hamiltonian (1) is $$H_j=\frac{1}{2}\frac{m}{12}(1+3\frac{M}{m})^2\dot{𝝈}^2+\frac{1}{2}k(𝝆^2+𝝀^2+\frac{1}{4}(1+3\frac{M}{m})^2𝝈^2)$$ (7) Only the variable $`𝝈`$ is dynamical. The energies (adiabatic potentials) are $$E_j=(n_\sigma +\frac{3}{2})\omega _\sigma ^{^{}}+\frac{1}{2}k(𝝆^2+𝝀^2)\omega _{\sigma }^{^{}}{}_{}{}^{2}=3\frac{k}{m}$$ (8) Now allow the quarks to move in their centre of mass frame, so that the quark motion hamiltonian is $`H_q`$ $`=`$ $`{\displaystyle \frac{1}{2}}M(\dot{𝐫}_1^2+\dot{𝐫}_2^2+\dot{𝐫}_3^2)+E_j`$ (9) $`=`$ $`{\displaystyle \frac{1}{2}}M(\dot{𝝆}^2+\dot{𝝀}^2)+(n_\sigma +{\displaystyle \frac{3}{2}})\omega _\sigma ^{^{}}+{\displaystyle \frac{1}{2}}k(𝝆^2+𝝀^2)`$ The energies are $$E_q=(n_\rho +\frac{3}{2})\omega _\rho +(n_\lambda +\frac{3}{2})\omega _\lambda +(n_\sigma +\frac{3}{2})\omega _\sigma ^{^{}}$$ (10) It is easy to see that these adiabatic approximation energies only agree with the exact energies (4) when $`\frac{M}{m}1`$, i.e. when the quarks are much heavier than the gluons. This is in accord with one’s naïve expectation for the validity of the adiabatic approximation. ### 2.3 Redefined adiabatic solution We follow the same procedure as for the adiabatic solution, with one critical change. The quark positions are still fixed will respect to each other, but all positions are now defined with respect to the overall centre of mass, before we allow the quarks to move. Define the overall centre of mass as $$𝐑=\frac{M(𝐫_1+𝐫_2+𝐫_3)+m𝐫_4}{m+3M}\dot{𝐑}=\frac{m}{m+3M}\dot{𝐫}_4=\frac{1}{2\sqrt{3}}\dot{𝝈}$$ (11) Define new coordinates which are the positions of each particle with respect to the overall centre of mass. The time derivatives of these coordinates are $$\dot{𝐫}_1^{^{}}=\dot{𝐫}_2^{^{}}=\dot{𝐫}_3^{^{}}=\dot{𝐑}=\frac{1}{2\sqrt{3}}\dot{𝝈}\dot{𝐫}_4^{^{}}=\dot{𝐫}_4\dot{𝐑}=\frac{\sqrt{3}}{2}\frac{M}{m}\dot{𝝈}$$ (12) The kinetic energy in terms of the new coordinates is $$\frac{1}{2}M(\dot{𝐫}_1^{{}_{}{}^{}\mathrm{\hspace{0.33em}2}}+\dot{𝐫}_2^{{}_{}{}^{}\mathrm{\hspace{0.33em}2}}+\dot{𝐫}_3^{{}_{}{}^{}\mathrm{\hspace{0.33em}2}})+\frac{1}{2}m\dot{𝐫}_4^{{}_{}{}^{}\mathrm{\hspace{0.33em}2}}=\frac{1}{2}\frac{M}{4}(1+3\frac{M}{m})\dot{𝝈}^2$$ (13) This kinetic energy combined with the potential in Eq. 1 is $$H_j=\frac{1}{2}\frac{M}{4}(1+3\frac{M}{m})\dot{𝝈}^2+\frac{1}{2}k(𝝆^2+𝝀^2+\frac{1}{4}(1+3\frac{M}{m})^2𝝈^2)$$ (14) We used the fact that the potential energy depends only on the differences between positions. The energies (redefined adiabatic potentials) are $$E_j=(n_\sigma +\frac{3}{2})\omega _\sigma +\frac{1}{2}k(𝝆^2+𝝀^2)$$ (15) and the junction wave functions $`\psi _{n_\sigma }(𝝈)`$. Allowing the quarks to move, the quark motion hamiltonian is $$H_q=\frac{1}{2}M(\dot{𝝆}^2+\dot{𝝀}^2)+(n_\sigma +\frac{3}{2})\omega _\sigma +\frac{1}{2}k(𝝆^2+𝝀^2)$$ (16) The energies are $$E_q=(n_\rho +\frac{3}{2})\omega _\rho +(n_\lambda +\frac{3}{2})\omega _\lambda +(n_\sigma +\frac{3}{2})\omega _\sigma $$ (17) which are identical to the exact solution (4). The quark wave functions are $`\psi _{n_\rho }(𝝆)\psi _{n_\lambda }(𝝀)`$. In order to obtain the full wave functions of the system, we take the direct product of the quark and junction wave functions, i.e. $`\psi _{n_\rho }(𝝆)\psi _{n_\lambda }(𝝀)\psi _{n_\sigma }(𝝈)`$. These are identical to the wave functions of the exact solution (6). In conclusion, the solution obtained via the redefined adiabatic approximation is exact. We shall now show how the specific dynamics (1) provide a toy model for the quark model, in the sense that many of the features needed for the validity of quark model, are exact in the toy model. Quark models enable the dynamics of quarks, while freezing out the dynamics of gluons. In the toy model this corresponds to claiming that gluons are in the same wave function for all conventional baryons. But this is manifestly the case in the redefined adiabatic approximation. The redefined adiabatic potential (15) is explicitly dependent on quark mass. This corresponds to quark models: the Coulomb interaction is usually postulated to have a coupling that depends on quark mass. On the other hand, the adiabatic potential (8) is not dependent on quark mass. The wave functions obtained from the adiabatic approximation differ from those in the redefined adiabatic approximation only for the gluons. When one freezes the dynamics of gluons, as in the quark model, one does not notice the difference between these two approximations. Even the functional forms of the gluon wave functions are the same: the wave functions only differ in the dimensionful scale they contain. Thus the forms can be used interchangeably. When a process is studied that involves more than one redefined adiabatic potential, e.g. hybrid baryon decay to a conventional baryon and meson, care has to be taken to use the redefined adiabatic gluon wave functions instead of the adiabatic ones. The former gluon wave functions are dependent on $`\omega _\sigma `$, which is itself dependent on quark mass. ## 3 Where can one search for hybrid baryons? We first outline the results of a recent flux–tube model calculation.$`^\mathrm{?}`$ The flavour, non–relativistic spin $`S`$ and $`J^P`$ of the seven low–lying hybrid baryons are $`(N,\mathrm{\Delta })^{2S+1}J^P=N^2\frac{1}{2}^+,N^2\frac{3}{2}^+,\mathrm{\Delta }^4\frac{1}{2}^+,\mathrm{\Delta }^4\frac{3}{2}^+,\mathrm{\Delta }^4\frac{5}{2}^+`$, where the first two states double. The bag model has the same number of states. The pair $`N^2\frac{1}{2}^+,N^2\frac{3}{2}^+`$ has the same quantum numbers as in the flux-tube model.$`^\mathrm{?}`$ The other five states in the bag model are $`\mathrm{\Delta }^2\frac{1}{2}^+,\mathrm{\Delta }^2\frac{3}{2}^+,N^4\frac{1}{2}^+,N^4\frac{3}{2}^+,N^4\frac{5}{2}^+`$. The state $`N\frac{1}{2}^+`$ was studied in QCD sum rules.$`^\mathrm{?}`$ The hyperfine interaction moves the $`\mathrm{\Delta }`$ states upwards and the $`N`$ states downwards.$`^\mathrm{?}`$ Hence there are four low–lying $`N`$ hybrid baryons with a mass of $`1870\pm 100`$ MeV. This is somewhat higher, but sometimes within errors, of bag models and QCD sum rules which find a mass around 1.5 GeV.$`^{\mathrm{?},\mathrm{?}}`$ The wave function sizes are estimated as $`\sqrt{\rho ^2}=\sqrt{\lambda ^2}=2.1`$ (conventional baryons) and $`2.5`$ (hybrid baryons) GeV<sup>-1</sup>. Hybrid baryons are hence larger than conventional baryons. The following techniques may enable the detection of hybrid baryons: $``$ Overabundance of states: This approach is troublesome, since not even all conventional baryons in the appropriate mass region have been discovered yet. However, hybrid baryon states are likely to be discovered before all conventional baryons in a relevant mass region, and hence need to be studied alongside conventional baryons. $``$ Decays: Except for a QCD sum rule motivated suggestion that hybrid baryons should decay to $`N\sigma `$,$`^\mathrm{?}`$ no decay calculations have been performed. However, decay of hybrid baryons to $`N\rho `$ and $`N\omega `$ is a priori interesting since it isolates states in the correct mass region, without contamination from lower–lying conventional baryons. Study of these decay modes will also yield information on photo– and electroproduction in Hall B and C at Jefferson Lab, since the $`\rho `$ and $`\omega `$ couple to the photon via vector meson dominance. The process $`N\pi `$ hybrid $`N\omega `$ can be studied at the D–line at Crystal Ball E913. $``$ Diffractive $`\gamma N`$ and $`\pi N`$ production: The detection of the hybrid meson candidate $`\pi (1800)`$ in diffractive $`\pi N`$ collisions by VES$`^\mathrm{?}`$ may indicate that hybrid mesons are producted abundantly via meson–pomeron fusion. If this is the case, one expects significant production of hybrid baryons via baryon–pomeron fusion, i.e. production in diffractive $`\gamma N`$ and $`\pi N`$ collisions. $``$ Production in $`\psi `$ decays: Naïve expectations are that the gluon–rich environment of $`\psi `$ decays should lead to dominant production of glueballs, but also signifant production of hybrid mesons and baryons. The large branching ratios$`^\mathrm{?}`$ $`Br(\psi p\overline{p}\omega ,p\overline{p}\eta ^{^{}})10^3`$ may contain hybrid baryons decaying to $`(p,\overline{p})\omega `$ especially. Recently a $`J^P=\frac{1}{2}^+`$ $`2\sigma `$ peak at mass $`1834_{55}^{+46}`$ MeV was seen in $`\mathrm{\Psi }p\overline{p}\eta `$.$`^\mathrm{?}`$ $``$ Electroproduction: In the flux–tube model, i.e. the adiabatic picture of a hybrid baryon, there is the qualitative conclusion that “$`epeX`$ should produce ordinary $`N^{}`$’s and hybrid baryons with comparable cross–sections”.$`^\mathrm{?}`$ However, the conclusions obtained from the three quark – gluon picture of a hybrid baryon is different. For large $`Q^2`$ electroproduction, the $`Q^2`$ dependence of the amplitudes is summarized in Table 1. Since the photon has both a transverse and longitudinal component, the amplitude for a conventional baryon is expected to dominate that of the hybrid baryon as $`Q^2`$ becomes large.$`^\mathrm{?}`$ For small $`Q^2`$ the conclusion agrees with the large $`Q^2`$ result for transverse photons, but is more dramatic for longitudinal photons: the amplitude vanishes.$`^\mathrm{?}`$ It has accordingly been concluded that the (radially excited) conventional baryon is dominantly electroproduced, with the hybrid baryon subdominant to the resonances $`S_{11}(1535),D_{13}(1520)`$ and $`\mathrm{\Delta }`$ as $`Q^2`$ increases.$`^\mathrm{?}`$ The $`Q^2`$ dependence of the electroproduction of a resonance can be measured at Jefferson Lab Hall B and an energy upgraded Jefferson Lab. A hybrid baryon is expected to behave different from nearby conventional baryons as a function of $`Q^2`$. One needs to perform partial wave analysis at different $`Q^2`$. For large $`Q^2`$ cross–sections are small, which would make this way of distinguishing conventional from hybrid baryons challenging. ## Acknowledgments The suggestion to study hybrid baryons with a simple harmonic oscillator potential was made by Nathan Isgur and Simon Capstick. The manuscript was carefully read by Noel Black. This research is supported by the Department of Energy under contract W-7405-ENG-36. ## References
warning/0004/math0004118.html	ar5iv	text	# Painlevé-Calogero Correspondence Revisited ## I Introduction The so called Painlevé equations are the following six equations discovered by Painlevé and Gambier : $`(\mathrm{P}_{\mathrm{VI}})`$ $`{\displaystyle \frac{d^2\lambda }{dt^2}}={\displaystyle \frac{1}{2}}\left({\displaystyle \frac{1}{\lambda }}+{\displaystyle \frac{1}{\lambda 1}}+{\displaystyle \frac{1}{\lambda t}}\right)\left({\displaystyle \frac{d\lambda }{dt}}\right)^2\left({\displaystyle \frac{1}{t}}+{\displaystyle \frac{1}{t1}}+{\displaystyle \frac{1}{\lambda t}}\right){\displaystyle \frac{d\lambda }{dt}}`$ $`+{\displaystyle \frac{\lambda (\lambda 1)(\lambda t)}{t^2(t1)^2}}\left(\alpha +{\displaystyle \frac{\beta t}{\lambda ^2}}+{\displaystyle \frac{\gamma (t1)}{(\lambda 1)^2}}+{\displaystyle \frac{\delta t(t1)}{(\lambda t)^2}}\right).`$ $`(\mathrm{P}_\mathrm{V})`$ $`{\displaystyle \frac{d^2\lambda }{dt^2}}=\left({\displaystyle \frac{1}{2\lambda }}+{\displaystyle \frac{1}{\lambda 1}}\right)\left({\displaystyle \frac{d\lambda }{dt}}\right)^2{\displaystyle \frac{1}{t}}{\displaystyle \frac{d\lambda }{dt}}`$ $`+{\displaystyle \frac{\lambda (\lambda 1)^2}{t^2}}\left(\alpha +{\displaystyle \frac{\beta }{\lambda ^2}}+{\displaystyle \frac{\gamma t}{(\lambda 1)^2}}+{\displaystyle \frac{\delta t^2(\lambda +1)}{(\lambda 1)^3}}\right).`$ $`(\mathrm{P}_{\mathrm{IV}})`$ $`{\displaystyle \frac{d^2\lambda }{dt^2}}={\displaystyle \frac{1}{2\lambda }}\left({\displaystyle \frac{d\lambda }{dt}}\right)^2+{\displaystyle \frac{3}{2}}\lambda ^3+4t\lambda ^2+2(t^2\alpha )\lambda +{\displaystyle \frac{\beta }{\lambda }}.`$ $`(\mathrm{P}_{\mathrm{III}})`$ $`{\displaystyle \frac{d^2\lambda }{dt^2}}={\displaystyle \frac{1}{\lambda }}\left({\displaystyle \frac{d\lambda }{dt}}\right)^2{\displaystyle \frac{1}{t}}{\displaystyle \frac{d\lambda }{dt}}+{\displaystyle \frac{\lambda ^2}{4t^2}}\left(\alpha +{\displaystyle \frac{\beta t}{\lambda ^2}}+\gamma \lambda +{\displaystyle \frac{\delta t^2}{4\lambda ^3}}\right).`$ $`(\mathrm{P}_{\mathrm{II}})`$ $`{\displaystyle \frac{d^2\lambda }{dt^2}}=2\lambda ^3+t\lambda +\alpha .`$ $`(\mathrm{P}_\mathrm{I})`$ $`{\displaystyle \frac{d^2\lambda }{dt^2}}=6\lambda ^2+t.`$ The third equation $`\mathrm{P}_{\mathrm{III}}`$ is slightly modified; the original equation can be reproduced by the simple change of variables $`(t,\lambda )(t^2,t\lambda )`$. It is well known that these equations are characterized by the absence of “movable singularities” other than poles. R. Fuchs proposed two more approaches to the sixth equation $`\mathrm{P}_{\mathrm{VI}}`$. One approach is the concept of isomonodromic deformations. In this approach, $`\mathrm{P}_{\mathrm{VI}}`$ is interpreted as a differential equation describing isomonodromic deformations of a linear ordinary differential equation on the Riemann sphere. This is the origin of many subsequent researches. Another approach relates $`\mathrm{P}_{\mathrm{VI}}`$ to an incomplete elliptic integral. Painlevé took the second approach, and derived a new expression of $`\mathrm{P}_{\mathrm{VI}}`$ in term of the Weierstrass $`\mathrm{}`$-function. This work of Painlevé is briefly reviewed in Okamoto’s work on affine Weyl group symmetries of $`\mathrm{P}_{\mathrm{VI}}`$ . Manin revived the almost forgotten work of Fuchs and Painlevé after nearly ninety years. Manin’s remarkable idea is to use the elliptic modulus $`\tau `$, rather than $`t`$, as an independent variable. The outcome is a Hamiltonian system with a Hamiltonian of the normal form $`=p^2/2+V(q)`$, where the potential is a linear combination of the Weierstrass $`\mathrm{}`$-function and its shift by three three half periods. This is a non-autonomous system, because the Hamiltonian depends on the “time” $`\tau `$ through the $`\tau `$-dependence of the $`\mathrm{}`$-function. Levin and Olshanetsky pointed out that Manin’s equation resembles the so called Calogero-Moser systems, i.e., the various extensions of the integrable many-body systems first discovered by Calogero . More precisely, the Hamiltonian $``$ is identical to a special case (the rank-one elliptic model) of Inozemtsev’s extensions of the Calogero-Moser systems. Levin and Olshanetsky called this relation the “Painlevé-Calogero correspondence”. One will naturally ask if this correspondence can be extended to the other Painlevé equations. Manin himself raised this problem in his paper. Olshanetsky conjectured that a degenerate version of Inozemtsev’s elliptic model will emerge therein. This paper aims to answer this question affirmatively. A guiding principle is the degeneration relation of the six Painlevé equations . This relation can be schematically expressed as follows: $`\begin{array}{ccccccc}\mathrm{P}_{\mathrm{VI}}& & \mathrm{P}_\mathrm{V}& & \mathrm{P}_{\mathrm{IV}}& & \\ & & & & & & \\ & & \mathrm{P}_{\mathrm{III}}& & \mathrm{P}_{\mathrm{II}}& & \mathrm{P}_\mathrm{I}\end{array}`$ (4) This diagram means, for instance, that $`\mathrm{P}_\mathrm{V}`$ can be derived from $`\mathrm{P}_{\mathrm{VI}}`$ by a degeneration process, which amounts to confluence of singular points of the aforementioned linear ordinary differential equation in the isomonodromic approach. We shall trace this process carefully on the “Calogero side”, and find a $`\mathrm{P}_\mathrm{V}`$-version of Manin’s equation. In principle, one can thus find an analogue of Manin’s equation for all the six Painlevé equations (though, actually, one can resort to a more direct approach that bypasses the complicated degeneration process). Remarkably (or rather naturally?), all the six equations on the Calogero side turn out to become a (non-autonomous) Hamiltonian system with a Hamiltonian of the normal form $`=p^2/2+V(q)`$. Furthermore, the Hamiltonians on the Calogero side of $`\mathrm{P}_\mathrm{V}`$ and $`\mathrm{P}_{\mathrm{IV}}`$ coincide with the Hamiltonians of the (rank one) hyperbolic and rational models in Inozemtsev’s classification (which were discovered by Levi and Wojciechowski before Inozemtsev’s work). Those corresponding to the other three Painlevé equations are not included therein, but may be thought of as a further degeneration of the hyperbolic and rational models. One can further proceed to the higher rank models, and ask if there is still a Painlevé-Calogero correspondence. We shall show that this is also the case. The Painlevé side of the correspondence is a kind of multi-dimensional extensions of the Painlevé equations. They are obviously different from another multi-dimensional extension called the “Garnier systems” . For this reason, we call our multi-dimensional extension a multi-component version of the Painlevé equations. This paper is organized as follows. Section 2 is a brief review of the work of Fuchs, Painlevé and Manin. Section 3 deals with $`\mathrm{P}_\mathrm{V}`$, $`\mathrm{P}_{\mathrm{IV}}`$ and $`\mathrm{P}_{\mathrm{III}}`$. The degeneration process is discussed in detail for the case of $`\mathrm{P}_\mathrm{V}`$. The direct approach is illustrated for the case of $`\mathrm{P}_{\mathrm{IV}}`$ and $`\mathrm{P}_{\mathrm{III}}`$. Section 4 shows a reformulation of the foregoing calculations in the Hamiltonian formalism. The status of $`\mathrm{P}_{\mathrm{II}}`$ and $`\mathrm{P}_\mathrm{I}`$ is also clarified therein. Section 5 is devoted to the higher rank Inozemtsev Hamiltonians and the multi-component Painlevé equations. Section 6 is for concluding remarks. Part of technical details are gathered in Appendices. ## II Painlevé-Calogero Correspondence for $`\mathrm{P}_{\mathrm{VI}}`$ We here briefly review the work of Fuchs, Painlevé and Manin. Fuchs rewrites $`\mathrm{P}_{\mathrm{VI}}`$ into the following form: $`t(1t)_t{\displaystyle _{\mathrm{}}^\lambda }{\displaystyle \frac{dz}{\sqrt{z(z1)(zt)}}}`$ (5) $`=`$ $`\sqrt{\lambda (\lambda 1)(\lambda t)}\left[\alpha +{\displaystyle \frac{\beta t}{\lambda ^2}}+{\displaystyle \frac{\gamma (t1)}{(\lambda 1)^2}}+\left(\delta {\displaystyle \frac{1}{2}}\right){\displaystyle \frac{t(t1)}{(\lambda t)^2}}\right].`$ Here $`_t`$ is the linear differential operator (Picard-Fuchs operator) $`_t=t(1t){\displaystyle \frac{d^2}{dt^2}}+(12t){\displaystyle \frac{d}{dt}}{\displaystyle \frac{1}{4}},`$ (6) which also appears in the Picard-Fuchs equation of complete elliptic integrals. In this respect, $`\mathrm{P}_{\mathrm{VI}}`$ may be thought of as an inhomogeneous (and nonlinear) analogue of the Picard-Fuchs equation. Painlevé and Manin make use of a parametrization of the elliptic curve $`y^2=z(z1)(zt)`$ (7) by the Weierstrass $`\mathrm{}`$-function. Let $`\mathrm{}(u)`$ be the $`\mathrm{}`$-function with primitive periods $`1`$ and $`\tau `$: $`\mathrm{}(u)=\mathrm{}(u1,\tau )={\displaystyle \frac{1}{u^2}}+{\displaystyle \underset{(m,n)(0,0)}{}}\left({\displaystyle \frac{1}{(u+m+n\tau )^2}}{\displaystyle \frac{1}{(m+n\tau )^2}}\right),`$ (8) The parametrization is now given by $`z={\displaystyle \frac{\mathrm{}(u)e_1}{e_2e_1}},y={\displaystyle \frac{\mathrm{}^{}(u)}{2(e_2e_1)^{3/2}}},`$ (9) where $`e_n=\mathrm{}(\omega _n)`$, $`n=1,2,3`$ are the values of $`\mathrm{}(u)`$ at the three half period points $`\omega _1=1/2`$, $`\omega _2=(1+\tau )/2`$, $`\omega _3=\tau /2`$. Manin’s excellent idea is to do a simultaneous change of the dependent variable $`\lambda q`$ by $`\lambda ={\displaystyle \frac{\mathrm{}(q)e_1}{e_2e_1}},`$ (10) and the independent vrariable $`t\tau `$ by $`t={\displaystyle \frac{e_3e_1}{e_2e_1}}.`$ (11) Manin presents the beautiful formula $`{\displaystyle \frac{d\tau }{dt}}={\displaystyle \frac{\pi i}{t(t1)(e_2e_1)}},`$ (12) for the Jacobian of the latter, which plays a key role in his calculations. $`\mathrm{P}_{\mathrm{VI}}`$ is thereby mapped to the equation $`(2\pi i)^2{\displaystyle \frac{d^2q}{d\tau ^2}}={\displaystyle \underset{n=0}{\overset{3}{}}}\alpha _n\mathrm{}^{}(q+\omega _n),`$ (13) where the parameters on the right hand side are connected with the parameters of $`\mathrm{P}_{\mathrm{VI}}`$ as $`\alpha _0=\alpha `$, $`\alpha _1=\beta `$, $`\alpha _2=\gamma `$, $`\alpha _3=\delta +1/2`$. This equation is equivalent to the Hamiltonian system $`2\pi i{\displaystyle \frac{dq}{d\tau }}={\displaystyle \frac{}{p}},2\pi i{\displaystyle \frac{dp}{d\tau }}={\displaystyle \frac{}{q}}`$ (14) with the Hamiltonian $`={\displaystyle \frac{p^2}{2}}{\displaystyle \underset{n=0}{\overset{3}{}}}\alpha _n\mathrm{}(q+\omega _n).`$ (15) ## III Correspondence for $`\mathrm{P}_\mathrm{V}`$, $`\mathrm{P}_{\mathrm{IV}}`$ and $`\mathrm{P}_{\mathrm{III}}`$ ### III.1 Degeneration of $`\mathrm{P}_{\mathrm{VI}}`$ to $`\mathrm{P}_\mathrm{V}`$ The degeneration of $`\mathrm{P}_{\mathrm{VI}}`$ to $`\mathrm{P}_\mathrm{V}`$ is achieved by rescaling the time variable and the parameters as $`t=1+ϵ\stackrel{~}{t},\alpha =\stackrel{~}{\alpha },\beta =\stackrel{~}{\beta },\gamma ={\displaystyle \frac{\stackrel{~}{\gamma }}{ϵ}}{\displaystyle \frac{\stackrel{~}{\delta }}{ϵ^2}},\delta ={\displaystyle \frac{\stackrel{~}{\delta }}{ϵ^2}}`$ (16) and letting $`ϵ0`$ while leaving $`\stackrel{~}{\alpha },\mathrm{},\stackrel{~}{\gamma }`$ and $`\stackrel{~}{t}`$ finite . The building blocks of Fuchs’ equation (5) turn out to survive this scaling limit as follows: 1. The Picard-Fuchs operator: $`t(1t)_t\stackrel{~}{t}^2{\displaystyle \frac{d^2}{d\stackrel{~}{t}^2}}+\stackrel{~}{t}{\displaystyle \frac{d}{d\stackrel{~}{t}}}=\left(\stackrel{~}{t}{\displaystyle \frac{d}{d\stackrel{~}{t}}}\right)^2.`$ 2. The sum $`\alpha +\mathrm{}`$ of four terms on the right hand side: $`\alpha +{\displaystyle \frac{\beta t}{\lambda ^2}}+{\displaystyle \frac{\gamma (t1)}{(\lambda 1)^2}}+\left(\delta {\displaystyle \frac{1}{2}}\right){\displaystyle \frac{t(t1)}{(\lambda t)^2}}\stackrel{~}{\alpha }+{\displaystyle \frac{\stackrel{~}{\beta }}{\lambda ^2}}+{\displaystyle \frac{\stackrel{~}{\gamma }\stackrel{~}{t}}{(\lambda 1)^2}}+{\displaystyle \frac{\stackrel{~}{\delta }\stackrel{~}{t}^2(\lambda +1)}{(\lambda 1)^3}}.`$ 3. The quare root on the right hand side: $`\sqrt{\lambda (\lambda 1)(\lambda t)}\sqrt{\lambda }(\lambda 1).`$ 4. The incomplete elliptic integral: $`{\displaystyle _{\mathrm{}}^\lambda }{\displaystyle \frac{dz}{\sqrt{z(z1)(zt)}}}{\displaystyle _{\mathrm{}}^\lambda }{\displaystyle \frac{dz}{\sqrt{z}(z1)}}.`$ In particular, the degeneration of $`\mathrm{P}_{\mathrm{VI}}`$ to $`\mathrm{P}_\mathrm{V}`$ is associated with the degeneration of the elliptic curve to a rational curve, $`y^2=z(z1)(zt)y^2=z(z1)^2,`$ (17) or, equivalently, the degeneration of the torus $`/(+\tau )`$ to the cylinder $`/`$. Thus, rewriting $`\stackrel{~}{\alpha },\stackrel{~}{\beta },\stackrel{~}{\gamma },\stackrel{~}{\delta }`$ and $`\stackrel{~}{t}`$ to $`\alpha ,\beta ,\gamma ,\delta `$ and $`t`$, we obtain the following equation as a $`\mathrm{P}_\mathrm{V}`$-version of Fuchs’ equation: $`\left(t{\displaystyle \frac{d}{dt}}\right)^2{\displaystyle _{\mathrm{}}^\lambda }{\displaystyle \frac{dz}{\sqrt{z}(z1)}}=\sqrt{\lambda }(\lambda 1)\left(\alpha +{\displaystyle \frac{\beta }{\lambda ^2}}+{\displaystyle \frac{\gamma t}{(\lambda 1)^2}}+{\displaystyle \frac{\delta t^2(\lambda +1)}{(\lambda 1)^3}}\right).`$ (18) ### III.2 Analogue of Manin’s equation for $`\mathrm{P}_\mathrm{V}`$ As an counterpart of the $`q`$-variable for $`\mathrm{P}_{\mathrm{VI}}`$, we now consider $`q={\displaystyle _{\mathrm{}}^\lambda }{\displaystyle \frac{dz}{\sqrt{z}(z1)}}.`$ (19) If one prefers to being more faithful to Manin’s parametrization, one should rather define $`q`$ as $`q={\displaystyle \frac{1}{2\pi i}}{\displaystyle _{\mathrm{}}^\lambda }{\displaystyle \frac{dz}{\sqrt{z}(z1)}},`$ because $`2(e_2e_1)^{1/2}2\pi i`$ as $`\mathrm{Im}\tau +\mathrm{}`$ (see Appendix B). Since there is no substantial difference, let us take the first definition that is slightly simpler for calculations. Let us rewrite (18) in terms of $`q`$. The integral can be readily calculated as $`q=\mathrm{log}\left({\displaystyle \frac{\sqrt{\lambda }1}{\sqrt{\lambda }+1}}\right),`$ (20) so that the inverse relation can be written $`\sqrt{\lambda }=\mathrm{coth}(q/2).`$ (21) Terms on the right hand side of (18) can be calculated as follows: $`\sqrt{\lambda }(\lambda 1)`$ $`=`$ $`{\displaystyle \frac{\mathrm{cosh}(q/2)}{\mathrm{sinh}^3(q/2)}},`$ $`\sqrt{\lambda }(\lambda 1){\displaystyle \frac{1}{\lambda ^2}}`$ $`=`$ $`{\displaystyle \frac{\mathrm{sinh}(q/2)}{\mathrm{cosh}^3(q/2)}},`$ $`\sqrt{\lambda }(\lambda 1){\displaystyle \frac{1}{(\lambda 1)^2}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{sinh}(q),`$ $`\sqrt{\lambda }(\lambda 1){\displaystyle \frac{(\lambda +1)}{(\lambda 1)^3}}`$ $`=`$ $`{\displaystyle \frac{\lambda ^{3/2}+\lambda ^{1/2}}{(\lambda 1)^2}}={\displaystyle \frac{1}{4}}\mathrm{sinh}(2q).`$ The differential equation for $`q`$ eventually takes the form $`\left(t{\displaystyle \frac{d}{dt}}\right)^2q={\displaystyle \frac{V(q)}{q}},`$ (22) where $`V(q)={\displaystyle \frac{\alpha }{\mathrm{sinh}^2(q/2)}}{\displaystyle \frac{\beta }{\mathrm{cosh}^2(q/2)}}+{\displaystyle \frac{\gamma t}{2}}\mathrm{cosh}(q)+{\displaystyle \frac{\delta t^2}{8}}\mathrm{cosh}(2q).`$ (23) This gives a $`\mathrm{P}_\mathrm{V}`$-version of Manin’s equation. Note that this equation can be readily converted to a Hamiltonian system with the Hamiltonian $`=p^2/2+V(q)`$. Remark. A very similar change of dependent variable for $`\mathrm{P}_\mathrm{V}`$ is discussed in the book of Iwasaki et al. . ### III.3 Idea of direct approach Although the degeneration process can be continued to the other Painlevé equations, we now present a more direct approach. Note that the integrand is connected with the coefficient of $`(d\lambda /dt)^2`$ in the original Painlevé equation by the following very simple relation: $`{\displaystyle \frac{1}{\sqrt{z(z1)(zt)}}}`$ $`=`$ $`\mathrm{exp}\left[{\displaystyle \frac{1}{2}\left(\frac{1}{z}+\frac{1}{z1}+\frac{1}{zt}\right)𝑑z}\right],`$ $`{\displaystyle \frac{1}{\sqrt{z}(z1)}}`$ $`=`$ $`\mathrm{exp}\left[{\displaystyle \left(\frac{1}{2z}+\frac{1}{z1}\right)𝑑z}\right].`$ If this is a correct prescription, one will be able to define the $`q`$-variable for $`\mathrm{P}_{\mathrm{III}}`$ and $`\mathrm{P}_{\mathrm{II}}`$ directly without the cumbersome degeneration process. This is indeed the case, as we shall show below. ### III.4 $`q`$-variable for $`\mathrm{P}_{\mathrm{IV}}`$ Since the expected integrand is given by $`\mathrm{exp}\left({\displaystyle \frac{dz}{2z}}\right)={\displaystyle \frac{1}{\sqrt{z}}},`$ (24) we define $`q={\displaystyle ^\lambda }{\displaystyle \frac{dz}{\sqrt{z}}}=2\sqrt{\lambda }.`$ (25) This can be solved for $`\lambda `$ as $`\lambda =\left({\displaystyle \frac{q}{2}}\right)^2.`$ (26) Honest calculations show that all derivative terms of $`\mathrm{P}_{\mathrm{IV}}`$ can be absorbed by the second derivative of $`q`$: $`{\displaystyle \frac{d^2q}{dt^2}}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{\lambda }}}{\displaystyle \frac{d^2\lambda }{dt^2}}{\displaystyle \frac{1}{2\lambda \sqrt{\lambda }}}\left({\displaystyle \frac{d\lambda }{dt}}\right)^2`$ (27) $`=`$ $`{\displaystyle \frac{1}{\sqrt{\lambda }}}\left({\displaystyle \frac{3}{2}}\lambda ^3+4t\lambda ^2+2(t^2\alpha )\lambda +{\displaystyle \frac{\beta }{\lambda }}\right).`$ Substituting $`\lambda =(q/2)^2`$ gives the second order differential equation $`{\displaystyle \frac{d^2q}{dt^2}}={\displaystyle \frac{V(q)}{q}}`$ (28) with the potential $`V(q)={\displaystyle \frac{1}{2}}\left({\displaystyle \frac{q}{2}}\right)^62t\left({\displaystyle \frac{q}{2}}\right)^42(t^2\alpha )\left({\displaystyle \frac{q}{2}}\right)^2+\beta \left({\displaystyle \frac{q}{2}}\right)^2.`$ (29) ### III.5 $`q`$-variable for $`\mathrm{P}_{\mathrm{III}}`$ The integrand is expected to be given by $`\mathrm{exp}\left({\displaystyle \frac{dz}{z}}\right)={\displaystyle \frac{1}{z}}.`$ (30) We consider $`q={\displaystyle ^\lambda }{\displaystyle \frac{dz}{z}}=\mathrm{log}\lambda `$ (31) and its inversion $`\lambda =e^q.`$ (32) All derivatives terms of $`\mathrm{P}_{\mathrm{III}}`$ are now absorbed by the second derivative of $`q`$ with respect to $`\mathrm{log}t`$: $`\left(t{\displaystyle \frac{d}{dt}}\right)^2q`$ $`=`$ $`{\displaystyle \frac{t^2}{\lambda }}{\displaystyle \frac{d^2\lambda }{dt^2}}+{\displaystyle \frac{t}{\lambda }}{\displaystyle \frac{d\lambda }{dt}}{\displaystyle \frac{t^2}{\lambda ^2}}\left({\displaystyle \frac{d\lambda }{dt}}\right)^2`$ (33) $`=`$ $`{\displaystyle \frac{\alpha \lambda }{4}}+{\displaystyle \frac{\beta t}{4\lambda }}+{\displaystyle \frac{\gamma \lambda ^2}{4}}+{\displaystyle \frac{\delta t^2}{4\lambda ^2}}.`$ Substituting $`\lambda =e^q`$ gives the second order equation $`\left(t{\displaystyle \frac{d}{dt}}\right)^2q={\displaystyle \frac{V(q)}{q}}`$ (34) with the potential $`V(q)={\displaystyle \frac{\alpha }{4}}e^q+{\displaystyle \frac{\beta t}{4}}e^q{\displaystyle \frac{\gamma }{8}}e^{2q}+{\displaystyle \frac{\delta t^2}{8}}e^{2q}.`$ (35) ### III.6 Summary Let us summarize the results of this section. ###### Theorem 1 The foregoing change of variable $`\lambda q`$ maps $`\mathrm{P}_\mathrm{V}`$, $`\mathrm{P}_{\mathrm{IV}}`$ and $`\mathrm{P}_{\mathrm{III}}`$ to a second order differential equation for the new dependent variable $`q`$. These equations are equivalent to a non-autonomous Hamiltonian system with a Hamiltonian of the normal form $`=p^2/2+V(q)`$: ($`\mathrm{P}_\mathrm{V}`$) The Hamiltonian system takes the form $`t{\displaystyle \frac{dq}{dt}}={\displaystyle \frac{}{p}},t{\displaystyle \frac{dp}{dt}}={\displaystyle \frac{}{q}}`$ (36) with the Hamiltonian $`={\displaystyle \frac{p^2}{2}}{\displaystyle \frac{\alpha }{\mathrm{sinh}^2(q/2)}}{\displaystyle \frac{\beta }{\mathrm{cosh}^2(q/2)}}+{\displaystyle \frac{\gamma t}{2}}\mathrm{cosh}(q)+{\displaystyle \frac{\delta t^2}{8}}\mathrm{cosh}(2q).`$ (37) ($`\mathrm{P}_{\mathrm{IV}}`$) The Hamiltonian system takes the form $`{\displaystyle \frac{dq}{dt}}={\displaystyle \frac{}{p}},{\displaystyle \frac{dp}{dt}}={\displaystyle \frac{}{q}}`$ (38) with the Hamiltonian $`={\displaystyle \frac{p^2}{2}}{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{q}{2}}\right)^62t\left({\displaystyle \frac{q}{2}}\right)^42(t^2\alpha )\left({\displaystyle \frac{q}{2}}\right)^2+\beta \left({\displaystyle \frac{q}{2}}\right)^2.`$ (39) ($`\mathrm{P}_{\mathrm{III}}`$) The Hamiltonian system takes the form $`t{\displaystyle \frac{dq}{dt}}={\displaystyle \frac{}{p}},t{\displaystyle \frac{dp}{dt}}={\displaystyle \frac{}{q}}`$ (40) with the Hamiltonian $`={\displaystyle \frac{p^2}{2}}{\displaystyle \frac{\alpha }{4}}e^q+{\displaystyle \frac{\beta t}{4}}e^q{\displaystyle \frac{\gamma }{8}}e^{2q}+{\displaystyle \frac{\delta t^2}{8}}e^{2q}.`$ (41) Remark. 1. The Hamiltonians for $`\mathrm{P}_\mathrm{V}`$ and $`\mathrm{P}_{\mathrm{IV}}`$ coincide with those of the hyperbolic and rational models of Inozemtsev , Levi and Wojciechowski . The Hamiltonian for $`\mathrm{P}_{\mathrm{III}}`$ has no counterpart in their work, but nowadays can be found in the literature . 2. The foregoing construction of the $`q`$-variable does not literally work for $`\mathrm{P}_{\mathrm{II}}`$ and $`\mathrm{P}_\mathrm{I}`$, because there is no $`(d\lambda /dt)^2`$ term. The status of these equations will be clarified in the next section from a different point of view. ## IV Hamiltonian formalism of correspondence ### IV.1 Hamiltonians of Painlevé equations All the six Painlevé equations are known to be expressed in the Hamiltonian form $`{\displaystyle \frac{d\lambda }{dt}}={\displaystyle \frac{H}{\mu }},{\displaystyle \frac{d\lambda }{dt}}={\displaystyle \frac{H}{\lambda }}`$ with a suitable choice of the canonical conjugate variable $`\mu `$ and the Hamiltonian $`H`$ . This expression is by no means unique; we here consider the following Hamiltonians . These Hamiltonians are referred to as the “polynomial Hamiltonians” because they are polynomials in $`\lambda `$ and $`\mu `$: $`(\mathrm{P}_{\mathrm{VI}})`$ $`H={\displaystyle \frac{\lambda (\lambda 1)(\lambda t)}{t(t1)}}\left[\mu ^2\left({\displaystyle \frac{\kappa _0}{\lambda }}+{\displaystyle \frac{\kappa _1}{\lambda 1}}+{\displaystyle \frac{\theta 1}{\lambda t}}\right)\mu +{\displaystyle \frac{\kappa }{\lambda (\lambda 1)}}\right].`$ $`(\mathrm{P}_\mathrm{V})`$ $`H={\displaystyle \frac{\lambda (\lambda 1)^2}{t}}\left[\mu ^2\left({\displaystyle \frac{\kappa _0}{\lambda }}+{\displaystyle \frac{\theta _1}{\lambda 1}}{\displaystyle \frac{\eta _1t}{(\lambda 1)^2}}\right)\mu +{\displaystyle \frac{\kappa }{\lambda (\lambda 1)}}\right].`$ $`(\mathrm{P}_{\mathrm{IV}})`$ $`H=2\lambda \left[\mu ^2\left({\displaystyle \frac{\lambda }{2}}+t+{\displaystyle \frac{\kappa _0}{\lambda }}\right)\mu +{\displaystyle \frac{\theta _{\mathrm{}}}{2}}\right].`$ $`(\mathrm{P}_{\mathrm{III}})`$ $`H={\displaystyle \frac{\lambda ^2}{t}}\left[\mu ^2\left(\eta _{\mathrm{}}+{\displaystyle \frac{\theta _0}{\lambda }}{\displaystyle \frac{\eta _0t}{\lambda ^2}}\right)\mu +{\displaystyle \frac{\eta _{\mathrm{}}(\theta _0+\theta _{\mathrm{}})}{2\lambda }}\right].`$ $`(\mathrm{P}_{\mathrm{II}})`$ $`H={\displaystyle \frac{\mu ^2}{2}}\left(\lambda ^2+{\displaystyle \frac{t}{2}}\right)\mu \left(\alpha +{\displaystyle \frac{1}{2}}\right)\lambda .`$ $`(\mathrm{P}_\mathrm{I})`$ $`H={\displaystyle \frac{\mu ^2}{2}}2\lambda ^3t\lambda .`$ Here $`\kappa _0,\kappa _1,\theta `$, etc. are constants that are connected with the parameters $`\alpha ,\beta ,\gamma ,\delta `$ of the Painlevé equations by simple algebraic relations: $`(\mathrm{P}_{\mathrm{VI}})`$ $`\alpha ={\displaystyle \frac{(\kappa _0+\kappa _1+\theta 1)^2}{2}}2\kappa ,\beta ={\displaystyle \frac{\kappa _0^2}{2}},\gamma ={\displaystyle \frac{\kappa _1^2}{2}},\delta ={\displaystyle \frac{1\theta ^2}{2}},`$ $`(\mathrm{P}_\mathrm{V})`$ $`\alpha ={\displaystyle \frac{(\kappa _0+\theta _1)^2}{2}}2\kappa ,\beta ={\displaystyle \frac{\kappa _0^2}{2}},\gamma =\eta _1(\theta _1+1),\delta ={\displaystyle \frac{\eta _1^2}{2}}.`$ $`(\mathrm{P}_{\mathrm{IV}})`$ $`\alpha =2\theta _{\mathrm{}}\kappa _0+1,\beta =2\kappa _0^2.`$ $`(\mathrm{P}_{\mathrm{III}})`$ $`\alpha =4\eta _{\mathrm{}}\theta _{\mathrm{}},\beta =4\eta _0(\theta _0+1),\gamma =4\eta _{\mathrm{}}^2,\delta =4\eta _0^2.`$ ### IV.2 How to find canonical transformations The goal of this section is to show that the Painlevé-Calogero correspondence is, in fact, a (time-dependent) canonical transformation of two Hamiltonian systems. By this, we mean that the functional relation between $`\lambda `$ and $`q`$ can be extended to $`(\lambda ,\mu )`$ and $`(q,p)`$ so as to satisfy the equation $`\mu d\lambda Hdt=\mathrm{constant}(pdqdT)+\text{exact form}.`$ (42) with a suitably redefined time variable $`T`$ (such as the logarithmic time $`\mathrm{log}t`$ in $`\mathrm{P}_\mathrm{V}`$ and $`\mathrm{P}_{\mathrm{III}}`$). The constant factor on the right hand side is inserted simply for convenience; if necessary, one can normalize the constant to $`1`$ by suitably rescaling $`p,q,`$ and $`T`$. For this reason, wel call this type of coordinate transformation a “canonical” transformation even if the constant factor is not equal to $`1`$. Let us illustrate, in the case of $`\mathrm{P}_{\mathrm{VI}}`$, how to find such a canonical transformation. Suppose that $`\lambda `$ and $`\mu `$ be a solution of $`\mathrm{P}_{\mathrm{VI}}`$ in the aforementioned Hamiltonian formalism, and that $`q`$ be a corresponding solution of Manin’s equation. The canonical equation for $`\lambda `$ takes the form $`{\displaystyle \frac{d\lambda }{dt}}={\displaystyle \frac{\lambda (\lambda 1)(\lambda t)}{t(t1)}}\left(2\mu {\displaystyle \frac{\kappa _0}{\lambda }}{\displaystyle \frac{\kappa _1}{\lambda 1}}{\displaystyle \frac{\theta 1}{\lambda t}}\right).`$ This equations can be solve for $`\mu `$: $`\mu ={\displaystyle \frac{t(t1)}{2\lambda (\lambda 1)(\lambda t)}}{\displaystyle \frac{d\lambda }{dt}}+{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{\kappa _0}{\lambda }}+{\displaystyle \frac{\kappa _1}{\lambda 1}}+{\displaystyle \frac{\theta 1}{\lambda t}}\right).`$ Our task is to rewrite the right hand side in terms of $`p`$ and $`q`$. We first consider $`d\lambda /dt`$. Differentiating (10) against $`t`$ gives $`{\displaystyle \frac{d\lambda }{dt}}=\left({\displaystyle \frac{\mathrm{}^{}(q)}{e_2e_1}}{\displaystyle \frac{dq}{d\tau }}+f_\tau (q)\right){\displaystyle \frac{d\tau }{dt}},`$ where we have introduced the functions $`f(u)={\displaystyle \frac{\mathrm{}(u)e_1}{e_2e_1}},f_\tau (u)={\displaystyle \frac{f(u)}{\tau }}.`$ (43) The derivative $`dq/d\tau `$ can be read off from the canonical equation for $`q`$: $`{\displaystyle \frac{dq}{d\tau }}={\displaystyle \frac{1}{2\pi i}}{\displaystyle \frac{}{p}}={\displaystyle \frac{p}{2\pi i}}.`$ As for the Jacobian $`d\tau /dt`$, Manin’s formula (12) is available. One can thus express $`d\lambda /dt`$ as a function of $`p,q`$ and $`\tau `$. The other part of the foregoing expression of $`\mu `$ contains $`\lambda `$ only, which can be readily converted to a function of $`q`$ and $`\tau `$ by (10). We thus obtain the following expression of $`\mu `$: $`\mu `$ $`=`$ $`{\displaystyle \frac{e_2e_1}{\mathrm{}^{}(q)}}p+{\displaystyle \frac{2\pi i(e_2e_1)^2}{\mathrm{}^{}(q)^2}}f_\tau (q)`$ (44) $`+{\displaystyle \frac{e_2e_1}{2}}\left({\displaystyle \frac{\kappa _0}{\mathrm{}(q)e_1}}+{\displaystyle \frac{\kappa _1}{\mathrm{}(q)e_2}}+{\displaystyle \frac{\theta 1}{\mathrm{}(q)e_3}}\right).`$ We now move the point of view, and think of (10) and (44) as defining a coordinate transformation $`(\lambda ,\mu )(q,p)`$. This gives a canonical transformation that we have sought for: ###### Theorem 2 (10) and (44) define a canonical transformation that connects the Hamiltonian form of $`\mathrm{P}_{\mathrm{VI}}`$ and Manin’s Hamiltonian system. The canonical coordinates and the Hamiltonians of the two systems obey the equation $`\mu d\lambda Hdt=pdq{\displaystyle \frac{d\tau }{2\pi i}}+\text{exact form}.`$ (45) ### IV.3 Proof of Theorem 2 Total differential of (10) gives $`d\lambda ={\displaystyle \frac{\mathrm{}^{}(q)}{e_2e_1}}dq+f_\tau (q)d\tau ,`$ so that $`\mu d\lambda `$ can be expressed as $`\mu d\lambda `$ $`=`$ $`\left({\displaystyle \frac{e_2e_1}{\mathrm{}^{}(q)}}p+{\displaystyle \frac{2\pi i(e_2e_1)^2}{\mathrm{}^{}(q)^2}}f_\tau (q)\right)\left({\displaystyle \frac{\mathrm{}^{}(q)}{e_2e_1}}dq+f_\tau (q)d\tau \right)`$ $`+{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{\kappa _0}{\lambda }}+{\displaystyle \frac{\kappa _1}{\lambda 1}}+{\displaystyle \frac{\theta 1}{\lambda t}}\right)d\lambda `$ $`=`$ $`pdq+(\mathrm{A})+(\mathrm{B})+(\mathrm{C}),`$ where $`(\mathrm{A})`$ $`=`$ $`{\displaystyle \frac{2\pi i(e_2e_1)}{\mathrm{}^{}(q)}}f_\tau (q)dq,`$ $`(\mathrm{B})`$ $`=`$ $`\left({\displaystyle \frac{e_2e_1}{\mathrm{}^{}(q)}}p+{\displaystyle \frac{2\pi i(e_2e_1)^2}{\mathrm{}^{}(q)^2}}f_\tau (q)\right)f_\tau (q)d\tau ,`$ $`(\mathrm{C})`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{\kappa _0}{\lambda }}+{\displaystyle \frac{\kappa _1}{\lambda 1}}+{\displaystyle \frac{\theta 1}{\lambda t}}\right)d\lambda .`$ As we shall prove in Appendix A, (A) can be further rewritten $`(\mathrm{A})=\left[{\displaystyle \frac{\mathrm{}(q+\omega _3)}{4\pi i}}\pi \left({\displaystyle \frac{f_\tau (q)}{f^{}(q)}}\right)^2\right]d\tau +\text{exact form},`$ (46) where $`f^{}(u)`$ denotes the $`u`$-derivative of $`f(u)`$: $`f^{}(u)={\displaystyle \frac{f(u)}{u}}={\displaystyle \frac{\mathrm{}^{}(u)}{e_2e_1}}.`$ (47) For (B) and (C), we have $`(\mathrm{B})`$ $`=`$ $`\left[{\displaystyle \frac{f_\tau (q)}{f^{}(q)}}p+2\pi i\left({\displaystyle \frac{f_\tau (q)}{f^{}(q)}}\right)^2\right]d\tau ,`$ $`(\mathrm{C})`$ $`=`$ $`{\displaystyle \frac{\theta 1}{2(\lambda t)}}dt+{\displaystyle \frac{1}{2}}\left(\kappa _0\mathrm{log}\lambda +\kappa _1\mathrm{log}(\lambda 1)+(\theta 1)\mathrm{log}(\lambda t)\right)`$ $`=`$ $`{\displaystyle \frac{\theta 1}{2(\lambda t)}}dt+\text{exact form}.`$ Thus we find that $`\mu d\lambda Hdt=pdq\stackrel{~}{}{\displaystyle \frac{d\tau }{2\pi i}}+\text{exact form},`$ (48) where $`\stackrel{~}{}=2\pi i{\displaystyle \frac{dt}{d\tau }}\left(H{\displaystyle \frac{\theta 1}{2(\lambda t)}}\right)2\pi i\left[{\displaystyle \frac{\mathrm{}(q+\omega _3)}{4\pi i}}+{\displaystyle \frac{f_\tau (q)}{f^{}(q)}}p+\pi i\left({\displaystyle \frac{f_\tau (q)}{f^{}(q)}}\right)^2\right].`$ (49) Our task is to prove that the transformed Hamiltonian $`\stackrel{~}{}`$ coincides, modulo irrelevant terms, with the Hamiltonian of Manin’s equation. Here “irrelevant” means that the term is a function of $`t`$ only. Such a “non-dynamical” term can be absorbed by the “exact form” part of the foregoing relation of $`1`$-forms, thereby being negligible. Let us evaluate the contribution of $`2\pi i(dt/d\tau )H`$. By Manin’s formula (12) of $`d\tau /dt`$, and also by the identity $`\lambda (\lambda 1)(\lambda t)={\displaystyle \frac{\mathrm{}^{}(q)^2}{4(e_2e_1)^3}},`$ we can rewrite $`2\pi i(dt/d\tau )H`$ as follows: $`2\pi i{\displaystyle \frac{dt}{d\tau }}H`$ $`=`$ $`{\displaystyle \frac{\mathrm{}^{}(q)^2}{2(e_2e_1)^2}}\left[\mu ^2\left({\displaystyle \frac{\kappa _0}{\lambda }}+{\displaystyle \frac{\kappa _1}{\lambda 1}}+{\displaystyle \frac{\theta 1}{\lambda t}}\right)\mu +{\displaystyle \frac{\kappa }{\lambda (\lambda 1)}}\right]`$ $`=`$ $`{\displaystyle \frac{\mathrm{}^{}(q)^2}{2(e_2e_1)^2}}\left[\mu {\displaystyle \frac{1}{2}}\left({\displaystyle \frac{\kappa _0}{\lambda }}+{\displaystyle \frac{\kappa _1}{\lambda 1}}+{\displaystyle \frac{\theta 1}{\lambda t}}\right)\right]^2`$ $`+{\displaystyle \frac{\mathrm{}^{}(q)^2}{2(e_2e_1)^2}}\left[{\displaystyle \frac{1}{4}}\left({\displaystyle \frac{\kappa _0}{\lambda }}+{\displaystyle \frac{\kappa _1}{\lambda 1}}+{\displaystyle \frac{\theta 1}{\lambda t}}\right)^2+{\displaystyle \frac{\kappa }{\lambda (\lambda 1)}}\right].`$ The first term on the right hand side is equal to $`{\displaystyle \frac{1}{2}}\left(p+2\pi i{\displaystyle \frac{f_\tau (q)}{f^{}(q)}}\right)^2={\displaystyle \frac{p^2}{2}}+2\pi i{\displaystyle \frac{f_\tau (q)}{f^{}(q)}}p+\left(2\pi i{\displaystyle \frac{f_\tau (q)}{f^{}(q)}}\right)^2,`$ by which the terms proportional to $`f_\tau (q)/f^{}(q)`$ and its square in the definition of $`\stackrel{~}{}`$ are cancelled out. The transformed Hamiltonian $`\stackrel{~}{}`$ can now be expressed as $`\stackrel{~}{}`$ $`=`$ $`{\displaystyle \frac{p^2}{2}}{\displaystyle \frac{\mathrm{}^{}(q)^2}{2(e_2e_1)^2}}{\displaystyle \frac{(\theta 1)t(t1)(e_2e_1)}{\lambda t}}`$ (50) $`+{\displaystyle \frac{\mathrm{}^{}(q)}{2(e_2e_1)^2}}\left[{\displaystyle \frac{1}{4}}\left({\displaystyle \frac{\kappa _0}{\lambda }}+{\displaystyle \frac{\kappa _1}{\lambda 1}}+{\displaystyle \frac{\theta 1}{\lambda t}}\right)^2+{\displaystyle \frac{\kappa }{\lambda (\lambda 1)}}\right].`$ Note that this is already of the normal form $`p^2/2+\stackrel{~}{V}(q)`$ with the potential $`\stackrel{~}{V}(q)`$ $`=`$ $`{\displaystyle \frac{\mathrm{}^{}(q)^2}{2(e_2e_1)^2}}{\displaystyle \frac{(\theta 1)t(t1)(e_2e_1)}{\lambda t}}`$ (51) $`+{\displaystyle \frac{\mathrm{}^{}(q)}{2(e_2e_1)^2}}\left[{\displaystyle \frac{1}{4}}\left({\displaystyle \frac{\kappa _0}{\lambda }}+{\displaystyle \frac{\kappa _1}{\lambda 1}}+{\displaystyle \frac{\theta 1}{\lambda t}}\right)^2+{\displaystyle \frac{\kappa }{\lambda (\lambda 1)}}\right].`$ What remains is to express $`\stackrel{~}{V}(q)`$ as an explicit function of $`q`$. To this end, we substitute the factor $`\mathrm{}^{}(q)^2/2(e_2e_1)^2`$ by $`2(e_2e_1)\lambda (\lambda 1)(\lambda t)`$, and rewrite the main part of $`\stackrel{~}{V}(q)`$ as a linear combination of $`\lambda `$, $`1/\lambda `$, $`1/(\lambda 1)`$ and $`1/(\lambda t)`$. This leads to the following expression of $`\stackrel{~}{V}(q)`$: $`\stackrel{~}{V}(q)`$ $`=`$ $`{\displaystyle \frac{(\kappa _0+\kappa _1+\theta 1)^24\kappa }{2}}(e_2e_1)\lambda `$ $`{\displaystyle \frac{\kappa _0^2}{2}}{\displaystyle \frac{(e_2e_1)t}{\lambda }}{\displaystyle \frac{\kappa _1^2}{2}}{\displaystyle \frac{(e_2e_1)(1t)}{\lambda 1}}{\displaystyle \frac{(\theta 1)^2+1}{2}}{\displaystyle \frac{(e_2e_1)t(t1)}{\lambda t}}`$ $`{\displaystyle \frac{1}{2}}\mathrm{}(q+\omega _3)+\text{function of }t\text{ only}.`$ The final piece of the ring is the general formula $`\mathrm{}(u+\omega _j)=e_j+{\displaystyle \frac{(e_je_k)(e_je_{\mathrm{}})}{\mathrm{}(u)e_j}}`$ (52) where $`(j,k,l)`$ is a cyclic permutation of $`(1,2,3)`$. This implies that $`{\displaystyle \frac{(e_2e_1)t}{\lambda }}`$ $`=`$ $`\mathrm{}(q+\omega _1)e_1,`$ $`{\displaystyle \frac{(e_2e_1)(1t)}{\lambda 1}}`$ $`=`$ $`\mathrm{}(q+\omega _2)e_2,`$ $`{\displaystyle \frac{(e_2e_1)t(t1)}{\lambda t}}`$ $`=`$ $`\mathrm{}(q+\omega _3)e_3,`$ so that $`\stackrel{~}{V}(q)`$ $`=`$ $`{\displaystyle \frac{(\kappa _0+\kappa _1+\theta 1)^24\kappa }{2}}\mathrm{}(q){\displaystyle \frac{\kappa _0^2}{2}}\mathrm{}(q+\omega _1)`$ (53) $`{\displaystyle \frac{\kappa _1^2}{2}}\mathrm{}(q+\omega _2){\displaystyle \frac{\theta ^2}{2}}\mathrm{}(q+\omega _3)+\text{function of }\tau \text{ only}.`$ Apart from the last term which is negligible, this potential is indeed the same as Manin’s potential $`V(q)`$ (recall the algebraic relations connecting the constants $`\kappa _0`$, etc. and the parameters of $`\mathrm{P}_{\mathrm{VI}}`$). This completes the proof of the theorem. Q.E.D. ### IV.4 Canonical transformation for $`\mathrm{P}_\mathrm{V}`$ This heuristic method for constructing a canonical transformation can be applied to the other Painlevé equations. We here consider the case of $`\mathrm{P}_\mathrm{V}`$. Let $`\lambda `$ be a solution of $`\mathrm{P}_\mathrm{V}`$, $`\mu `$ the canonical conjugate variable, and $`q`$ the corresponding solution of (22). The canonical equation for $`\lambda `$ can be written $`{\displaystyle \frac{d\lambda }{dt}}={\displaystyle \frac{\lambda (\lambda 1)^2}{t}}\left(2\mu {\displaystyle \frac{\kappa _0}{\lambda }}{\displaystyle \frac{\theta _1}{\lambda 1}}+{\displaystyle \frac{\eta _1t}{(\lambda 1)^2}}\right).`$ This equation can be solved for $`\mu `$ as $`\mu ={\displaystyle \frac{1}{2\lambda (\lambda 1)^2}}t{\displaystyle \frac{d\lambda }{dt}}+{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{\kappa _0}{\lambda }}+{\displaystyle \frac{\theta _1}{\lambda 1}}{\displaystyle \frac{\eta _1t}{(\lambda 1)^2}}\right).`$ By differentiating (21) against $`t`$ and using the canonical equation $`tdq/dt=/p=p`$, we obtain the identity $`t{\displaystyle \frac{d\lambda }{dt}}=\sqrt{\lambda }(\lambda 1)p,`$ which can be used to rewrite the expression of $`\mu `$ as $`\mu `$ $`=`$ $`{\displaystyle \frac{p}{2\sqrt{\lambda }(\lambda 1)}}+{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{\kappa _0}{\lambda }}+{\displaystyle \frac{\theta _1}{\lambda 1}}{\displaystyle \frac{\eta _1t}{(\lambda 1)^2}}\right).`$ (54) We now reinterpret (21) and (54) as defining a coordinate transformation $`(\lambda ,\mu )(q,p)`$. This indeed turns out to give a canonical transformation that we have sought for: ###### Theorem 3 (21) and (54) define a canonical transformation that connects $`\mathrm{P}_\mathrm{V}`$ and the $`\mathrm{P}_\mathrm{V}`$-version of Manin’s Hamiltonian system. The canonical coordinates and the Hamiltonians of the two systems obey the equation $`\mu d\lambda Hdt={\displaystyle \frac{1}{2}}\left(pdq{\displaystyle \frac{dt}{t}}\right)+\text{exact form}.`$ (55) Proof. Since $`d\lambda `$ and $`dq`$ are connected by the relation $`d\lambda =\sqrt{\lambda }(\lambda 1)dq,`$ $`\mu d\lambda `$ can be expressed as $`\mu d\lambda `$ $`=`$ $`{\displaystyle \frac{1}{2}}pdq+{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{\kappa _0}{\lambda }}+{\displaystyle \frac{\theta _1}{\lambda 1}}{\displaystyle \frac{\eta _1t}{(\lambda 1)^2}}\right)d\lambda `$ $`=`$ $`{\displaystyle \frac{1}{2}}pdq{\displaystyle \frac{\eta _1}{2(\lambda 1)}}dt+{\displaystyle \frac{1}{2}}d\left(\kappa _0\mathrm{log}\lambda +\theta _1\mathrm{log}(\lambda 1)+{\displaystyle \frac{\eta _1t}{\lambda 1}}\right),`$ so that $`\mu d\lambda Hdt={\displaystyle \frac{1}{2}}\left(pdq\stackrel{~}{}{\displaystyle \frac{dt}{t}}\right)+\text{exact form},`$ (56) where $`\stackrel{~}{}=2Ht+{\displaystyle \frac{\eta _1t}{\lambda 1}}.`$ (57) We can rewrite $`\stackrel{~}{}`$ to a normal form as $`\stackrel{~}{}`$ $`=`$ $`2\lambda (\lambda 1)^2\left[\mu {\displaystyle \frac{1}{2}}\left({\displaystyle \frac{\kappa _0}{\lambda }}+{\displaystyle \frac{\theta _1}{\lambda 1}}{\displaystyle \frac{\theta _1t}{(\lambda 1)^2}}\right)\right]^2`$ (58) $`+2\lambda (\lambda 1)^2\left[{\displaystyle \frac{1}{4}}\left({\displaystyle \frac{\kappa _0}{\lambda }}+{\displaystyle \frac{\theta _1}{\lambda 1}}{\displaystyle \frac{\eta _1t}{(\lambda 1)^2}}\right)^2+{\displaystyle \frac{\kappa }{\lambda (\lambda 1)}}\right]+{\displaystyle \frac{\eta _1t}{\lambda 1}}`$ $`=`$ $`{\displaystyle \frac{p^2}{2}}+\stackrel{~}{V}(q),`$ where $`\stackrel{~}{V}(q)`$ $`=`$ $`{\displaystyle \frac{\lambda (\lambda 1)^2}{2}}\left({\displaystyle \frac{\kappa _0}{\lambda }}+{\displaystyle \frac{\theta _1}{\lambda 1}}{\displaystyle \frac{\eta _1t}{(\lambda 1)^2}}\right)^2+2\kappa (\lambda 1)+{\displaystyle \frac{\eta _1t}{\lambda 1}}.`$ (59) $`=`$ $`\left({\displaystyle \frac{\kappa _0}{2}}+{\displaystyle \frac{\theta _1^2}{2}}+\kappa _1\theta _12\kappa \right){\displaystyle \frac{1}{\mathrm{sinh}^2(q/2)}}+{\displaystyle \frac{\kappa _0^2}{2}}{\displaystyle \frac{1}{\mathrm{cosh}^2(q/2)}}`$ $`+{\displaystyle \frac{\eta _1(\theta _1+1)t}{2}}\mathrm{cosh}(q){\displaystyle \frac{\eta _1^2t^2}{2}}\mathrm{cosh}(2q)+\text{function of }t\text{ only}.`$ Apart from the last negligible term, this coincides with the potential $`V(q)`$ in the statement of the theorem. Q.E.D. ### IV.5 Canonical transformation for $`\mathrm{P}_{\mathrm{IV}}`$ We now consider the case of $`\mathrm{P}_{\mathrm{IV}}`$. Let $`\lambda `$ be a solution of $`\mathrm{P}_{\mathrm{IV}}`$, $`\mu `$ the canonical conjugate variable, and $`q`$ the corresponding solution of (28). The canonical equation for $`\lambda `$ can be written $`{\displaystyle \frac{d\lambda }{dt}}=4\lambda \mu (\lambda ^2+2t\lambda +2\kappa _0),`$ which can be solved for $`\mu `$ as $`\mu ={\displaystyle \frac{1}{4\lambda }}{\displaystyle \frac{d\lambda }{dt}}+{\displaystyle \frac{1}{4}}\left(\lambda +2t+{\displaystyle \frac{2\kappa _0}{\lambda }}\right).`$ By (26) and the canonical equation $`dq/dt=/p=p`$, we have the identity $`{\displaystyle \frac{d\lambda }{dt}}=\sqrt{\lambda }{\displaystyle \frac{dq}{dt}}=\sqrt{\lambda }p,`$ so that $`\mu ={\displaystyle \frac{p}{4\sqrt{\lambda }}}+{\displaystyle \frac{1}{4}}\left(\lambda +2t+{\displaystyle \frac{2\kappa _0}{\lambda }}\right).`$ (60) ###### Theorem 4 (26) and (60) define a canonical transformation that connects $`\mathrm{P}_{\mathrm{IV}}`$ and the $`\mathrm{P}_{\mathrm{IV}}`$-version of Manin’s Hamiltonian system. The canonical coordinates and Hamiltonians of the two systems obey the equation $`\mu d\lambda Hdt={\displaystyle \frac{1}{4}}(pdqdt)+\text{exact form}.`$ (61) Proof. Since $`d\lambda `$ and $`dq`$ are connected by the relation $`d\lambda =\sqrt{\lambda }dq,`$ $`\mu d\lambda `$ can be expressed as $`\mu d\lambda `$ $`=`$ $`{\displaystyle \frac{1}{4}}pdq+{\displaystyle \frac{1}{4}}\left(\lambda +2t+{\displaystyle \frac{2\kappa _0}{\lambda }}\right)d\lambda `$ $`=`$ $`{\displaystyle \frac{1}{4}}pdq{\displaystyle \frac{1}{2}}\lambda dt+{\displaystyle \frac{1}{4}}d\left({\displaystyle \frac{\lambda ^2}{2}}+2t\lambda +2\kappa _0\mathrm{log}\lambda \right),`$ so that $`\mu d\lambda Hdt={\displaystyle \frac{1}{4}}(pdq\stackrel{~}{}dt)+\text{exact form},`$ (62) where $`\stackrel{~}{}=4H+2\lambda .`$ (63) We can rewrite the transformed Hamiltonian $`\stackrel{~}{}`$ to a normal form as $`\stackrel{~}{}`$ $`=`$ $`8\lambda \left[\mu {\displaystyle \frac{1}{2}}\left({\displaystyle \frac{\lambda }{2}}+t+{\displaystyle \frac{\kappa _0}{\lambda }}\right)\right]^2+8\lambda \left[{\displaystyle \frac{1}{4}}\left({\displaystyle \frac{\lambda }{2}}+t+{\displaystyle \frac{\kappa _0}{\lambda }}\right)^2+{\displaystyle \frac{\theta _{\mathrm{}}}{2}}\right]+2\lambda `$ (64) $`=`$ $`{\displaystyle \frac{p^2}{2}}+\stackrel{~}{V}(q),`$ where $`\stackrel{~}{V}(q)`$ $`=`$ $`2\lambda \left({\displaystyle \frac{\lambda }{2}}+t+{\displaystyle \frac{\kappa _0}{\lambda }}\right)^2+4\theta _{\mathrm{}}\lambda +2\lambda `$ (65) $`=`$ $`{\displaystyle \frac{1}{2}}\lambda ^32t\lambda ^22(t^2+\kappa _02\theta _{\mathrm{}}1)\lambda 2\kappa _0^2\lambda ^1`$ $`+\text{function of }t\text{ only}.`$ Substituting $`\lambda =(q/2)^2`$ gives the potential $`V(q)`$ modulo an irrelevant term. Q.E.D. ### IV.6 Canonical transformations for $`\mathrm{P}_{\mathrm{III}}`$ The situation of $`\mathrm{P}_{\mathrm{III}}`$ is somewhat similar to $`\mathrm{P}_\mathrm{V}`$. Let $`\lambda `$, again, be a solution of $`\mathrm{P}_{\mathrm{III}}`$, $`\lambda `$ the canonical conjuage variable, and $`q`$ be the corresponding solution of (34). The canonical equation for $`\lambda `$ takes the form $`{\displaystyle \frac{d\lambda }{dt}}={\displaystyle \frac{\lambda ^2}{t}}\left(2\mu \eta _{\mathrm{}}{\displaystyle \frac{\theta _0}{\lambda }}+{\displaystyle \frac{\eta _0t}{\lambda ^2}}\right),`$ which can be solved for $`\mu `$ as $`\mu ={\displaystyle \frac{t}{2\lambda ^2}}{\displaystyle \frac{d\lambda }{dt}}+{\displaystyle \frac{1}{2}}\left(\eta _{\mathrm{}}+{\displaystyle \frac{\theta _0}{\lambda }}{\displaystyle \frac{\eta _0t}{\lambda ^2}}\right).`$ By differentiating (32) and using the canonical equation $`tdq/dt=/p=p`$, the $`t`$-derivative of $`\lambda `$ can be written $`t{\displaystyle \frac{d\lambda }{dt}}=\lambda p,`$ so that we obtain $`\mu ={\displaystyle \frac{p}{2\lambda }}+{\displaystyle \frac{1}{2}}\left(\eta _{\mathrm{}}+{\displaystyle \frac{\theta _0}{\lambda }}{\displaystyle \frac{\eta _0t}{\lambda ^2}}\right).`$ (66) This relation, again, can be used to define a canonical transformation: ###### Theorem 5 (32) and (66) define a canonical transformation that connects $`\mathrm{P}_{\mathrm{III}}`$ and the $`\mathrm{P}_{\mathrm{III}}`$-version of Manin’s Hamiltonian system. The canonical coordinates and the Hamiltonians of the two systems obey the equation $`\mu d\lambda Hdt={\displaystyle \frac{1}{2}}\left(pdq{\displaystyle \frac{dt}{t}}\right)+\text{exact form}.`$ (67) Proof. Since $`d\lambda `$ and $`dq`$ are connected by the relation $`d\lambda =\lambda dq,`$ $`\mu d\lambda `$ can be written $`\mu d\lambda `$ $`=`$ $`{\displaystyle \frac{1}{2}}pdq+{\displaystyle \frac{1}{2}}\left(\eta _{\mathrm{}}+{\displaystyle \frac{\theta _0}{\lambda }}{\displaystyle \frac{\eta _0t}{\lambda ^2}}\right)d\lambda `$ $`=`$ $`{\displaystyle \frac{1}{2}}pdq{\displaystyle \frac{\eta _0}{2\lambda }}dt+{\displaystyle \frac{1}{2}}d\left(\eta _{\mathrm{}}\lambda +\theta _0\mathrm{log}\lambda +{\displaystyle \frac{\eta _0t}{\lambda }}\right),`$ so that $`\mu d\lambda Hdt={\displaystyle \frac{1}{2}}\left(pdq\stackrel{~}{}{\displaystyle \frac{dt}{t}}\right)+\text{exact form},`$ (68) where $`\stackrel{~}{}=2Ht+{\displaystyle \frac{\eta _0t}{\lambda }}.`$ (69) We can convert the transformed Hamiltonian $`\stackrel{~}{}`$ to a normal form as $`\stackrel{~}{}`$ $`=`$ $`2\lambda ^2\left[\mu {\displaystyle \frac{1}{2}}\left(\eta _{\mathrm{}}+{\displaystyle \frac{\eta _0}{\lambda }}{\displaystyle \frac{\eta _0t}{\lambda ^2}}\right)\right]^2`$ (70) $`+2\lambda ^2\left[{\displaystyle \frac{1}{2}}\left(\eta _{\mathrm{}}+{\displaystyle \frac{\eta _0}{\lambda }}{\displaystyle \frac{\eta _0t}{\lambda ^2}}\right)^2+{\displaystyle \frac{\eta _{\mathrm{}}(\theta _0+\theta _{\mathrm{}})}{2\lambda }}\right]+{\displaystyle \frac{\eta _0t}{\lambda }}`$ $`=`$ $`{\displaystyle \frac{p^2}{2}}+\stackrel{~}{V}(q),`$ where $`\stackrel{~}{V}(q)`$ $`=`$ $`{\displaystyle \frac{\lambda ^2}{2}}\left(\eta _{\mathrm{}}+{\displaystyle \frac{\theta _0}{\lambda }}{\displaystyle \frac{\eta _0t}{\lambda ^2}}\right)^2+\eta _{\mathrm{}}(\theta _0+\theta _{\mathrm{}})\lambda +{\displaystyle \frac{\eta _0t}{\lambda }}`$ (71) $`=`$ $`\eta _{\mathrm{}}\theta _{\mathrm{}}e^q+\eta _0(\theta _0+1)te^q{\displaystyle \frac{\eta _{\mathrm{}}^2}{2}}e^{2q}{\displaystyle \frac{\eta _0^2t^2}{2}}e^{2q}`$ $`+\text{function of }t\text{ only}.`$ Thus, apart from the last irrelevant term, $`\stackrel{~}{V}(q)`$ coincides with the potential $`V(q)`$ in the statement of the theorem. Q.E.D. ### IV.7 Status of $`\mathrm{P}_{\mathrm{II}}`$ and $`\mathrm{P}_\mathrm{I}`$ Let us turn to $`\mathrm{P}_{\mathrm{II}}`$ and $`\mathrm{P}_\mathrm{I}`$. The Hamiltonian of $`\mathrm{P}_\mathrm{I}`$ is already of the normal form $`=\frac{p^2}{2}+V(q)`$ with $`\lambda =q`$, $`\mu =p`$ and $`H=`$. Although this is not the case for $`\mathrm{P}_{\mathrm{II}}`$, one can directly find a canonical transformation that converts the Hamiltonian $`H`$ to a normal form: ###### Theorem 6 A $`\mathrm{P}_{\mathrm{II}}`$-version of Manin’s Hamiltonian system is defined by the Hamiltonian $`={\displaystyle \frac{p^2}{2}}{\displaystyle \frac{1}{2}}\left(q^2+{\displaystyle \frac{t}{2}}\right)^2\alpha q.`$ (72) This system is connected with $`\mathrm{P}_{\mathrm{II}}`$ by the canonical transformation $`\lambda =q,\mu =p+\lambda ^2+{\displaystyle \frac{t}{2}}.`$ (73) The canonical coordinates and the Hamiltonians of the two systems obey the equation $`\mu d\lambda Hdt=pdqdt+\text{exact form}.`$ (74) Proof. The foregoing relation between $`(\lambda ,\mu )`$ and $`(q,p)`$ implies that $`\mu d\lambda =pdq+\left(\lambda ^2+{\displaystyle \frac{t}{2}}\right)d\lambda =pdq{\displaystyle \frac{\lambda }{2}}dt+d\left({\displaystyle \frac{\lambda ^3}{3}}+{\displaystyle \frac{t\lambda }{2}}\right),`$ so that $`\mu d\lambda Hdt=pdq\stackrel{~}{}dt+\text{exact form},`$ (75) where $`\stackrel{~}{}`$ $`=`$ $`H+{\displaystyle \frac{\lambda }{2}}`$ (76) $`=`$ $`{\displaystyle \frac{1}{2}}\left[\mu \left(\lambda ^2+{\displaystyle \frac{t}{2}}\right)\right]^2{\displaystyle \frac{1}{2}}\left(\lambda ^2+{\displaystyle \frac{t}{2}}\right)^2\left(\alpha +{\displaystyle \frac{1}{2}}\right)\lambda +{\displaystyle \frac{\lambda }{2}}`$ $`=`$ $`{\displaystyle \frac{p^2}{2}}{\displaystyle \frac{1}{2}}\left(q^2+{\displaystyle \frac{t}{2}}\right)^2\alpha q.`$ This is nothing but the Hamiltonian in the statement of the theorem. Q.E.D. ## V Multi-component Painlevé equations ### V.1 Inozemtsev Hamiltonians of higher rank The rank $`\mathrm{}`$ version of Inozemtsev’s Hamiltonians have $`\mathrm{}`$ coordinates $`q_1,\mathrm{},q_{\mathrm{}}`$ and canonical conjugate momenta $`p_1,\mathrm{},p_{\mathrm{}}`$. The Hamiltonians of the elliptic, hyperbolic and rational models take the following form : * Elliptic model: $`={\displaystyle \underset{j=1}{\overset{\mathrm{}}{}}}\left({\displaystyle \frac{p_j^2}{2}}+{\displaystyle \underset{n=0}{\overset{3}{}}}g_n^2\mathrm{}(q_j+\omega _n)\right)+g_4^2{\displaystyle \underset{jk}{}}\left(\mathrm{}(q_jq_k)+\mathrm{}(q_j+q_k)\right).`$ * Hyperbolic model: $``$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{\mathrm{}}{}}}\left({\displaystyle \frac{p_j^2}{2}}+{\displaystyle \frac{g_0^2}{\mathrm{sinh}^2(q_j/2)}}+{\displaystyle \frac{g_1^2}{\mathrm{cosh}^2(q_j/2)}}+g_2^2\mathrm{cosh}(q_j)+g_3^2\mathrm{cosh}(2q_j)\right)`$ $`+g_4^2{\displaystyle \underset{jk}{}}\left({\displaystyle \frac{1}{\mathrm{sinh}^2((q_jq_k)/2)}}+{\displaystyle \frac{1}{\mathrm{sinh}^2((q_j+q_k)/2)}}\right).`$ * Rational model: $`={\displaystyle \underset{j=1}{\overset{\mathrm{}}{}}}\left({\displaystyle \frac{p_j^2}{2}}+g_0^2q_j^6+g_1^2q_j^4+g_2^2q_j^2+g_3^2q_j^2\right)+g_4^2{\displaystyle \underset{jk}{}}\left({\displaystyle \frac{1}{(q_jq_k)^2}}+{\displaystyle \frac{1}{(q_j+q_k)^2}}\right).`$ Here $`g_0,g_1,g_2,g_3`$ and $`g_4`$ are coupling constants. The Painlevé-Calogero correspondence for $`\mathrm{P}_{\mathrm{III}}`$, $`\mathrm{P}_{\mathrm{II}}`$ and $`\mathrm{P}_\mathrm{I}`$ suggests the existence of further degeneration of these models. The goal of this section is to extend the the Painlevé-Calogero correspondence to these higher rank models. Since a complete exposition will become inevitably lengthy, we shall illustrate the elliptic and hyperbolic models in detail, leaving the other cases rather sketchy. The strategy is as follows: The point of departure is the Hamiltonian of Inozemtsev’s rank $`\mathrm{}`$ elliptic model. This gives rise to a rank $`\mathrm{}`$ version of Manin’s equation. Starting with this non-autonomous Hamiltonian system, we seek for an analogue of the degeneration process for the Painlevé equations. We can thus obtain six types of non-autonomous Hamiltonian systems. At each stage of the degeneration process, we confirm that the non-autonomous Hamiltonian system on the Calogero side can be mapped, by a canonical transformation, to a multi-component analogue of the Painlevé equation of the corresponding type. ### V.2 Elliptic model and multi-component $`\mathrm{P}_{\mathrm{VI}}`$ We now consider the non-autonomous Hamiltonian system $`2\pi i{\displaystyle \frac{dq_j}{d\tau }}={\displaystyle \frac{}{p_j}},2\pi i{\displaystyle \frac{dp_j}{d\tau }}={\displaystyle \frac{}{q_j}}`$ (77) defined by the Hamiltonian of Inozemtsev’s elliptic model. This is a rank $`\mathrm{}`$ version of Manin’s equation. This non-autonomous system is known to describe a family of isomonodromic deformations on the torus . An honest generalization of the canonical transformation for the case of $`\mathrm{}=1`$ leads to a multi-component version of $`\mathrm{P}_{\mathrm{VI}}`$ as follows: ###### Theorem 7 The time-dependent canonical transformation defined by $`\lambda _j`$ $`=`$ $`{\displaystyle \frac{\mathrm{}(q_j)e_1}{e_2e_1}},`$ $`\mu _j`$ $`=`$ $`{\displaystyle \frac{e_2e_1}{\mathrm{}^{}(q)}}p_j+{\displaystyle \frac{2\pi i(e_2e_1)^2}{\mathrm{}^{}(q_j)^2}}f_\tau (q_j)`$ (78) $`+{\displaystyle \frac{e_2e_1}{2}}\left({\displaystyle \frac{\kappa _0}{\mathrm{}(q_j)e_1}}+{\displaystyle \frac{\kappa _1}{\mathrm{}(q_j)e_2}}+{\displaystyle \frac{\theta 1}{\mathrm{}(q_j)e_3}}\right),`$ and $`t`$ $`=`$ $`{\displaystyle \frac{e_3e_1}{e_2e_1}}.`$ (79) maps (77) to the Hamiltonian system $`{\displaystyle \frac{d\lambda _j}{dt}}={\displaystyle \frac{H}{\mu _j}},{\displaystyle \frac{d\mu _j}{dt}}={\displaystyle \frac{H}{\lambda _j}}`$ (80) with the Hamiltonian $`H`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\lambda _j(\lambda _j1)(\lambda _jt)}{t(t1)}}\left[\mu _j^2\left({\displaystyle \frac{\kappa _0}{\lambda _j}}+{\displaystyle \frac{\kappa _1}{\lambda _j1}}+{\displaystyle \frac{\theta 1}{\lambda _jt}}\right)\mu _j+{\displaystyle \frac{\kappa }{\lambda _j(\lambda _j1)}}\right]`$ (81) $`+{\displaystyle \frac{g_4^2}{2t(t1)}}{\displaystyle \underset{jk}{}}\left[{\displaystyle \frac{\lambda _j(\lambda _j1)(\lambda _jt)+\lambda _k(\lambda _k1)(\lambda _kt)}{8(\lambda _j\lambda _k)^2}}2(\lambda _j+\lambda _k)\right].`$ Proof. The method of proof for the case of $`\mathrm{}=1`$ can be applied to the present case as well, yielding the equality $`{\displaystyle \underset{j=1}{\overset{\mathrm{}}{}}}p_jdq_j{\displaystyle \frac{d\tau }{2\pi i}}={\displaystyle \underset{j=1}{\overset{\mathrm{}}{}}}\mu _jd\lambda _j\stackrel{~}{H}dt+\text{exact form},`$ (82) where $`\stackrel{~}{H}`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\lambda _j(\lambda _j1)(\lambda _jt)}{t(t1)}}\left[\mu _j^2\left({\displaystyle \frac{\kappa _0}{\lambda _j}}+{\displaystyle \frac{\kappa _1}{\lambda _j1}}+{\displaystyle \frac{\theta 1}{\lambda _jt}}\right)\mu _j+{\displaystyle \frac{\kappa }{\lambda _j(\lambda _j1)}}\right]`$ (83) $`+{\displaystyle \frac{g_4^2}{2t(t1)(e_2e_1)}}{\displaystyle \underset{jk}{}}\left(\mathrm{}(q_jq_k)+\mathrm{}(q_j+q_k)\right).`$ What remains is to express the “two-body potential” part in terms of $`\lambda _j`$. To this end, let us recall the addition formula $`\mathrm{}(uv)+\mathrm{}(u+v)=2\mathrm{}(u)2\mathrm{}(v)+{\displaystyle \frac{\mathrm{}^{}(u)^2+\mathrm{}^{}(v)^2}{2(\mathrm{}(u)\mathrm{}(v))^2}}`$ (84) of the $`\mathrm{}`$-function. Applying it to the case where $`(u,v)=(\lambda _j,\lambda _k)`$, and substituting $`\mathrm{}(q_j)`$ $`=`$ $`e_1+(e_2e_1)\lambda _j,`$ $`\mathrm{}(q_k)`$ $`=`$ $`e_1+(e_2e_1)\lambda _k,`$ $`\mathrm{}^{}(q_j)^2`$ $`=`$ $`{\displaystyle \frac{(e_2e_1)^3}{4}}\lambda _j(\lambda _j1)(\lambda _jt),`$ $`\mathrm{}^{}(q_k)^2`$ $`=`$ $`{\displaystyle \frac{(e_2e_1)^3}{4}}\lambda _k(\lambda _k1)(\lambda _kt),`$ we can rewrite the two-body potential terms as $`\mathrm{}(q_jq_k)+\mathrm{}(q_j+q_k)`$ $`=`$ $`2(e_1+(e_2e_1)\lambda _j)2(e_1+(e_2e_1)\lambda _k)`$ (85) $`+{\displaystyle \frac{(e_2e_1)^3}{8}}{\displaystyle \frac{\lambda _j(\lambda _j1)(\lambda _jt)+\lambda _k(\lambda _k1)(\lambda _kt)}{(e_1+(e_2e_1)\lambda _je_1(e_2e_1)\lambda _k)^2}}`$ $`=`$ $`4e_12(e_2e_1)(\lambda _j+\lambda _k)`$ $`+{\displaystyle \frac{e_2e_1}{8}}{\displaystyle \frac{\lambda _j(\lambda _j1)(\lambda _jt)+\lambda _k(\lambda _k1)(\lambda _kt)}{(\lambda _j\lambda _k)^2}}.`$ The first term $`4e_1`$ is non-dynamical, thereby negligible (i.e., can be absorbed by the “exact form” part). Removing these terms from $`\stackrel{~}{H}`$, we obtain the Hamiltonian $`H`$. Q.E.D. ### V.3 Degeneration of elliptic model to hyperbolic model The degeneration of the elliptic model is achieved by letting $`\mathrm{Im}\tau +\mathrm{}`$. Like the degeneration process from $`\mathrm{P}_{\mathrm{VI}}`$ to $`\mathrm{P}_\mathrm{V}`$, this is a kind of scaling limit, namely, the coupling constants $`g_n`$ and the elliptic modulus $`\tau `$ have to be suitably rescaled. To this end, we have to understand the asymptotic behavior of the constants $`e_1,e_2,e_3`$ and the $`\mathrm{}`$-function in the limit as $`\mathrm{Im}\tau +\mathrm{}`$. All necessary data are collected in Appendix B. For instance, the asymptotic expression of $`e_1,e_2`$ and $`e_3`$ imply that $`t=1+{\displaystyle \frac{e_3e_2}{e_2e_1}}=1+16\pi ^2e^{\pi i\tau }+O(e^{2\pi i\tau }).`$ (86) This is indeed consistent with the scaling rule $`t=1+ϵ\stackrel{~}{t}`$ in the degeneration process of $`\mathrm{P}_{\mathrm{VI}}`$ to $`\mathrm{P}_\mathrm{V}`$. Having these data, we now rescale the coupling constants and the elliptic modulus as $`g_0^2=\stackrel{~}{g}_0^2,g_1^2=\stackrel{~}{g}_1^2,g_2^2={\displaystyle \frac{\stackrel{~}{g}_2^2}{ϵ}}+{\displaystyle \frac{\stackrel{~}{g}_3^2}{ϵ^2}},g_3^3={\displaystyle \frac{\stackrel{~}{g}_3^2}{ϵ^2}},g_4^2=\stackrel{~}{g}_4^2`$ (87) and $`16e^{\pi i\tau }=ϵ\stackrel{~}{t},`$ (88) and consider the limit as $`ϵ0`$ while leaving $`\stackrel{~}{g}_n`$ and $`\stackrel{~}{t}`$ finite. Note that letting $`ϵ0`$ amounts to letting $`\mathrm{Im}\tau +\mathrm{}`$. The asymptotic expression of $`\mathrm{}(u)`$ and $`\mathrm{}(u+\omega _n)`$ in Appendix B show that the potential $`V(q)`$ of the elliptic model behaves as $`V(q)`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{\mathrm{}}{}}}\left({\displaystyle \frac{\stackrel{~}{g}_0^2\pi ^2}{\mathrm{sin}^2(\pi q_j)}}+{\displaystyle \frac{\stackrel{~}{g}_1^2\pi ^2}{\mathrm{cos}^2(\pi q_j)}}+{\displaystyle \frac{\stackrel{~}{g}_2^2\pi ^2\stackrel{~}{t}}{2}}\mathrm{cos}(2\pi q_j){\displaystyle \frac{\stackrel{~}{g}_3^2\pi ^2\stackrel{~}{t}^2}{8}}\mathrm{cos}(4\pi q_j)\right)`$ $`+\stackrel{~}{g}_4^2{\displaystyle \underset{jk}{}}\left({\displaystyle \frac{1}{\mathrm{sin}^2(\pi (q_jq_k))}}+{\displaystyle \frac{1}{\mathrm{sin}^2(\pi (q_j+q_k))}}\right)`$ $`+\text{function of }ϵ\text{ and }\stackrel{~}{t}\text{ only}+O(ϵ).`$ Thus, removing negligible terms, we obtain the following Hamiltonian in the limit: $`\stackrel{~}{}`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{\mathrm{}}{}}}\left({\displaystyle \frac{p_j^2}{2}}+{\displaystyle \frac{\stackrel{~}{g}_0^2\pi ^2}{\mathrm{sin}^2(\pi q_j)}}+{\displaystyle \frac{\stackrel{~}{g}_1^2\pi ^2}{\mathrm{cos}^2(\pi q_j)}}+{\displaystyle \frac{\stackrel{~}{g}_2^2\pi ^2\stackrel{~}{t}}{2}}\mathrm{cos}(2\pi q_j){\displaystyle \frac{\stackrel{~}{g}_3^2\pi ^2\stackrel{~}{t}^2}{8}}\mathrm{cos}(4\pi q_j)\right)`$ (89) $`+\stackrel{~}{g}_4^2{\displaystyle \underset{jk}{}}\left({\displaystyle \frac{1}{\mathrm{sin}^2(\pi (q_jq_k))}}+{\displaystyle \frac{1}{\mathrm{sin}^2(\pi (q_j+q_k))}}\right).`$ The asymptotic expression of $`t`$ determines the equation of motion in the limit. In fact, since $`{\displaystyle \frac{d\tau }{dt}}={\displaystyle \frac{\pi }{t(t1)(e_2e_1)}}={\displaystyle \frac{\pi i}{(1+ϵ\stackrel{~}{t})(ϵ\stackrel{~}{t})(\pi ^2+O(ϵ))}}`$ and $`2\pi i{\displaystyle \frac{d}{d\tau }}=2\pi i{\displaystyle \frac{dt}{d\tau }}{\displaystyle \frac{d\stackrel{~}{t}}{dt}}{\displaystyle \frac{d}{dt}}=\left(2\pi ^2\stackrel{~}{t}+O(ϵ^2)\right){\displaystyle \frac{d}{dt}},`$ we find that the equations of motion take the following form: $`2\pi ^2\stackrel{~}{t}{\displaystyle \frac{dq_j}{d\stackrel{~}{t}}}={\displaystyle \frac{\stackrel{~}{}}{p_j}},2\pi ^2\stackrel{~}{t}{\displaystyle \frac{dp_j}{d\stackrel{~}{t}}}={\displaystyle \frac{\stackrel{~}{}}{q_j}}.`$ (90) The final step is to rescale the variables and the Hamiltonian as $`q_j{\displaystyle \frac{q_j}{2\pi i}},p_j\pi iq_j,\stackrel{~}{}\pi ^2\stackrel{~}{},`$ (91) and to rename $`\stackrel{~}{t}`$ and $`\stackrel{~}{}`$ to $`t`$ and $``$. Let us also define the new constants $`\alpha ={\displaystyle \frac{\stackrel{~}{g}_0^2}{2}},\beta ={\displaystyle \frac{\stackrel{~}{g}_1^2}{2}},\gamma ={\displaystyle \frac{\stackrel{~}{g}_2^2}{2}},\delta ={\displaystyle \frac{\stackrel{~}{g}_3^2}{2}},`$ (92) which are to be identified with the four parameters of $`\mathrm{P}_\mathrm{V}`$. The outcome is the non-autonomous Hamiltonian system $`t{\displaystyle \frac{dq_j}{dt}}={\displaystyle \frac{}{p_j}},t{\displaystyle \frac{dp_j}{dt}}={\displaystyle \frac{}{q_j}}`$ (93) with the Hamiltonian $``$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{\mathrm{}}{}}}\left({\displaystyle \frac{p_j^2}{2}}{\displaystyle \frac{\alpha }{\mathrm{sinh}^2(q_j/2)}}{\displaystyle \frac{\beta }{\mathrm{cosh}^2(q_j/2)}}+{\displaystyle \frac{\gamma t}{2}}\mathrm{cosh}(q_j)+{\displaystyle \frac{\delta t^2}{8}}\mathrm{cosh}(2q_j)\right)`$ (94) $`+g_4^2{\displaystyle \underset{jk}{}}\left({\displaystyle \frac{1}{\mathrm{sinh}^2((q_jq_k)/2)}}+{\displaystyle \frac{1}{\mathrm{sinh}^2((q_j+q_k)/2)}}\right).`$ This gives a rank $`\mathrm{}`$ version of the non-autonomous Hamiltonian system on the Calogero side of $`\mathrm{P}_\mathrm{V}`$. Note that the Hamiltoian is essentially the same as the Hamiltonian of Inozemtsev’s hyperbolic model, except that the effective coupling constants are now time-dependent. Remark. The foregoing prescription of scaling limit of the coupling constants and the elliptic modulus is reminiscent of “renormalization” in quantum field theories. In this analogy, one can interpret the equations of motion of the Hamiltonian system as “renormalization group equations”, in which $`\stackrel{~}{t}`$ plays the role of a “mass scale” parameter. ### V.4 Canonical transformation to multi-component $`\mathrm{P}_\mathrm{V}`$ Again, an honest generalization of the canonical transformation for the case of $`\mathrm{}=1`$ leads to a multi-component version of $`\mathrm{P}_\mathrm{V}`$: ###### Theorem 8 The time-dependent canonical transformation defined by $`\sqrt{\lambda _j}`$ $`=`$ $`\mathrm{coth}(q_j/2),`$ $`\mu _j`$ $`=`$ $`{\displaystyle \frac{p_j}{2\sqrt{\lambda _j}(\lambda _j1)}}+{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{\kappa _0}{\lambda _j}}+{\displaystyle \frac{\theta _1}{\lambda _j1}}{\displaystyle \frac{\eta _1t}{(\lambda _j1)^2}}\right)`$ (95) maps (93) to the Hamiltonian system $`{\displaystyle \frac{d\lambda _j}{dt}}={\displaystyle \frac{H}{\mu _j}},{\displaystyle \frac{d\mu _j}{dt}}={\displaystyle \frac{H}{\lambda _j}}`$ (96) with the Hamiltonian $`H`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\lambda _j(\lambda _j1)^2}{t}}\left[\mu _j^2\left({\displaystyle \frac{\kappa _0}{\lambda _j}}+{\displaystyle \frac{\theta _1}{\lambda _j1}}{\displaystyle \frac{\eta _1t}{(\lambda _j1)^2}}\right)\mu _j+{\displaystyle \frac{\kappa }{\lambda _j(\lambda _j1)}}\right]`$ (97) $`+{\displaystyle \frac{g_4^2}{2t}}{\displaystyle \underset{jk}{}}{\displaystyle \frac{2(\lambda _j1)(\lambda _k1)(\lambda _j+\lambda _k)}{(\lambda _j\lambda _k)^2}}.`$ Proof. The method of proof for the case of $`\mathrm{}=1`$ can be used as it is. The outcome is the equality $`{\displaystyle \underset{j=1}{\overset{\mathrm{}}{}}}p_jdq_j{\displaystyle \frac{dt}{t}}=2\left({\displaystyle \underset{j=1}{\overset{\mathrm{}}{}}}\mu _jd\lambda _jHdt\right)+\text{exact form},`$ (98) where $`H`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\lambda _j(\lambda _j1)^2}{t}}\left[\mu _j^2\left({\displaystyle \frac{\kappa _0}{\lambda _j}}+{\displaystyle \frac{\theta _1}{\lambda _j1}}{\displaystyle \frac{\eta _1t}{(\lambda _j1)^2}}\right)\mu _j+{\displaystyle \frac{\kappa }{\lambda _j(\lambda _j1)}}\right]`$ (99) $`+{\displaystyle \frac{g_4^2}{2t}}{\displaystyle \underset{jk}{}}\left({\displaystyle \frac{1}{\mathrm{sinh}^2((q_jq_k)/2)}}+{\displaystyle \frac{1}{\mathrm{sinh}^2((q_j+q_k)/2)}}\right).`$ The two-body potential part can be rewritten by use of the identity $`{\displaystyle \frac{1}{\mathrm{sinh}^2(uv)}}+{\displaystyle \frac{1}{\mathrm{sinh}^2(u+v)}}=4{\displaystyle \frac{\mathrm{cosh}(2u)\mathrm{cosh}(2v)1}{(\mathrm{cosh}(2u)\mathrm{cosh}(2v))^2}}.`$ (100) Substituting $`u=q_j/2`$, $`v=q_k/2`$, and also using the equality $`\mathrm{cosh}(q_j)=(\lambda _j+1)/(\lambda _j1)`$, we find that $`{\displaystyle \frac{1}{\mathrm{sinh}^2((q_jq_k)/2)}}+{\displaystyle \frac{1}{\mathrm{sinh}^2((q_j+q_k)/2)}}={\displaystyle \frac{2(\lambda _j1)(\lambda _k1)(\lambda _j+\lambda _k)}{(\lambda _j\lambda _k)^2}},`$ (101) which gives the two-body potential term in $`H`$. Q.E.D. ### V.5 Other models The degeneration process can be further continued, and leads to four more models that correspond to a multi-component version of $`\mathrm{P}_{\mathrm{IV}}`$, $`\mathrm{P}_{\mathrm{III}}`$, $`\mathrm{P}_{\mathrm{II}}`$ and $`\mathrm{P}_\mathrm{I}`$. Since the details of derivation are more or less parallel, we show the final results only. The Hamiltonian of each model, like those in the foregoing cases, becomes a sum of $`\mathrm{}`$ copies of the one-component Hamiltonian and Calogero-like two-body potential terms. #### V.5.1 Rational model and multi-component $`\mathrm{P}_{\mathrm{IV}}`$ This model can be derived from the hyperbolic model by degeneration. The degeneration process consists of putting the variables and the parameters as $`t=1+2ϵ\stackrel{~}{t},q_j=\pi i+ϵ^{1/2}\stackrel{~}{q_j},p_j={\displaystyle \frac{\stackrel{~}{p_j}}{2ϵ^{1/2}}},`$ (102) and $`\alpha ={\displaystyle \frac{1}{8ϵ^4}},\beta ={\displaystyle \frac{\stackrel{~}{\beta }}{4}},\gamma ={\displaystyle \frac{1}{4ϵ^4}},\delta ={\displaystyle \frac{1}{8ϵ^4}}+{\displaystyle \frac{\stackrel{~}{\alpha }}{2ϵ^2}},`$ (103) and letting $`ϵ0`$ while leaving the “renormalized” quantities $`\stackrel{~}{t}`$, etc. finite. The equations of motion of this model takes the canonical form $`{\displaystyle \frac{dq_j}{dt}}={\displaystyle \frac{}{p_j}},{\displaystyle \frac{dp_j}{dt}}={\displaystyle \frac{}{q_j}}`$ (104) with the Hamiltonian $``$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{\mathrm{}}{}}}\left[{\displaystyle \frac{p_j^2}{2}}{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{q_j}{2}}\right)^62t\left({\displaystyle \frac{q_j}{2}}\right)^42(t^2\alpha )\left({\displaystyle \frac{q_j}{2}}\right)^2+\beta \left({\displaystyle \frac{q_j}{2}}\right)^2\right]`$ (105) $`+g_4^2{\displaystyle \underset{jk}{}}\left({\displaystyle \frac{1}{(q_jq_k)^2}}+{\displaystyle \frac{1}{(q_j+q_k)^2}}\right).`$ The canonical transformation defined by $`\lambda _j=\left({\displaystyle \frac{q_j}{2}}\right)^2,\mu _j={\displaystyle \frac{p_j}{4\sqrt{\lambda _j}}}+{\displaystyle \frac{1}{4}}\left(\lambda _j+2t+{\displaystyle \frac{2\kappa _0}{\lambda _j}}\right)`$ (106) maps the foregoing non-autonomous system to the Hamiltonian system $`{\displaystyle \frac{d\lambda _j}{dt}}={\displaystyle \frac{H}{\mu _j}},{\displaystyle \frac{d\mu _j}{dt}}={\displaystyle \frac{H}{\lambda _j}}`$ (107) with the Hamiltonian $`H={\displaystyle \underset{j=1}{\overset{\mathrm{}}{}}}2\lambda _j^2\left[\mu _j^2\left({\displaystyle \frac{\lambda _j}{2}}+t+{\displaystyle \frac{\kappa _0}{\lambda }}\right)\mu _j+{\displaystyle \frac{\theta _0}{2}}\right]+{\displaystyle \frac{g_4^2}{4}}{\displaystyle \underset{jk}{}}{\displaystyle \frac{2(\lambda _j+\lambda _k)}{(\lambda _j\lambda _k)^2}}.`$ (108) #### V.5.2 Exponential-hyperbolic model and multi-component $`\mathrm{P}_{\mathrm{III}}`$ This model, too, can be derived from the hyperbolic model by degeneration. This degeneration is achieved by the putting the variables and the parameters as $`q_j=\stackrel{~}{q_j}\mathrm{log}{\displaystyle \frac{ϵ}{4}},p_j=\stackrel{~}{p_j},`$ (109) and $`\alpha ={\displaystyle \frac{\stackrel{~}{\alpha }}{4ϵ}}+{\displaystyle \frac{\stackrel{~}{\gamma }}{8ϵ^2}},\beta ={\displaystyle \frac{\stackrel{~}{\gamma }}{8ϵ^2}},\gamma ={\displaystyle \frac{\stackrel{~}{\beta }ϵ}{4}},\delta ={\displaystyle \frac{\stackrel{~}{\delta }ϵ^2}{8}},`$ (110) and letting $`ϵ0`$. The equations of motion of this model takes the canonical form $`t{\displaystyle \frac{dq_j}{dt}}={\displaystyle \frac{}{p_j}},t{\displaystyle \frac{dp_j}{dt}}={\displaystyle \frac{}{q_j}}`$ (111) with the Hamiltonian $``$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{\mathrm{}}{}}}\left({\displaystyle \frac{p_j^2}{2}}{\displaystyle \frac{\alpha }{4}}e^{q_j}+{\displaystyle \frac{\beta t}{4}}e^{q_j}{\displaystyle \frac{\gamma }{8}}e^{2q_j}+{\displaystyle \frac{\delta t^2}{8}}e^{2q_j}\right)`$ (112) $`+g_4^2{\displaystyle \underset{jk}{}}{\displaystyle \frac{1}{\mathrm{sinh}^2((q_jq_k)/2)}}.`$ The canonical transformation defined by $`\lambda _j=e^{q_j},\mu _j={\displaystyle \frac{p_j}{2\lambda _j}}+{\displaystyle \frac{1}{2}}\left(\eta _{\mathrm{}}+{\displaystyle \frac{\theta _0}{\lambda _j}}{\displaystyle \frac{\eta _0t}{\lambda _j^2}}\right)`$ (113) maps the foregoing onn-autonomous system to the Hamiltonian system $`{\displaystyle \frac{d\lambda _j}{dt}}={\displaystyle \frac{H}{\mu _j}},{\displaystyle \frac{d\mu _j}{dt}}={\displaystyle \frac{H}{\lambda _j}}`$ (114) with the Hamiltonian $`H={\displaystyle \underset{j=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\lambda _j^2}{t}}\left[\mu _j^2\left(\eta _{\mathrm{}}+{\displaystyle \frac{\theta _0}{\lambda _j}}{\displaystyle \frac{\eta _0t}{\lambda _j^2}}\right)\mu _j+{\displaystyle \frac{\eta _{\mathrm{}}(\theta _0+\theta _{\mathrm{}})}{2\lambda _j}}\right]+{\displaystyle \frac{g_4^2}{2t}}{\displaystyle \underset{jk}{}}{\displaystyle \frac{4\lambda _j\lambda _k}{(\lambda _j\lambda _k)^2}}.`$ (115) #### V.5.3 Second rational model and multi-component $`\mathrm{P}_{\mathrm{II}}`$ This model can be derived from both the rational model and the exponential-hyperbolic model by degeneration. For the degeneration from the rational model, we write the variables and the parameters as $`t={\displaystyle \frac{1+4^{1/3}ϵ^4\stackrel{~}{t}}{ϵ}},{\displaystyle \frac{q_j}{2}}={\displaystyle \frac{1+2^{1/3}ϵ^2\stackrel{~}{q_j}}{ϵ^{3/2}}},p_j={\displaystyle \frac{4^{2/3}\stackrel{~}{p}_j}{ϵ^{1/2}}}`$ (116) and $`\alpha =2\stackrel{~}{\alpha }{\displaystyle \frac{1}{2ϵ^6}},\beta ={\displaystyle \frac{1}{2ϵ^{12}}},`$ (117) and let $`ϵ0`$. The degeneration from the exponential-hyperbolic model is similarly achieved by putting $`t=1+2ϵ^2\stackrel{~}{t},q_j=2ϵ\stackrel{~}{q_j},p_j={\displaystyle \frac{\stackrel{~}{p_j}}{ϵ}},`$ (118) and $`\alpha ={\displaystyle \frac{1}{2ϵ^6}},\beta ={\displaystyle \frac{1+4ϵ^3\stackrel{~}{\alpha }}{2ϵ^6}},\gamma ={\displaystyle \frac{1}{4ϵ^6}},\delta ={\displaystyle \frac{1}{4ϵ^6}},`$ (119) and again letting $`ϵ0`$. The equations of motion of this model takes the canonical form $`{\displaystyle \frac{dq_j}{dt}}={\displaystyle \frac{}{p_j}},{\displaystyle \frac{dp_j}{dt}}={\displaystyle \frac{}{q_j}}`$ (120) with the Hamiltonian $`={\displaystyle \underset{j=1}{\overset{\mathrm{}}{}}}\left[{\displaystyle \frac{p_j^2}{2}}{\displaystyle \frac{1}{2}}\left(q_j^2+{\displaystyle \frac{t}{2}}\right)^2\alpha q_j\right]+g_4^2{\displaystyle \underset{jk}{}}{\displaystyle \frac{1}{(q_jq_k)^2}}.`$ (121) The canonical transformation defined by $`\lambda _j=q_j,\mu _j=p_j+\lambda _j^2+{\displaystyle \frac{t}{2}}`$ (122) maps the foregoing non-autonomous system to the Hamiltonian system $`{\displaystyle \frac{d\lambda _j}{dt}}={\displaystyle \frac{H}{\mu _j}},{\displaystyle \frac{d\mu _j}{dt}}={\displaystyle \frac{H}{\lambda _j}}`$ (123) with the Hamiltonian $`H={\displaystyle \underset{j=1}{\overset{\mathrm{}}{}}}\left[{\displaystyle \frac{\mu _j^2}{2}}\left(\lambda _j^2+{\displaystyle \frac{t}{2}}\right)\mu _j\left(\alpha +{\displaystyle \frac{1}{2}}\right)\lambda _j\right]+g_4^2{\displaystyle \underset{jk}{}}{\displaystyle \frac{1}{(\lambda _j\lambda _k)^2}}.`$ (124) #### V.5.4 Multi-component $`\mathrm{P}_\mathrm{I}`$ This model can be derived from the second rational model, and takes the same form on both the Painlevé and Calogero sides. The degeneration process is achieved by putting $`t={\displaystyle \frac{6+ϵ^{12}\stackrel{~}{t}}{ϵ^{10}}},q_j={\displaystyle \frac{1+ϵ^6\stackrel{~}{q_j}}{ϵ^5}},p_j={\displaystyle \frac{\stackrel{~}{p_j}}{ϵ}},\alpha =4ϵ^{15}`$ (125) and letting $`ϵ0`$. The equations of motion takes the canonical form $`{\displaystyle \frac{dq_j}{dt}}={\displaystyle \frac{H}{p_j}},{\displaystyle \frac{dp_j}{dt}}={\displaystyle \frac{H}{q_j}}`$ (126) with the Hamiltonian $`H={\displaystyle \underset{j=1}{\overset{\mathrm{}}{}}}\left({\displaystyle \frac{p_j^2}{2}}2q_j^3tq_j\right)+g_4^2{\displaystyle \underset{jk}{}}{\displaystyle \frac{1}{(q_jq_k)^2}}.`$ (127) ## VI Concluding remarks We have shown that the Painlevé-Calogero correspondence persists for all the six Painlevé equations and their multi-component generalizations. The Calogero side of this correspondence is a non-autonomous version of Inozemtsev’s elliptic model and its various degenerations. Those for $`\mathrm{P}_\mathrm{V}`$ and $`\mathrm{P}_{\mathrm{IV}}`$ are a non-autonomous version of Inozemtsev’s hyperbolic and rational models. The others corresponding to $`\mathrm{P}_{\mathrm{III}}`$, $`\mathrm{P}_{\mathrm{II}}`$ and $`\mathrm{P}_\mathrm{I}`$ are further degenerations of the hyperbolic and rational models. The pattern of degeneration on the Calogero side repeats the degeneration diagram $`\begin{array}{ccccccc}\mathrm{P}_{\mathrm{VI}}& & \mathrm{P}_\mathrm{V}& & \mathrm{P}_{\mathrm{IV}}& & \\ & & & & & & \\ & & \mathrm{P}_{\mathrm{III}}& & \mathrm{P}_{\mathrm{II}}& & \mathrm{P}_\mathrm{I}\end{array}`$ (131) of the Painlevé equations. This picture applies to the autonomous systems as well. Actually, such degeneration relations in the autonomous case have been more or less well known to experts of Calogero-Moser systems (see the Introduction of van Diejen’s paper ). The autonomous systems are defined by a Hamiltonian of the same form with the time-dependent coupling constants being replaced by absolute constants (except for the elliptic model, in which case an independent time variable is introduced). Those in the position of the first row of the degeneration diagram are, of course, Inozemtsev’s elliptic, hyperbolic and rational models (see Section 5). Those in the position of $`\mathrm{P}_{\mathrm{III}}`$ and $`\mathrm{P}_{\mathrm{II}}`$ are defined by the following Hamiltonians: * Exponential-hyperbolic model: $`={\displaystyle \underset{j=1}{\overset{\mathrm{}}{}}}\left({\displaystyle \frac{p_j^2}{2}}+g_0^2e^{q_j}+g_1^2e^{2q_j}+g_2^2e^{q_j}+g_3^2e^{2q_j}\right)+g_4^2{\displaystyle \underset{jk}{}}{\displaystyle \frac{1}{\mathrm{sinh}^2((q_jq_k)/2)}}.`$ * Second rational model: $`={\displaystyle \underset{j=1}{\overset{\mathrm{}}{}}}\left({\displaystyle \frac{p_j^2}{2}}+g_0^2q_j^4+g_1^2q_j^3+g_2^2q_j^2+g_3^2q_j\right)+g_4^2{\displaystyle \underset{jk}{}}{\displaystyle \frac{1}{(q_jq_k)^2}}.`$ The Hamiltonian in the position of $`\mathrm{P}_\mathrm{I}`$ is redundant in the automonous case, because it is a specialization, rather than a degeneration, of the last Hamiltonian. Note that the Hamiltonian of the second rational model is a quartic perturbation of the usual ($`A_{\mathrm{}}`$ type) rational Calogero Hamiltonian. According to recent work of Caseiro, Françoise and Sasaki , such a quartic (integrable) perturbation always exists for any rational Calogero-Moser system. Inozemtsev’s rational model, which is a sextic perturbation of the $`D_{\mathrm{}}`$ type rational Calogero-Moser system, might admit a similar interpretation. Back to the Painlevé equations, the extended Painlevé-Calogero correspondence raises many interesting problems. A central issue will be to find an isomonodromic description of the multi-component Painlevé equations. If such an isomonodromic description does exist, it should be related to a new geometric structure. ### Acknwlegements I am grateful to Marta Mazzocco, Davide Guzzetti, Kazuo Okamoto, Ryu Sasaki, Shun Shimomura and Jan Felipe van Diejen, for useful comments. This work was partly supported by the Grant-in-Aid for Scientific Researches (No. 10640165) from the Ministry of Education, Science and Culture. ## Appendix A Proof of (46) Let us introduce the two auxiliary functions $`g(u)={\displaystyle \frac{f_\tau (u)}{f^{}(u)}},h(u)={\displaystyle \frac{\vartheta ^{}(u+\omega _1)}{\vartheta (u+\omega _1)}},`$ (A.1) associated with the function $`f(u)={\displaystyle \frac{\mathrm{}(u)e_1}{e_2e_1}}`$ (A.2) and the standard elliptic theta function $`\vartheta (u)={\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}\mathrm{exp}(\pi i\tau n^2+2\pi inu).`$ (A.3) ###### Lemma 1 $`g(u)`$ is a meromorphic function on the $`u`$-plane with additive quasi-periodicity $`g(u+1)=g(u),g(u+\tau )=g(u)1.`$ (A.4) All poles are of the first order and contained in the lattice $`\omega _3++\tau `$. Furthermore, $`g(u)`$ has zeros at $`u=0`$ and $`u=\omega _1`$. Proof. Since $`f(u)`$ is a doubly periodic function with primitive periods 1 and $`\tau `$, $`f^{}(u)`$ and $`f_\tau (u)`$ transform as $`f^{}(u+1)=f^{}(u),`$ $`f^{}(u+\tau )=f^{}(u),`$ $`f_\tau (u+1)=f_\tau (u),`$ $`f_\tau (u+\tau )=f_\tau (u)f^{}(u)`$ under the shift by $`1`$ and $`\tau `$. This implies the additive quasi-periodicity of $`g(u)`$. Furthermore, by the construction, $`g(u)`$ is a meromorphic function on the $`u`$-plane, and all possible poles are of the first order and located at the points of $`\omega _k++\tau `$. Let us examine the behavior of $`g(u)`$ at the representative points $`u=\omega _0,\omega _1,\omega _2,\omega _3`$: * As $`u\omega _0=0`$, $`f(u)={\displaystyle \frac{1}{(e_2e_1)u^2}}+O(1),`$ thereby $`f^{}(u)={\displaystyle \frac{2}{(e_2e_1)u^3}}+O(1),f_\tau (u)={\displaystyle \frac{e_{2,\tau }e_{1,\tau }}{(e_2e_1)^2u^2}}+O(1),`$ so that $`g(u)`$ has rather a zero at $`u=0`$: $`g(u)=O(u).`$ (A.5) * As $`u\omega _1=\frac{1}{2}`$, $`f(u)`$ $`=`$ $`{\displaystyle \frac{1}{e_2e_1}}\left(\mathrm{}(\omega _1)e_1+\mathrm{}^{}(\omega _1)(u\omega _1)+O((u\omega _1)^2)\right)`$ $`=`$ $`O((u\omega _1)^2),`$ thereby $`f^{}(u)=O(u\omega _1),f_\tau (u)=O((u\omega _1)^2),`$ so that $`g(u)`$ has another zero at $`u=\omega _1`$: $`g(u)=O(u\omega _1).`$ (A.6) * As $`u\omega _2=\frac{1}{2}+\frac{\tau }{2}`$, $`f(u)`$ $`=`$ $`{\displaystyle \frac{1}{e_2e_1}}\left(\mathrm{}(\omega _2)e_1+\mathrm{}^{}(\omega _2)(u+\omega _2)+O((u+\omega _2)^2)\right)`$ $`=`$ $`O((u+\omega _2)^2),`$ thereby $`f^{}(u)=O(u+\omega _2),f_\tau (u)=O(u+\omega _2),`$ so that $`g(u)`$ behaves as $`g(u)=O(1).`$ (A.7) * As $`u\omega _3=\frac{\tau }{2}`$, $`f(u)`$ $`=`$ $`{\displaystyle \frac{1}{e_2e_1}}\left(\mathrm{}(\omega _3)e_1+\mathrm{}^{}(\omega _3)(u\omega _3)+O((u\omega _3)^2)\right)`$ $`=`$ $`t+O((u\omega _3)^2),`$ thereby $`f^{}(u)=O(u\omega _3),f_\tau (u)=O(1),`$ so that $`g(u)`$ turns out to have a pole of the first order at $`u=\omega _3`$: $`g(u)=0((u\omega _3)^1).`$ (A.8) The behavior of $`g(u)`$ at the other points of $`\omega _n++\tau `$ can be deduced from these results by the additive quasi-periodicity of $`g(u)`$. Q.E.D. ###### Lemma 2 $`h(u)`$ is a meromorphic function on the $`u`$-plane with additive quasi-periodicity $`h(u+1)=h(u),h(u+\tau )=h(u)2\pi i.`$ (A.9) All poles are of the first order and contained in the lattice $`\omega _3++\tau `$. Furthermore, $`h(u)`$ has zeros at $`u=0`$ and $`u=\omega _1`$. Proof. Let us recall the fundamental properties of $`\vartheta (u)`$: * $`\vartheta (u)`$ is an entire function on the $`u`$-plane with zeros of the first order at the lattice points $`\omega _2+m+n\tau `$ ($`m,n`$). * $`\vartheta (u)`$ is quasi-periodic, $`\vartheta (u+1)=\vartheta (u),\vartheta (u+\tau )=e^{\pi i\tau 2\pi iu}\vartheta (u).`$ * $`\theta (u)`$ and $`\vartheta (u+1/2)`$ are even under the reflection $`uu`$. All the properties of $`h(u)`$ in the statement of the lemma are an immediate consequence of these properties of $`\vartheta (u)`$. Q.E.D. ###### Lemma 3 The function $`f(u)`$ satisfies the equation $`2\pi i{\displaystyle \frac{f_\tau (u)}{f^{}(u)}}={\displaystyle \frac{\vartheta ^{}(u+\omega _1)}{\vartheta (u+\omega _1)}},`$ (A.10) where the prime stands for $`/u`$. Proof. The foregoing properties of $`g(u)`$ and $`h(u)`$ imply the following: * $`2\pi ig(u)h(u)`$ is a doubly periodic meromorphic function with fundamental period $`1`$ and $`\tau `$. * All poles of $`2\pi ig(u)h(u)`$ are of the first order and contained in the lattice $`\omega _3++\tau `$. * $`2\pi ig(u)h(u)`$ has zeros at $`u=0`$ and $`u=\omega _1`$. The first two properties imply that $`2\pi ig(u)h(u)`$ is a constant. By the last one, this constant has to be zero. We thus find that $`2\pi ig(u)h(u)=0`$. Q.E.D. ###### Lemma 4 $`\vartheta (u)`$ satisfies the equation $`\left(\mathrm{log}\vartheta (u+\omega _1)\right)^{\prime \prime }=\mathrm{}(u+\omega _3)+\text{function of }\tau \text{ only}.`$ (A.11) Proof. The aforementioned complex analytic properties of $`\vartheta (u)`$ imply the following: * $`\left(\mathrm{log}\vartheta (u+\omega _1)\right)^{\prime \prime }`$ is a doubly periodic meromorphic function with primitive period $`1`$ and $`\tau `$. * All poles of this meromorphic function are contained in the lattice $`\omega _3++\tau `$. * As $`u\omega _3`$, this function behaves as $`\left(\mathrm{log}\vartheta (u+\omega _1)\right)^{\prime \prime }={\displaystyle \frac{1}{(u+\omega _3)^2}}+O(1).`$ The function $`\mathrm{}(u+\omega _3)`$, too, has these properties. Accordingly, their difference is a constant function on the $`u`$-plane, namely, a function of $`\tau `$ only. Q.E.D. We now return to the proof of (46). By the third lemma, we have the identity $`2\pi i{\displaystyle \frac{f_\tau (u)}{f^{}(u)}}du={\displaystyle \frac{\vartheta ^{}(u+\omega _1)}{\vartheta (u+\omega _1)}}du={\displaystyle \frac{d\vartheta (u+\omega _1)}{\vartheta (u+\omega _1)}}{\displaystyle \frac{\vartheta (u+\omega _1)/\tau }{\vartheta (u+\omega _1)}}d\tau `$ (A.12) On the other hand, the well known “heat equation” $`4\pi i{\displaystyle \frac{\vartheta (u)}{\tau }}=\vartheta (u)^{\prime \prime }`$ (A.13) implies that $`{\displaystyle \frac{\vartheta (u+\omega _1)/\tau }{\vartheta (u+\omega _1)}}={\displaystyle \frac{1}{4\pi }}{\displaystyle \frac{\vartheta (u+\omega _1)^{\prime \prime }}{\vartheta (u+\omega _2)}}={\displaystyle \frac{1}{4\pi i}}\left[\left(\mathrm{log}\vartheta (u+\omega _1)\right)^{\prime \prime }+\left({\displaystyle \frac{\vartheta ^{}(u+\omega _1)}{\vartheta (u+\omega _1)}}\right)^2\right].`$ By the third and forth lemmas, the last line can be rewritten $`{\displaystyle \frac{1}{4\pi i}}\left[\mathrm{}(u+\omega _3)+\left(2\pi i{\displaystyle \frac{f_\tau (u)}{f^{}(u)}}\right)^2\right]+\text{function of }\tau \text{ only},`$ so that $`2\pi i{\displaystyle \frac{f_\tau (u)}{f^{}(u)}}du={\displaystyle \frac{1}{4\pi i}}\left[\mathrm{}(u+\omega _3)\left(2\pi i{\displaystyle \frac{f_\tau (u)}{f^{}(u)}}\right)^2\right]d\tau +\text{exact form}.`$ (A.14) Substituting $`u=q`$ gives (46) ## Appendix B Asymptotics of elliptic functions The asymptotic behavior of the $`\mathrm{}`$-function $`\mathrm{}(u)`$, the shifted $`\mathrm{}`$-functions $`\mathrm{}(u+\omega _k)`$ and the constants $`e_k=\mathrm{}(\omega _k)`$, in the limit as $`\mathrm{Im}\tau +\mathrm{}`$, can be deduced from the well known formula $`\mathrm{}(u)={\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\pi ^2}{\mathrm{sin}^2(\pi (u+n\tau ))}}{\displaystyle \frac{\pi ^2}{3}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{2\pi ^2}{\mathrm{sin}^2(\pi n\tau )}}.`$ (B.1) Let us first consider the asymptotic behavior of $`\mathrm{}(u)`$ itself. The constant ($`n=0`$) term in the first sum is of order $`1`$ and the $`n`$-th term is of order $`e^{2n\pi i\tau }`$. Similarly, the $`n`$-th term in the second sum is of order $`e^{2n\pi i\tau }`$. Therefore $`\mathrm{}(u)={\displaystyle \frac{\pi ^2}{\mathrm{sin}^2(\pi u)}}{\displaystyle \frac{\pi ^2}{3}}+O(e^{2\pi i\tau }).`$ (B.2) A similar estimate leads to the following asymptotic expression for the shifted $`\mathrm{}`$-functions: $`\mathrm{}(u+\omega _1)`$ $`=`$ $`{\displaystyle \frac{\pi ^2}{\mathrm{cos}^2(\pi u)}}{\displaystyle \frac{\pi ^2}{3}}+O(e^{2\pi i\tau }),`$ $`\mathrm{}(u+\omega _2)`$ $`=`$ $`{\displaystyle \frac{\pi ^2}{3}}+8\pi ^2\mathrm{cos}(2\pi u)e^{\pi i\tau }+O(e^{2\pi i\tau }),`$ $`\mathrm{}(u+\omega _3)`$ $`=`$ $`{\displaystyle \frac{\pi ^2}{3}}8\pi ^2\mathrm{cos}(2\pi u)e^{2\pi i\tau }+O(e^{2\pi i\tau }).`$ (B.3) In fact, the degeneration process of the elliptic model requires us to know the asymptotic expression of $`\mathrm{}(u+\omega _2)+\mathrm{}(u+\omega _3)`$ to the order $`e^{2\pi i\tau }`$. This can be achieved by the following calculations: $`\mathrm{}(u+\omega _2)+\mathrm{}(u+\omega _3)`$ (B.4) $`=`$ $`{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\pi ^2}{\mathrm{cos}^2(u+\frac{\tau }{2}+n\tau )\mathrm{sin}^2(u+\frac{\tau }{2}+n\tau )}}{\displaystyle \frac{2\pi ^2}{3}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{4\pi ^2}{\mathrm{sin}^2(\pi n\tau )}}`$ $`=`$ $`{\displaystyle \frac{2\pi ^2}{3}}32\pi ^2\mathrm{cos}(2\pi u)e^{2\pi i\tau }+16\pi ^2e^{2\pi i\tau }+O(e^{3\pi i\tau }).`$ We now consider the constants $`e_k`$. For instance, $`e_1`$ can be written $`e_1`$ $`=`$ $`{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\pi ^2}{\mathrm{cos}^2(\pi n\tau )}}{\displaystyle \frac{\pi ^2}{3}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{2\pi ^2}{\mathrm{sin}^2(\pi n\tau )}}`$ (B.5) $`=`$ $`{\displaystyle \frac{2}{3}}\pi ^2+{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{2\pi ^2}{\mathrm{cos}^2(\pi n\tau )}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{2\pi ^2}{\mathrm{sin}^2(\pi n\tau )}}.`$ The constant $`2\pi ^2/3`$ becomes the leading term; the leading ($`n=1`$) terms of the last two series give the next-leading term of the order $`e^{2\pi i\tau }`$. $`e_2`$ and $`e_3`$ can be similarly analyzed. Thus the following asymptotic formulas are obtained: $`e_1`$ $`=`$ $`{\displaystyle \frac{2\pi ^2}{3}}+16\pi ^2e^{2\pi i\tau }+O(e^{4\pi i\tau }),`$ $`e_2`$ $`=`$ $`{\displaystyle \frac{\pi ^2}{3}}+8\pi ^2e^{\pi i\tau }+O(e^{2\pi i\tau })`$ $`e_3`$ $`=`$ $`{\displaystyle \frac{\pi ^2}{3}}8\pi ^2e^{\pi i\tau }+O(e^{2\pi i\tau }).`$ (B.6) In particular, $`e_2e_1\pi ^2`$, as expected.
warning/0004/hep-th0004199.html	ar5iv	text	# 1 Introduction ## 1 Introduction Studying brane solutions of low energy supergravity equations is proved to be an important way of searching non-perturbative properties of string/M theories. These solutions play a crucial role both in the conjectured duality symmetries between seemingly different string/M theories and in the AdS/CFT correspondence. Generically, a brane solution has the interpretation of a $`p`$-dimensional extended black hole, which can be characterized by few constants like the mass and the charges of the antisymmetric tensor fields. The metric along the world-volume directions have Poincare, and along the transverse directions have rotational invariances. The geometry is asymptotically flat and the number of transverse directions dictates the radial coordinate dependence of the metric functions. By smearing out some transverse directions, solutions with slower fall-off properties can also be constructed (for a review see, for instance, ) Generalizations of the usual brane solutions have been studied from different points of view. In , it has been shown that the transverse 7-sphere of the membrane solution can be replaced with any Einstein manifold. The brane solutions having transverse hyper-Kähler spaces have been studied in . In , a fivebrane solution wrapping on the manifold K3 has been constructed. Brane solutions with Ricci flat world-volumes have been obtained in . In , more general brane solutions with curved world-volume directions have been studied. An intrinsic metric, representing non-trivially embedded D3-branes, has been given in . Brane solutions which are product of an AdS space with an Einstein space have been constructed in . In this paper, we will systematically extend these observations by using a theorem proved in which states that, under certain conditions, the existence of a Killing spinor implies the Einstein equations. This enables one to concentrate on the first order Killing spinor equations without worrying about more complicated, second order Einstein equations. The spaces having Killing spinors will play a crucial role in the constructions. The organization of the paper is as follows. In section 2, we review the theorem of . In sections 3 and 4, we specifically consider M2, M5 and D3-branes. In section 3, we generalize the well known solutions and obtain branes having Ricci flat world-volumes and smeared transverse directions, and non-spherical cross sections. In section 4, using Kähler forms of Ricci flat Kähler spaces, we construct intersecting brane solutions which can also be interpreted as wrapping branes over the cycles determined by Kähler forms. In section 5, we show that a generic brane background still obeys the field equations when certain directions in the metric are replaced with more general spaces. In section 6, we first construct new, singular, Ricci flat manifolds, which have covariantly constant spinors, as (generalized) cones over U(1) bundles over Ricci flat Kähler manifolds. These manifolds give rise to new supersymmetric compactifications of non-gauged supergravities. In the same section, we find M2 and D3-brane solutions which asymptotically approach these vacua. We conclude with some brief remarks in section 7. ## 2 BPS solutions from Killing spinor equations: a theorem It is well known that existence of a covariantly constant spinor on a Euclidean manifold implies Ricci flatness. It turns out, this observation has a very useful generalization in the supergravity theory context . Let us consider a bosonic background ($`g_{MN}`$,$`F_{MNPQ}`$) of $`D=11`$ supergravity theory which obeys: $$R_{MN}=\frac{1}{3}(F_M{}_{}{}^{PQR}F_{NPQR}^{}\frac{1}{12}g_{MN}F^{PQRS}F_{PQRS}),$$ (2.1) $$_QF^{QMNP}=\frac{1}{(24)^2}ϵ^{MNPA_1\mathrm{}A_8}F_{A_1..A_4}F_{A_5..A_8}.$$ (2.2) The linearized Rarita-Schwinger equations on this background may be written as: $$\mathrm{\Gamma }^{MNP}D_N\psi _P=0,$$ (2.3) where the supercovariant derivative $`D_M`$ is given by $$D_M=_M+\frac{1}{144}(\mathrm{\Gamma }^{PQRS}{}_{M}{}^{}\frac{1}{8}\delta _M^P\mathrm{\Gamma }^{QRS})F_{PQRS},$$ (2.4) and $`_M`$ is the usual covariant derivative acting on spinors. Let us further consider a linearized spin 3/2 field $`\psi _M`$, which is obtained by the action of supercovariant derivative on an arbitrary spinor $`ϵ`$: $$\psi _M=D_Mϵ.$$ (2.5) Due to the invariance of D=11 supergravity at the linearized fermionic level, this spin 3/2 field solves (2.3). To verify this claim, we insert (2.5) into (2.3) and, $`using`$ $`only`$ the 4-form field equations, obtain $$\mathrm{\Gamma }_M{}_{}{}^{NP}D_{N}^{}D_Pϵ=\frac{1}{2}(G_{MN}T_{MN})\mathrm{\Gamma }^Nϵ=0,$$ (2.6) where $`G_{MN}=R_{MN}\frac{1}{2}g_{MN}R`$ is the Einstein tensor and $$T_{MN}=\frac{1}{3}(F_M{}_{}{}^{PQR}F_{NPQR}^{}\frac{1}{8}g_{MN}F^{PQRS}F_{PQRS})$$ (2.7) is the energy momentum tensor. In setting (2.6) to zero, we have used the fact that the background is chosen to obey Einstein equations (2.1). Now consider a background which satisfies the 4-form field equations (2.2) but $`not`$ necessarily obeys the Einstein equations. Let us further assume the existence of at least one Killing spinor obeying $$D_Mϵ_0=0.$$ (2.8) Following (2.6), it is easy to see that $$(G_{MN}T_{MN})\mathrm{\Gamma }^Nϵ_0=0.$$ (2.9) Note that in deriving (2.6), we have $`only`$ used the 4-form field equations which are also assumed to be satisfied by the new background. (2.6) has been set to zero since that background was chosen to obey Einstein equations. On the other hand, (2.9) is satisfied since $`ϵ_0`$ is the Killing spinor. Therefore, (2.9) is still valid, even if the background is not chosen to obey Einstein equations. Although (2.9) is very close to the Einstein equations, one needs to impose further conditions to proceed due to the Lorentzian signature of the metric. Since $`(G_{MN}T_{MN})`$ is symmetric, one can point-wise diagonalize it by elements of the group O(11)<sup>2</sup><sup>2</sup>2In claiming this, we assume that there is no topological obstruction to have a continuous map from space-time to the group manifold O(11). . However, this does not give any new information, since O(11) transformation of O(1,10) gamma matrices are not nice objects. Therefore, one should be able to split space and time directions in (2.9). It turns out it is sufficient to assume existence of an orthonormal basis such that $`(G_{0i}T_{0i})=0`$, where $`0`$ is the time and $`i`$ is the spatial direction. When the indices of (2.9) refer to this orthonormal basis, $`M=0`$ component implies $`G_{00}T_{00}=0`$. Along spatial directions, when $`M=i`$, one can point-wise diagonalize the symmetric matrix $`(G_{ij}T_{ij})`$ by O(10) transformations, under which the spatial gamma matrices still satisfy the same Clifford algebra and thus invertible. Therefore, the diagonal entries of $`(G_{ij}T_{ij})`$ should also be zero, which then implies Einstein equations $`G_{MN}T_{MN}=0`$. Summarizing above considerations, we obtain the following result: (i) if a background is known to solve the 4-form field equations (2.2), (ii) if there exist an orthonormal basis such that $`G_{0i}T_{0i}=0`$; <sup>3</sup><sup>3</sup>3Here, we relax the corresponding condition imposed in . I would like to thank C.N. Pope for the discussions about this point. (iii) if it is known that there exist a Killing spinor on this background, then this background, being preserve some fraction of supersymmetries, also solves the Einstein equations of D=11 supergravity. The main advantage of the theorem for applications is that, to find a solution of the second order Einstein equations, one can instead concentrate on the first order Killing spinor equations which are of course easier to solve. One can also convince himself that it is not hard to satisfy conditions (i) and (ii). Indeed, in constructing $`p`$-brane solutions, one starts with ansatzs having the property (i) and (ii) (see, for instance, ). As we will see for a moment, condition (ii) can easily be satisfied by choosing the background to be static. We finally note that a solution obtained in this fashion already preserves some fraction of supersymmetries. Although explicitly proven for $`D=11`$ supergravity, it is possible to argue that such a theorem can be established for all supergravities. The linearized supersymmetry invariance of any supergravity theory requires $$\mathrm{\Gamma }^{MNP}D_ND_Pϵ=0,$$ (2.10) where $`D_M=_M+\mathrm{}.`$ is the supercovariant derivative, $`ϵ`$ is an arbitrary spinor and bosonic fields refer to a background which obeys equations of motion. In (2.10), the derivatives of the metric are expected to appear in the combination of the Einstein tensor<sup>4</sup><sup>4</sup>4This is due to the identity $`\mathrm{\Gamma }_M{}_{}{}^{NP}_{N}^{}_Pϵ=1/2G_{MN}\mathrm{\Gamma }^Nϵ`$. . Therefore, after imposing all but Einstein equations, (2.10) should become $$\mathrm{\Gamma }_M{}_{}{}^{NP}D_{N}^{}D_Pϵ=\frac{1}{2}(G_{MN}T_{MN})\mathrm{\Gamma }^Nϵ=0,$$ (2.11) where $`T_{MN}`$ is the appropriate energy momentum tensor of the theory and in the last step we have used the fact that the background obeys Einstein equations. With this information, it is very easy to prove the theorem; dropping the condition that the background satisfies Einstein equations and further assuming the existence of a Killing spinor, one obtains (2.11) evaluated for the Killing spinor. With condition (iii), this implies the Einstein equations. Therefore, if we replace (i) with ($`i^{^{}}`$) if the background satisfies all but Einstein equations, then together with conditions (ii) and (iii) above , the result of the theorem should apply to all supergravities. To conclude this section, we finally comment on a simple way to satisfy condition (ii). For a static background there exist an orthonormal basis such that $`G_{0i}=0`$. To see this, using the fact that the background is static, we write the line element as $$ds^2=A^2(dt^2)+ds_M^2,$$ (2.12) where $`A`$ is a time independent function and $`ds_M^2`$ is a line element on the Euclidean manifold $`M`$. One can easily show that, when expressed in the basis $`e^0=Adt`$ and $`e^i`$ (an orthonormal basis on $`M`$), the Ricci tensor of $`ds^2`$ obeys $`R_{0i}=0`$. This gives $`G_{0i}=0`$ since in any orthonormal basis $`g_{0i}=0`$. On the other hand, if the matter fields are chosen to be static, then $`T_{0i}=0`$<sup>5</sup><sup>5</sup>5This can be regarded as the definition for matter fields to be static.. Therefore, one immediate way to satisfy condition (ii) is to take background fields to be static. However, one should not rule out existence of interesting non-static cases for which only the combination $`(G_{0i}T_{0i})=0`$. ## 3 Generalized brane solutions As mentioned in the introduction, a generic $`p`$-brane solution has Poincare invariance on the world-volume and spherical symmetry along transverse directions. In this section, we will focus on the M2, M5 and D3-branes and try to obtain generalizations of the well known solutions, by using the theorem of section 2. ### Definitions Let $`L_d`$, $`M_m`$ and $`X_n`$ be Lorentzian Ricci flat, Euclidean Ricci flat and Einstein manifolds of dimensions $`d`$, $`m`$ and $`n`$, respectively. We will denote the basis one-forms on these spaces as $`e^\mu `$, $`e^a`$ and $`e^\alpha `$, where the indices $`\mu ,\nu \mathrm{}`$ refer to $`L_d`$; $`a,b..`$ refer to $`M_m`$ and $`\alpha ,\beta ..`$ refer to $`X_n`$. We will assume that $`L_d`$ and $`M_m`$ have covariantly constant Killing spinors, $$_\mu ϵ=0,$$ (3.1) $$_aϵ=0,$$ (3.2) and $`X_n`$ has Killing spinors obeying $$_\alpha ϵ=\frac{1}{2}\mathrm{\Gamma }_\alpha \mathrm{\Gamma }_rϵ,$$ (3.3) where we set the “inverse radius” of $`X_n`$ to 1. In these equations, we will not restrict representations of $`\mathrm{\Gamma }`$-matrices to be irreducible. Indeed, we will view the Clifford algebras on $`L_d`$, $`M_m`$ and $`X_n`$ as sub-algebras of a bigger Clifford algebra. For instance, in (3.3), $`\mathrm{\Gamma }_r`$ can be any element of this bigger algebra which anti-commutes with $`\mathrm{\Gamma }_\alpha `$ and squares to identity. We will write the line elements on $`L_d`$, $`M_m`$ and $`X_n`$ as $`ds_{L_d}^2`$, $`ds_{M_m}^2`$ and $`ds_{X_n}^2`$, respectively. The volume forms will be denoted by $`V_L`$, $`V_M`$ and $`V_X`$. We will use to denote differentiation with respect to the argument of the function, and $``$ to mean equality up to a constant. The indices in tensor equations will refer to the tangent space. Specifically, $`0`$ and $`i`$ will be used for time-like and spatial directions, respectively. ### M2-brane Let us start from the eleven dimensional membrane. Our aim is to write an ansatz, which satisfies the conditions (i) and (ii) of the theorem of section 2, and then work out the Killing spinor equations. For the metric and 4-form field we assume $`ds^2`$ $`=`$ $`A^2ds_{L_3}^2+B^2dr^2+C^2ds_{M_m}^2+D^2ds_{X_n}^2,`$ (3.4) $`F`$ $``$ $`V_MV_X,`$ (3.5) where $`m+n=7`$, $``$ is the Hodge dual corresponding to $`ds^2`$ and the metric functions $`A,B,C`$ and $`D`$ are chosen to depend only on the coordinate $`r`$. We first note that the 4-form field equations are identically satisfied without imposing any condition on the metric functions. It is also easy to see that $`G_{0i}=T_{0i}=0`$ in the orthonormal basis $`E^\mu =Ae^\mu `$, $`E^r=Bdr`$, $`E^a=Ce^a`$ and $`E^\alpha =De^\alpha `$. Therefore, the ansatz obeys the conditions (i) and (ii) of section 2. To find a supersymmetric solution of $`D=11`$ supergravity, one needs to impose conditions on metric functions which will ensure existence of at least one Killing spinor obeying $`D_\mu ϵ`$ $`=`$ $`_\mu ϵ+{\displaystyle \frac{A^{^{}}}{2AB}}\mathrm{\Gamma }_\mu {}_{}{}^{r}ϵ+{\displaystyle \frac{q_e}{18C^mD^n}}ϵ_{\nu \rho \sigma }\mathrm{\Gamma }^{\nu \rho \sigma r}\mathrm{\Gamma }_\mu ϵ=0,`$ (3.6) $`D_rϵ`$ $`=`$ $`_rϵ+{\displaystyle \frac{q_e}{18C^mD^n}}ϵ_{\nu \rho \sigma }\mathrm{\Gamma }^{\nu \rho \sigma }ϵ=0,`$ (3.7) $`D_aϵ`$ $`=`$ $`_aϵ+{\displaystyle \frac{C^{^{}}}{2CB}}\mathrm{\Gamma }_a{}_{}{}^{r}ϵ+{\displaystyle \frac{q_e}{36C^mD^n}}ϵ_{\nu \rho \sigma }\mathrm{\Gamma }^{\nu \rho \sigma r}\mathrm{\Gamma }_aϵ=0,`$ (3.8) $`D_\alpha ϵ`$ $`=`$ $`_\alpha ϵ+{\displaystyle \frac{D^{^{}}}{2DB}}\mathrm{\Gamma }_\alpha {}_{}{}^{r}ϵ+{\displaystyle \frac{q_e}{36C^mD^n}}ϵ_{\nu \rho \sigma }\mathrm{\Gamma }^{\nu \rho \sigma r}\mathrm{\Gamma }_\alpha ϵ=0,`$ (3.9) where the indices and spinors refer to the basis one-forms defined above and the electrical charge $`q_e`$ is defined to be the proportionality constant in (3.5) so that $$F_{\mu \nu \rho r}=\frac{q_e}{C^mD^n}ϵ_{\mu \nu \rho }.$$ (3.10) The presence of the functions $`C`$ and $`D`$ in (3.10) is due to the fact that $``$ in (3.5) refers to the line element (3.4). An easy way to solve these equations is to impose (3.1)-(3.3) and $$ϵ_{\mu \nu \rho }\mathrm{\Gamma }^{\mu \nu \rho }ϵ=6ϵ,$$ (3.11) which are of course consistent with each other. Then, the Killing spinor equations (3.6)-(3.9) imply $$\frac{A^{^{}}}{AB}=\frac{2q_e}{3C^mD^n},\frac{C^{^{}}}{CB}=\frac{q_e}{3C^mD^n},\frac{D^{^{}}}{DB}=\frac{q_e}{3C^mD^n}+\frac{1}{D}.$$ (3.12) Therefore, when (3.12) is satisfied, the background has at least one Killing spinor obeying (3.1),(3.3) and (3.11). Now, by the theorem of section 2, Einstein equations should be satisfied identically. At this point it is instructive to verify this claim by a direct calculation. The Einstein equations can be written as $$2\left(\frac{A^{^{}}}{AB}\right)^2+\left(\frac{A^{^{}}}{B}\right)^{^{}}\frac{1}{AB}+m\frac{A^{^{}}C^{^{}}}{AC}\frac{1}{B^2}+n\frac{A^{^{}}D^{^{}}}{AD}\frac{1}{B^2}=\frac{4q_e^2}{3C^{2m}D^{2n}},$$ (3.13) $$3\left(\frac{A^{^{}}}{B}\right)^{^{}}\frac{1}{AB}+m\left(\frac{C^{^{}}}{B}\right)^{^{}}\frac{1}{CB}+n\left(\frac{D^{^{}}}{B}\right)^{^{}}\frac{1}{DB}=\frac{4q_e^2}{3C^{2m}D^{2n}},$$ (3.14) $$(m1)\left(\frac{C^{^{}}}{CB}\right)^2+\left(\frac{C^{^{}}}{B}\right)^{^{}}\frac{1}{CB}+3\frac{A^{^{}}C^{^{}}}{AC}\frac{1}{B^2}+n\frac{C^{^{}}D^{^{}}}{CD}\frac{1}{B^2}=\frac{2q_e^2}{3C^{2m}D^{2n}},$$ (3.15) $$\frac{1}{D^2}+(n1)\left(\frac{D^{^{}}}{DB}\right)^2+\left(\frac{D^{^{}}}{B}\right)^{^{}}\frac{1}{DB}+3\frac{A^{^{}}D^{^{}}}{AD}\frac{1}{B^2}+m\frac{C^{^{}}D^{^{}}}{CD}\frac{1}{B^2}=\frac{2q_e^2}{3C^{2m}D^{2n}}.$$ (3.16) where we have used the fact that $`L_3`$ and $`M_m`$ are Ricci flat and $`X_n`$ is Einstein. It is now straightforward to check that, when the unknown functions $`A,B,C`$ and $`D`$ obey (3.12) they also solve the Einstein equations, as required by the theorem. There are four independent functions and three differential equations, which may be thought to imply that functions are not constrained enough. However, this is simply a manifestation of reparametrization invariance in coordinate $`r`$. We want to fix this reparametrization freedom in a way, which will help us in solving the differential equations. A convenient choice is $$C^mD^{n1}=r^{n1}+M,$$ (3.17) where $`M=2q_e/(n1)`$. Then, the unknown functions obeying the first order coupled equations (3.12) can be solved to give following solution $`ds^2`$ $`=`$ $`H^{2/3}ds_{L_3}^2+H^{1/3}(dr^2+ds_{M_m}^2+r^2ds_{X_n}^2),`$ (3.18) $`F_{\mu \nu \rho r}`$ $`=`$ $`{\displaystyle \frac{1}{2}}H^{7/6}(_rH)ϵ_{\mu \nu \rho },`$ (3.19) where $`H=(1+M/r^{n1})`$. ### M5-brane For the generalized M5-brane we write the following ansatz, $`ds^2`$ $`=`$ $`A^2ds_{L_6}^2+B^2dr^2+C^2ds_{M_m}^2+D^2ds_{X_n}^2,`$ (3.20) $`F`$ $``$ $`V_MV_X,`$ (3.21) where $`m+n=4`$ and the metric functions $`A,B,C`$ and $`D`$ are chosen to depend only on the coordinate $`r`$. As for the M2-brane, the 4-from field equations are satisfied identically, and $`G_{0i}=T_{0i}=0`$ in the orthonormal basis $`E^\mu =Ae^\mu `$, $`E^r=Bdr`$, $`E^a=Ce^a`$ and $`E^\alpha =De^\alpha `$. Therefore, the background satisfies conditions (i) and (ii) of section 2. To solve the Killing spinor equations, one can choose $`ϵ`$ to obey (3.1)-(3.3) and $$ϵ_{a..b}ϵ_{\alpha ..\beta }\mathrm{\Gamma }^{a..b\alpha ..\beta r}ϵ=m!n!ϵ.$$ (3.22) which are consistent with each other. One can then check that the Killing spinor equations imply $$\frac{A^{^{}}}{AB}=\frac{q_m}{3C^mD^n},\frac{C^{^{}}}{CB}=\frac{2q_m}{3C^mD^n},\frac{D^{^{}}}{DB}=\frac{2q_e}{3C^mD^n}+\frac{1}{D},$$ (3.23) where the magnetic charge $`q_m`$ is defined to be the proportionality constant in (3.21). Like in the membrane case, a convenient way to fix the $`r`$-reparametrization invariance is to define $`C^mD^{n1}=r^{n1}+M`$, where $`M=2q_m/(n1)`$. Then, (3.23) can be solved to give the following solution $`ds^2`$ $`=`$ $`H^{1/3}ds_{L_6}^2+H^{2/3}(dr^2+ds_{M_m}^2+r^2ds_{X_n}^2),`$ (3.24) $`F_{a..b\alpha ..\beta }`$ $`=`$ $`{\displaystyle \frac{1}{2}}H^{4/3}(_rH)ϵ_{a..b}ϵ_{\alpha ..\beta },`$ (3.25) where $`H=(1+M/r^{n1})`$. ### D3-brane The $`D3`$-brane solution of IIB supergravity in 10-dimensions has only non-vanishing (anti)-self dual 5-form and the metric. The equations governing dynamics of these fields can be written as $$R_{MN}=\frac{1}{96}F^{PQRS}{}_{M}{}^{}F_{PQRSN}^{}$$ (3.26) $$dF=0,F=F.$$ (3.27) The Killing spinors on such a background obey $$D_Mϵ=_Mϵ+\frac{i}{4\times 480}\mathrm{\Gamma }^{NP..Q}\mathrm{\Gamma }_MF_{NP..Q}ϵ=0,$$ (3.28) where $`ϵ`$ is a Weyl spinor $`\mathrm{\Gamma }_{(11)}ϵ=ϵ`$; $`\mathrm{\Gamma }_{(11)}=\mathrm{\Gamma }^0\mathrm{}\mathrm{\Gamma }^9`$, $`\mathrm{\Gamma }_{(11)}^{}=\mathrm{\Gamma }_{(11)}`$ and $`\mathrm{\Gamma }_{(11)}^2=I`$. For the generalized $`D3`$-brane we assume the form $`ds^2`$ $`=`$ $`A^2ds_{L_4}^2+B^2dr^2+C^2ds_{M_m}^2+D^2ds_{X_n}^2,`$ (3.29) $`F`$ $``$ $`(V_MV_X)(V_MV_X),`$ (3.30) where $`m+n=5`$, and the metric functions depend only on the coordinate $`r`$. The 5-form field equations are identically satisfied and $`G_{0i}=T_{0i}=0`$ in the basis $`E^\mu =Ae^\mu `$, $`E^r=Bdr`$, $`E^a=Ce^a`$ and $`E^\alpha =De^\alpha `$. Therefore, the ansatz obeys the conditions $`(i^{^{}})`$ and (ii) of section 2. On the other hand, if one choose $`ϵ`$ to obey (3.1)-(3.3) and impose further $$ϵ_{\mu \nu \rho \sigma }\mathrm{\Gamma }^{\mu \nu \rho \sigma }ϵ=24iϵ,$$ (3.31) the Killing spinor equations imply $$\frac{A^{^{}}}{AB}=\frac{q}{4C^mD^n},\frac{C^{^{}}}{CB}=\frac{q}{4C^mD^n},\frac{D^{^{}}}{DB}=\frac{q}{4C^mD^n}+\frac{1}{D},$$ (3.32) where the dyonic charge $`q`$ is defined to be the proportionality constant in (3.30). Therefore, when $`A,B,C`$ and $`D`$ obey (3.32), all conditions of the theorem are satisfied and the background should obey Einstein equations. Fixing $`r`$-reparametrization invariance by $`C^mD^{n1}=r^{n1}+M`$, we obtain the following solution, $`ds^2`$ $`=`$ $`H^{1/2}ds_{L_4}^2+H^{1/2}(dr^2+ds_{M_m}^2+r^2ds_{X_n}^2),`$ (3.33) $`F_{\mu \nu \rho \sigma r}`$ $`=`$ $`{\displaystyle \frac{1}{2}}H^{5/4}(_rH)ϵ_{\mu \nu \rho \sigma },`$ (3.34) where $`M=q/(n1)`$ and $`H=(1+M/r^{n1})`$. ### Interpretation and special cases It is clear from the structure of the antisymmetric tensor fields and Killing spinor projections that, in all solutions, $`L_d`$ represents the curved world-volumes of the branes. The dependence of the metric functions on the radial coordinate $`r`$ implies that $`M_m`$ and $`X_n`$ correspond to the smeared and actual transverse directions, respectively. The solutions obtained so far have a very similar structure with the well known brane solutions. Indeed, when $`X_n`$ is chosen to be the $`n`$-sphere, $`(r,X_n)`$ space becomes flat. In this case, $`H`$ can be generalized to be any harmonic function on this flat space. Choosing, furthermore, $`L_d`$ and $`M_m`$ to be flat gives the well known M2, M5 and D3-brane solutions which have $`m`$ smeared transverse directions. Therefore, one can think of the new solutions as the branes having curved world-volumes and smeared transverse directions, and non-spherical cross sections. For $`m=0`$, i.e. when $`M_m`$ is empty, one obtains the solutions of , which can thus be viewed to be the members of a more general family found in this paper. For $`n=1`$, reparametrization fixing conditions should clearly be modified. A convinient choice is to impose $`D=rC`$, which then gives the same solutions above with $`H=q\mathrm{log}r+const.`$ The number of Killing spinors on $`L_d`$, $`M_m`$ and $`X_n`$ determines the number of unbroken supersymmetries. Finally, it is also worth to mention that, field equations are still satisfied even when $`L_d`$, $`M_m`$ and $`X_n`$ have no Killing spinors. As we will show in section 5, this is not a coincidence, and indeed there is another simple way of generating new solutions. ## 4 Generalized intersections To generalize intersecting M2, M5 and D3-brane solutions, we will make use of Ricci flat Kähler spaces. In this section, we will still use the definitions of section 3, but furthermore assume that $`M_m`$ has a covariantly constant complex structure obeying $$J_a{}_{}{}^{b}J_{b}^{}{}_{}{}^{c}=\delta _a{}_{}{}^{c},$$ (4.1) $$_cJ_{ab}=0.$$ (4.2) This implies that $`M_m`$ is a Ricci flat Kähler space, $`m`$ is an even integer and $`J`$ is the Kähler two-form. Our basic strategy will still be the same; we will write an ansatz obeying conditions (i) and (ii) of section 2, and then work out the Killing spinor equations. ### M2-brane intersections We start with the following ansatz $`ds^2`$ $`=`$ $`A^2(dt^2)+B^2dr^2+C^2ds_{M_m}^2+D^2ds_{X_n}^2,`$ (4.3) $`F`$ $``$ $`(_MJ)V_X,`$ (4.4) where $`n+m=9`$ with $`n=1,3,5,7`$, $`_M`$ is the Hodge dual on the manifold $`M_m`$ and the metric functions are assumed to depend only on $`r`$. For $`n=7`$, the ansatz becomes a special case of the generalized M2-brane ansatz of section 3, and thus we will mainly consider $`n=1,3,5`$ cases. The 4-form field equations are identically satisfied since $`J`$ is both closed and co-closed on $`M_m`$. Furthermore, $`G_{0i}=T_{0i}=0`$ in the basis $`E^0=Adt`$, $`E^r=Bdr`$, $`E^a=Ce^a`$ and $`E^\alpha =De^\alpha `$. Therefore, the conditions (i) and (ii) of section 2 are satisfied. To solve the Killing spinor equations, $`ϵ`$ can consistently be chosen to obey (3.1), (3.2) and $$_tϵ=0,\mathrm{\Gamma }^0ϵ=iϵ,J_{ab}\mathrm{\Gamma }^bϵ=i\mathrm{\Gamma }_aϵ,$$ (4.5) which implies $$\frac{A^{^{}}}{AB}=\frac{(9n)q_e}{3C^{7n}D^n},\frac{C^{^{}}}{CB}=\frac{(n3)q_e}{6C^{7n}D^n},\frac{D^{^{}}}{DB}=\frac{(9n)q_e}{6C^{7n}D^n}+\frac{1}{D},$$ (4.6) where $`q_e`$ is the proportionality constant in (4.4). When $`n1`$, one can fix $`r`$-reparametrization invariance by imposing $`C^{m2}D^{n1}=r^{n1}+M`$ and this gives the following solution $`ds^2`$ $`=`$ $`H^{\frac{(n9)}{3}}(dt^2)+H^{\frac{(3n)}{6}}ds_{M_m}^2+H^{\frac{(9n)}{6}}(dr^2+r^2ds_{X_n}^2),`$ (4.7) $`F_{0rab}`$ $`=`$ $`{\displaystyle \frac{1}{2}}H^{\frac{(n21)}{12}}(_rH)J_{ab},`$ (4.8) where $`M=2q_e/(n1)`$ and $`H=(1+M/r^{n1})`$. For $`n=1`$, one should modify reparametrization fixing condition. A convinient choice is to demand $`D=rC^4`$, which then gives the above solution with $`H=2q_e\mathrm{log}r+const.`$ As it will be made more clear when we discuss the special cases, for $`n=5,3,1`$, the solutions describe two, three and four M2-brane intersections over a line, respectively. The structure of the background 4-form field implies that membranes also wrap over the 2-cycle dual to Kähler form of the space $`M_m`$. ### M5-brane intersections The discussion for intersecting M5-branes is not very different from M2-brane intersections. Considering the fact that the 4-form field should give rise to magnetic type of charges, we start with the following ansatz $`ds^2`$ $`=`$ $`A^2ds_{L_d}^2+B^2dr^2+C^2ds_{M_m}^2+D^2ds_{X_2}^2,`$ (4.9) $`F`$ $``$ $`JV_X,`$ (4.10) where $`m+d=8`$ with $`m=2,4,6`$ and the metric functions are assumed to depend only on $`r`$. For $`m=2,`$ this reduces to a special case of the generalized M5-brane ansatz of section 3. The 4-form field equations are identically satisfied since $`J`$ is both closed and co-closed on $`M_m`$. On the other hand $`G_{0i}=T_{0i}=0`$ in the basis $`E^0=Adt`$, $`E^r=Bdr`$, $`E^a=Ce^a`$, $`E^\alpha =De^\alpha `$ and therefore, the conditions (i) and (ii) of section 2 are satisfied. To solve the Killing spinor equations $`ϵ`$ can consistently be chosen to obey (3.1)-(3.3) and $$ϵ_{\alpha \beta }\mathrm{\Gamma }^{\alpha \beta r}ϵ=2iϵ,J_{ab}\mathrm{\Gamma }^bϵ=i\mathrm{\Gamma }_aϵ,$$ (4.11) which implies $$\frac{A^{^{}}}{AB}=\frac{mq_m}{6C^2D^2},\frac{C^{^{}}}{CB}=\frac{(m6)q_m}{6C^2D^2},\frac{D^{^{}}}{DB}=\frac{mq_m}{3C^2D^2}+\frac{1}{D},$$ (4.12) where $`q_m`$ is the proportionality constant in (4.4). To fix the $`r`$-reparametrization invariance we impose $`C^2D=r+M`$, where $`M=2q_m`$. Then, (4.12) can be solved to give $`ds^2`$ $`=`$ $`H^{m/6}ds_{L_{8m}}^2+H^{\frac{(6m)}{6}}ds_{M_m}^2+H^{m/3}(dr^2+r^2ds_{X_2}^2),`$ (4.13) $`F_{ab\alpha \beta }`$ $`=`$ $`{\displaystyle \frac{1}{2}}H^{\frac{(m+6)}{6}}(_rH)J_{ab}ϵ_{\alpha \beta },`$ (4.14) where $`H=(1+M/r)`$. As indicated above, for $`m=2`$ the solution becomes one of the M5-brane solution of section 3. On the other hand, for $`m=4`$ and $`m=6`$, the solutions describe two and three M5-branes intersecting over a three-brane and a string, respectively. The structure of the background 4-form field implies that M5-branes also wrap over the $`(m2)`$-cycle dual to $`_MJ`$. ### D3-brane intersections Let us start by discussing possible ansatzs for the (anti)-self dual 5-form involving the Kähler two-form of $`M_m`$. The first obvious choice is to assume $`FJV_{X_3}`$. Among the possible cases, $`m=2`$ corresponds to the smeared D3-brane of section 3 and $`m=6`$ is not allowed since this does not leave any room for a time-like direction. One may also try to write an ansatz involving $`_MJ`$ such as $`F_MJV_{X_n}`$. For $`m=4`$, this becomes equal to the choice $`FJV_{X_3}`$ and for $`m=6`$, $`X_n`$ space becomes one-dimensional. On the other hand, $`m=2`$ case turns out to be the same with the D3-brane ansatz of section 3. Summarizing, as a non-trivial ansatz representing intersecting D3-branes, one can write $`ds^2`$ $`=`$ $`A^2ds_{L_2}^2+B^2dr^2+C^2ds_{M_4}^2+D^2ds_{X_3}^2,`$ (4.15) $`F`$ $``$ $`(JV_X)(JV_X),`$ (4.16) where the metric functions are assumed to depend only on $`r`$. The 5-form field equations are identically satisfied and one can check that the background obeys condition (ii) of the section 2. To determine unknown functions, we solve the Killing spinor equations by imposing (3.1)-(3.3) and $$ϵ_{\mu \nu }\mathrm{\Gamma }^{\mu \nu }ϵ=2ϵ,J_{ab}\mathrm{\Gamma }^bϵ=i\mathrm{\Gamma }_aϵ,$$ (4.17) which implies $$\frac{A^{^{}}}{AB}=\frac{q}{2C^2D^3},\frac{C^{^{}}}{CB}=0,\frac{D^{^{}}}{DB}=\frac{q}{2C^2D^3}+\frac{1}{D},$$ (4.18) where $`q`$ is the proportionality constant in (4.16). Reparametrization invariance can be fixed by demanding $`C^2D^2=r^2+M`$, where $`M=q/2`$. Then, (4.18) can be solved to obtain the solution $`ds^2`$ $`=`$ $`H^1ds_{L_2}^2+ds_{M_4}^2+H(dr^2+r^2ds_{X_3}^2),`$ (4.19) $`F_{\mu \nu abr}`$ $`=`$ $`{\displaystyle \frac{1}{2}}H^{3/2}(_rH)J_{ab}ϵ_{\mu \nu },`$ (4.20) where $`H=(1+M/r^2)`$, which describes two D3-branes intersecting over a string and wrapping over the cycle dual to $`J`$. ### Interpretation and special cases To be able to interpret these solutions properly, let us choose $`X_n`$ to be the $`n`$-sphere. In this case, one can introduce Cartesian coordinates to span $`(r,X_n)`$ space and $`H`$ can be generalized to be any harmonic function of these coordinates. When $`L_d`$ and $`M_m`$ are also chosen to be flat, the solutions become the well known intersecting brane solutions in which $`all`$ harmonic functions are $`equal`$<sup>6</sup><sup>6</sup>6In the next section, we will illustrate with an example how to generalize intersecting brane solutions when harmonic functions are not equal.. Comparing with this special case, it is easy to argue that $`L_d`$, $`M_m`$ and $`X_n`$ correspond to common tangent, relative transverse and overall transverse directions. As mentioned earlier, the branes also wrap over the cycles dual to the Kähler two-form $`J`$ or its Hodge dual $`_MJ`$. For a given solution, this cycle can be written as a union of $`m/2`$ different submanifolds. To see this we note that, in an orthonormal basis, $`J`$ can be written as $`J=e^1e^2+..+e^{m1}e^m`$. Therefore, the cycle dual to $`J`$ or $`_MJ`$ becomes the sum of $`m/2`$ submanifolds dual to two-forms $`e^1e^2`$,..,$`e^{m1}e^m`$ or $`(m2)`$-forms $`_M(e^1e^2)`$, ..,$`_M(e^{m1}e^m)`$, respectively. On the other hand, the Killing spinor projections imply that there are $`m/2`$ intersecting branes each wrapping over one of these submanifolds. We note that this interpretation is consistent with the special case where $`M_m`$ is flat. When $`X_n`$ is different from the $`n`$-sphere, one obtains new solutions which have not been encountered before. These solutions can be viewed to be the intersecting brane counterparts of the brane solutions constructed in , which have non-spherical transverse spaces. In all cases the number of unbroken supersymmetries depend on the number of Killing spinors on $`L_d`$, $`M_m`$ and $`X_n`$. Like for the generalized brane solutions of the previous section, one can check that the field equations are still satisfied, even when the spaces $`L_d`$, $`M_m`$ and $`X_n`$ have no Killing spinors. ## 5 A general argument In the last two sections, we have obtained supersymmetric, generalized brane solutions by writing suitable ansatzs and working out Killing spinor equations. As pointed earlier, the field equations are still satisfied, even when the manifolds $`L_d`$, $`M_m`$ and $`X_n`$ have no Killing spinors. This suggests existence of a general relation which may hold at the level of field equations. Let us consider a brane solution which has a metric of the form, $$ds^2=A^2ds_M^2+ds_N^2,$$ (5.1) where $`ds_M^2`$ and $`ds_N^2`$ are line elements of $`m`$,$`n`$ dimensional manifolds $`M`$ and $`N`$, respectively, and $`A`$ is a function on $`N`$. We assume that the anti-symmetric tensor fields of the solution are of the form $`FV_M\mathrm{}`$ or $`F\mathrm{}`$, where $`V_M`$ is the volume form of $`M`$ and the dotted terms depend only on $`N`$. Furthermore, we consider the cases in which scalar fields are independent of $`M`$ and the anti-symmetric tensor field equations reduce to $`dF=0`$ and $`dF=0`$ <sup>7</sup><sup>7</sup>7It seems one can relax the last assumption, but here, for simplicity, we do not consider more general cases.. Our claim is that, for such a brane solution, the field equations are still satisfied when $`M`$ is replaced with a different manifold $`\stackrel{~}{M}`$, provided the Ricci tensors of both manifolds have the same form. For instance, if $`M`$ is flat, $`\stackrel{~}{M}`$ can be any manifold which is Ricci flat. Or, if $`M`$ is a sphere, $`\stackrel{~}{M}`$ can be any Einstein manifold. To prove this, let $`Ae^a`$ and $`e^\alpha `$ be a basis for the tangent space, where $`e^a`$ and $`e^\alpha `$ are the basis-one forms of $`M`$ and $`N`$. We calculate the Ricci tensor of (5.1) as $`R_a^b`$ $`=`$ $`{\displaystyle \frac{1}{A^2}}R_{(M)a}{}_{}{}^{b}(m1){\displaystyle \frac{A_\alpha A^\alpha }{A^2}}\delta _a{}_{}{}^{b}{\displaystyle \frac{F_\alpha ^\alpha }{A}}\delta _a{}_{}{}^{b},`$ (5.2) $`R_\alpha ^\beta `$ $`=`$ $`R_{(N)\alpha }{}_{}{}^{\beta }n{\displaystyle \frac{F_\alpha ^\beta }{A}},`$ (5.3) where the indices refer to the tangent space and $`R_{(M)ab}`$, $`R_{(N)\alpha \beta }`$ are the Ricci tensors of $`M`$ and $`N`$, respectively. The tensor quantities $`A_\alpha `$ and $`F_{\alpha \beta }`$ are defined by the relations $$dA=A_\alpha e^\alpha ,$$ (5.4) $$dA_\alpha =F_{\alpha \beta }e^\beta .$$ (5.5) We note that the Ricci tensor of (5.1) depends only on the dimension $`m`$ and Ricci tensor $`R_{(M)ab}`$ of the manifold $`M`$. Let us now analyze how the field equations may change when one replaces $`M`$ with $`\stackrel{~}{M}`$. We first note that $`F`$ is closed or co-closed irrespective of the choice of $`M`$, thus the form equations are still satisfied. The possible terms in scalar and Einstein equations have at most second order covariant derivatives of scalars. One can check that when $`M`$ is replaced with $`\stackrel{~}{M}`$, only $`\omega ^a_{bc}`$ components of the spin connection may change. This ensures that the second order covariant derivatives of scalars remain the same, since they are assumed to be independent of $`M`$. Referring to the tangent space, it is easy to show that any tensor field which is constructed from $`F`$, the terms containing up to second order covariant derivatives of scalars and the metric have the same form. By (5.2) and (5.3), the last statement ensures that both scalar and Einstein equations are still satisfied. In some cases, $`M`$ can be a Kähler space and the anti-symmetric tensor fields may be of the form $`FJ..`$, where $`J`$ is the Kähler two-form. In such cases, $`\stackrel{~}{M}`$ should also be chosen to be Kähler and $`J`$ should be replaced with the Kähler two-form of $`\stackrel{~}{M}`$. One can easily argue that the field equations are not affected from these replacements. Let us now present some examples which enables one to recover some solutions obtained in the literature and in the previous sections. Consider, for instance, the well known smeared M2-brane solution, $`ds^2`$ $`=`$ $`H^{2/3}(dt^2+dx_1^2+dx_2^2)+H^{1/3}(dy_1^2+..+dy_m^2+dy_{m+1}^2+..+dy_{10}^2),`$ (5.6) $`F_{\mu \nu \rho \alpha }`$ $`=`$ $`{\displaystyle \frac{1}{2}}H^{7/6}(_\alpha H)ϵ_{\mu \nu \rho },`$ (5.7) where $`H`$ is harmonic on $`(y_{m+1},..,y_{10})`$. Referring to the above discussion, it is easy to see that the world-volume directions $`(t,x_1,x_2)`$ and the smeared transverse directions $`(y_1,..,y_m)`$ can be replaced with more general Ricci flat manifolds. In this way, one can obtain $`ds^2`$ $`=`$ $`H^{2/3}ds_{L_3}^2+H^{1/3}(ds_{M_m}^2+dy_{m+1}^2+..+dy_{10}^2),`$ (5.8) $`F_{\mu \nu \rho \alpha }`$ $`=`$ $`{\displaystyle \frac{1}{2}}H^{7/6}(_\alpha H)ϵ_{\mu \nu \rho },`$ (5.9) which becomes one of the generalized M2-brane solutions constructed in section 3. In IIB theory, one can start from the single centered D3-brane solution $`ds^2`$ $`=`$ $`H^{1/2}(dt^2+dx_1^2+..+dx_3^2)+H^{1/2}(dr^2+r^2d\mathrm{\Omega }_5^2),`$ (5.10) $`F_{\mu \nu \rho \sigma r}`$ $`=`$ $`{\displaystyle \frac{1}{2}}H^{5/4}(_rH)ϵ_{\mu \nu \rho \sigma },`$ (5.11) where $`d\mathrm{\Omega }_5^2`$ is the line element on the $`5`$-sphere and $`H=(1+M/r^4)`$. The metric and the 5-form field are of the form discussed above with the flat world-volume $`(t,x_1,..,x_3)`$ and the transverse $`5`$-sphere play the role of manifold $`M`$ in (5.1). By the claim we have just proved, one can generalize the well known solution as $`ds^2`$ $`=`$ $`H^{1/2}ds_{L_4}^2+H^{1/2}(dr^2+r^2ds_{X_5}^2),`$ (5.12) $`F_{\mu \nu \rho \sigma r}`$ $`=`$ $`{\displaystyle \frac{1}{2}}H^{5/4}(_rH)ϵ_{\mu \nu \rho \sigma },`$ (5.13) which corresponds to the solution obtained in when $`L_4`$ is flat. Let us finally consider an example to illustrate how the intersecting solutions can be generalized when the harmonic functions are not chosen to be equal. The background representing two M5-branes intersecting over a three-brane is given by $`ds^2`$ $`=`$ $`(H_1H_2)^{1/2}(dt^2+dx_1^2+dx_2^2+dx_3^2)+H_1^{1/3}H_2^{2/3}(dx_4^2+dx_5^2)`$ (5.14) $`+`$ $`H_1^{2/3}H_2^{1/3}(dx_6^2+dx_7^2)+(H_1H_2)^{2/3}(dx_8^2+..+dx_{10}^2),`$ $`F`$ $``$ $`dx_4dx_5dH_2+dx_6dx_7dH_1,`$ (5.15) where, $`H_1`$ and $`H_2`$ are harmonic and $``$ is the Hodge dual on $`(x_8,..,x_{10})`$ space. Here, it is easy to see that, $`(t,x_1,x_2,x_3)`$, $`(x_4,x_5)`$ and $`(x_6,x_7)`$ spaces play the role of manifold $`M`$ in (5.1). Since all these spaces are flat, one can replace them with arbitrary Ricci flat manifolds $`L_4`$, $`M_2`$ and $`\stackrel{~}{M}_2`$ to obtain <sup>8</sup><sup>8</sup>8In two-dimensions, Ricci flatness implies flatness. Here, we would like to emphasize that for field equations to be satisfied Ricci flatness of these manifolds is sufficient. $`ds^2`$ $`=`$ $`(H_1H_2)^{1/3}ds_{L_4}^2+H_1^{1/3}H_2^{2/3}ds_{M_2}^2+H_1^{2/3}H_2^{1/3}ds_{\stackrel{~}{M}_2}^2`$ (5.16) $`+`$ $`(H_1H_2)^{2/3}(dx_5^2+..+dx_{10}^2),`$ $`F`$ $``$ $`V_{M_2}dH_2+V_{\stackrel{~}{M}_2}dH_1.`$ (5.17) When $`H_1`$=$`H_2=H`$, the 4-form field (5.15) can be written as $$FJdH,$$ (5.18) where $`J`$ is the complex structure of the flat space $`(x_4,x_5,x_6,x_7)`$. In this case, these four flat directions can be replaced with any 4-dimensional Ricci flat Kähler space, which corresponds to one class of intersections found in section 4. As can be inferred from these examples, in a brane solution the flat world-volume and smeared transverse directions, and the transverse sphere at a fixed radial distance can be replaced with more general Ricci flat and Einstein manifolds. In intersecting brane solutions the common tangent and relative transverse directions have this property. We note that, these replacements can be done at the level of field equations and supersymmetry of the new solution obtained in this way is not manifest. ## 6 Solutions from $`U(1)`$ bundles over Ricci flat Kähler spaces In the brane solutions obtained in section 3 the warping factors are the same for the transverse and world-volume directions. In this section, by using the theorem of section 2, we will take a step in obtaining solutions which fail to have this property, at least along the transverse directions. We first construct singular, Ricci flat manifolds having a cone-like structure over U(1) bundles over Ricci flat Kähler spaces, which give rise to new supersymmetric vacua for non-gauged supergravities. Since the conformal factors multiplying U(1) fibers and the base spaces turn out to be neither equal to each other nor equal to the square of the coordinate parametrizing the cone, we name these spaces as (generalized) cones. As we will see in a moment, there are two singularities associated with the cone, since at the origin the base space and at infinity the U(1) fibers shrink to zero size. Also, in viewing the solutions as Ricci flat compactifications, unlike the conventional Kaluza-Klein picture, there is no natural way of assuming the internal spaces to be small compared to the space-time. At the end of this section, we will construct brane solutions which asymptotically approach these singular compactifications and preserve half of the available supersymmtries. This is in fact not surprising, if one assumes stability of the vacua, since the fundamental branes underlying the supergravity theories reveals themselves as supersymmetric soliton solutions. ### A singular Ricci flat space It is well known that the cone over a manifold $`X`$ $$ds^2=dr^2+r^2ds_X^2$$ (6.1) is Ricci flat if and only if $`X`$ is Einstein. In this section, we would like to construct Ricci flat manifolds having a cone-like structure over Ricci flat Kähler spaces. It is clear that one should modify (6.1) in a non-trivial way. Let us consider a metric of the form $$ds^2=B^2dr^2+C^2ds_{M_m}^2+D^2(d\tau 𝒜)^2,$$ (6.2) where $`M_m`$ is an $`m`$-dimensional Ricci flat Kähler space, $`𝒜`$ is the one-form potential for the Kähler two-form so that $`d𝒜=J`$, $`\tau `$ is a periodic coordinate and the metric functions are assumed to depend only on $`r`$. It is clear that we have a fiber bundle structure and $`\tau `$ is the coordinate on U(1) fibers. We would like to determine $`B`$,$`C`$ and $`D`$ which give rise to a Ricci flat space. Instead of calculating Ricci tensor and solving second order, coupled differential equations we demand existence of a covariantly constant spinor which, by considerations of section 2, will imply Ricci flatness. In the tangent space basis $`E^r=Bdr`$, $`E^a=Ce^a`$ and $`E^D=D(d\tau 𝒜)`$, a covariantly constant spinor obeys $`Dϵ`$ $`=`$ $`dϵ+{\displaystyle \frac{1}{4}}(\omega _{ab}+{\displaystyle \frac{D}{2C^2}}J_{ab}E^D)\mathrm{\Gamma }^{ab}ϵ{\displaystyle \frac{1}{4}}{\displaystyle \frac{D}{C^2}}J_{ab}E^b\mathrm{\Gamma }^{Da}ϵ`$ (6.3) $`+`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{C^{^{}}}{CB}}E_a\mathrm{\Gamma }^{ar}ϵ+{\displaystyle \frac{1}{2}}{\displaystyle \frac{D^{^{}}}{DB}}E_D\mathrm{\Gamma }^{Dr}ϵ=0,`$ where $`\omega _{ab}`$ is the spin connection on $`M_m`$, and we have used the differential form notation. Conistently imposing $$_rϵ=_\tau ϵ=0,dϵ+\frac{1}{4}\omega _{ab}\mathrm{\Gamma }^{ab}ϵ=0,$$ (6.4) $$J_{ab}\mathrm{\Gamma }^bϵ=i\mathrm{\Gamma }_aϵ,\mathrm{\Gamma }^{Dr}ϵ=iϵ,$$ (6.5) one obtains $$\frac{C^{^{}}}{CB}=\frac{D}{2C^2},\frac{D^{^{}}}{DB}=\frac{mD}{4C^2}.$$ (6.6) Therefore, when the metric functions $`B`$, $`C`$ and $`D`$ obey (6.6), the Ricci tensor of (6.2) vanishes. Fixing $`r`$-reparametrization invariance by $`B=4/(m+4)`$, (6.6) can be solved to give the following Ricci flat metric $$ds^2=\frac{16}{(m+4)^2}dr^2+r^{4/(m+4)}ds_{M_m}^2+r^{2m/(4+m)}(d\tau 𝒜)^2,$$ (6.7) which we name as the space $`M_C`$. Topologically, $`M_C`$ is a product of a line parametrized by $`r`$ and a U(1) bundle over a Ricci flat Kähler space. The metric is singular, since at the origin, when $`r=0`$, the base space and at infinity, as $`r\mathrm{}`$, the U(1) fibers shrink to zero size. On the other hand, $`M_C`$ have covariantly constant spinors, and thus gives rise to supersymmetric compactifications of non-gauged supergravities. In these compactifications, one can view the coordinate $`r`$ as a radial coordinate in the uncompactified space-time. Therefore, the U(1) bundle over $`M_m`$ plays the role of internal space. As can be inferred from the structure of the metric, there is no natural way of assuming the internal directions to be “small”. Thus these compactifications are very different from the conventional Kaluza-Klein ones. Although the metric is singular, the existence of unbroken supersymmetries is a sign for stability. We now construct M2 and D3-brane solutions which belong to these singular vacua. ### M2 and D3-branes on $`M_C`$ compactifications In searching brane solutions which asymptotically belong to above Ricci flat compactifications, we consider the cases in which $`M_C`$ plays the role of the total transverse space. Since $`M_C`$ is an even dimensional manifold, this leaves the possibility of an even dimensional brane in $`D=11`$, and an odd dimensional brane in $`D=10`$. We think of the coordinate $`r`$ as the radial coordinate in the transverse space. M2-brane is the only even dimensional brane in $`D=11`$, and this suggests the following ansatz $`ds^2`$ $`=`$ $`A^2ds_{L_3}^2+B^2dr^2+C^2ds_{M_6}^2+D^2(d\tau 𝒜)^2,`$ (6.8) $`F`$ $``$ $`V_M(d\tau 𝒜),`$ (6.9) where all metric functions are assumed to depend on $`r`$. We remind the reader that we are still using the definitions of section 3. One can check that, the 4-form field equations are identically satisfied, and $`G_{0i}=T_{0i}=0`$ in the basis $`E^\mu =Ae^\mu `$, $`E^r=Bdr`$, $`E^a=Ce^a`$ and $`E^D=D(d\tau 𝒜)`$. Therefore, the conditions (i) and (ii) of section 2 are satisfied. By imposing <sup>9</sup><sup>9</sup>9The covariant derivatives in (6.10) refer to the spaces $`L_3`$ and $`M_6`$. $$_\tau ϵ=0,_\mu ϵ=0,_aϵ=0,$$ (6.10) and $$ϵ_{\mu \nu \rho }\mathrm{\Gamma }^{\mu \nu \rho }ϵ=6ϵ,J_{ab}\mathrm{\Gamma }^bϵ=i\mathrm{\Gamma }_aϵ,\mathrm{\Gamma }^{Dr}ϵ=iϵ,$$ (6.11) the Killing spinor equations imply $$\frac{A^{^{}}}{AB}=\frac{2q_e}{3C^6D},\frac{C^{^{}}}{CB}=\frac{D}{2C^2}\frac{q_e}{3C^6D},\frac{D^{^{}}}{DB}=\frac{3D}{2C^2}\frac{q_e}{3C^6D},$$ (6.12) where $`q_e`$ is the proportionality constant in (6.9). Since the conditions imposed on the Killing spinor are consistent with each other, when (6.12) is satisfied, the background should obey the Einstein equations by the theorem of section 2. We fix the $`r`$-reparametrization invariance by imposing $`C^4D^2=r`$. Using this condition, (6.12) can be solved to give the following solution $`ds^2`$ $`=`$ $`\left(1{\displaystyle \frac{M}{r}}\right)^{2/3}ds_{L_3}^2+{\displaystyle \frac{1}{r^7}}\left(1{\displaystyle \frac{M}{r}}\right)^{22/3}dr^2+{\displaystyle \frac{1}{r}}\left(1{\displaystyle \frac{M}{r}}\right)^{4/3}ds_{M_6}^2`$ (6.13) $`+`$ $`r^3\left(1{\displaystyle \frac{M}{r}}\right)^{8/3}(d\tau 𝒜)^2,`$ $`F_{\mu \nu \rho r}`$ $`=`$ $`{\displaystyle \frac{M}{2}}r^{3/2}\left(1{\displaystyle \frac{M}{r}}\right)^{8/3}ϵ_{\mu \nu \rho },`$ (6.14) where $`M=2q_e`$. It is clear that, in the solution, $`L_3`$ represents the world-volume of the M2-brane and $`r`$ is a radial coordinate along transverse directions. Unlike all solutions obtained before, there are three different <sup>10</sup><sup>10</sup>10Indeed, by a coordinate change, the warping factor along the coordinate $`r`$ can be made equal, for instnce, to the warping factor along U(1) fiber. warping factors multiplying the transverse directions, which have also a non-trivial topological structure due to the presence of U(1) bundle. There is a horizon located at $`r=M`$ and asymptotically, as $`r\mathrm{}`$, the solution becomes $`L_3\times M_C`$, which can be seen by a coordinate change $`d\stackrel{~}{r}r^{7/2}dr`$, near infinity. The solution is a half supersymmetry preserving state of the vacuum $`L_3\times M_C`$ since the presence of the M2-brane brakes only half of the available supersymmetries. As an example in $`D=10`$, we consider the D3-brane of IIB theory and start with the following ansatz $`ds^2`$ $`=`$ $`A^2ds_{L_4}^2+B^2dr^2+C^2ds_{M_4}^2+D^2(d\tau 𝒜)^2,`$ (6.15) $`F`$ $``$ $`V_M(d\tau 𝒜)[V_M(d\tau 𝒜)],`$ (6.16) where, as usual, we assume that $`A`$, $`B`$, $`C`$ and $`D`$ depend only on $`r`$. One can check that the ansatz obeys the condition $`(i^{^{}})`$ and (ii) of section 2, and therefore to obtain a supersymmetric solution one needs to work out Killing spinor equations. Imposing <sup>11</sup><sup>11</sup>11The covariant derivatives in (6.17) refer to the spaces $`L_4`$ and $`M_4`$. $$_\tau ϵ=0,_\mu ϵ=0,_aϵ=0,$$ (6.17) and $$ϵ_{\mu \nu \rho \sigma }\mathrm{\Gamma }^{\mu \nu \rho \sigma }ϵ=24iϵ,J_{ab}\mathrm{\Gamma }^bϵ=i\mathrm{\Gamma }_aϵ,\mathrm{\Gamma }^{Dr}ϵ=iϵ,$$ (6.18) the Killing spinor equations imply $$\frac{A^{^{}}}{AB}=\frac{q}{4C^4D},\frac{C^{^{}}}{CB}=\frac{D}{2C^2}\frac{q}{4C^4D},\frac{D^{^{}}}{DB}=\frac{D}{C^2}\frac{q}{4C^4D},$$ (6.19) where the dyonic charge $`q`$ is defined to be the proportionality constant in (6.16). Fixing $`r`$-reparametrization invariance by $`C^2D^2=r`$, one can solve the differential equations to obtain the following D3-brane solution $`ds^2`$ $`=`$ $`\left(1{\displaystyle \frac{q}{r}}\right)^{1/2}ds_{L_4}^2+{\displaystyle \frac{1}{r^6}}\left(1{\displaystyle \frac{q}{r}}\right)^{13/2}dr^2+{\displaystyle \frac{1}{r}}\left(1{\displaystyle \frac{q}{r}}\right)^{3/2}ds_{M_4}^2`$ (6.20) $`+`$ $`r^2\left(1{\displaystyle \frac{q}{r}}\right)^{3/2}(d\tau 𝒜)^2,`$ $`F_{\mu \nu \rho \sigma r}`$ $`=`$ $`qr\left(1{\displaystyle \frac{q}{r}}\right)^{9/4}ϵ_{\mu \nu \rho \sigma }.`$ (6.21) It is clear that, $`L_4`$ corresponds to the world-volume of the D3-brane. Asymptotically, as $`r\mathrm{}`$, the solution becomes $`L_4\times M_C`$ which can be seen by a coordinate change $`d\stackrel{~}{r}r^3dr`$, near infinity. The solution is half supersymmetry preserving state of the vacuum $`L_4\times M_C`$, since the presence of D3-brane brakes half of the available supersymmetries. There is also a horizon located at $`r=q`$. We finally note the non-orthogonal decomposition of the transverse directions due to the U(1) bundle structure. In this paper we do not attempt to determine singularity structures of these solutions. Due to the presence of unbroken supersymmetries we believe that the solutions are stable. Another important open problem is to find realizations of unbroken supersymmetries on the solutions, which enable one to determine how Bogomolny bounds are saturated. ## 7 Conclusions The brane solutions of supergravity theories have played a crucial role in recent developments in the non-perturbative string/M theories. With this experience, it is reasonable to claim that finding new solutions will also be important for future developments. In this paper, we have constructed new classes of solutions in $`D=11`$ and $`D=10`$ dimensions. The solutions obtained in sections 3 and 4 have a very similar structure with the well known solutions and can be viewed as the generalizations of them. Solutions obtained in section 6 belong to a different class in which the transverse directions do not have a single warping factor. In section 5, we have presented a general argument which allows one to construct new solutions from the old ones by replacing certain directions with more general manifolds. In constructing new solutions, we have mainly used the theorem proved in . The results of the present paper show that the theorem reviewed in section 2 is an important way of obtaining new supersymmetric solutions. Here, we would like to mention two open problems which can possibly be attacked by using the theorem; the first one is to find non-static cases which may lead to time dependent or stationary supersymmetric solutions, and the second one is to construct explicit examples of non-trivially embedded brane solutions, as formally discussed in . By the well known solution generating techniques, like applying S or T-dualities, or by dimensional reduction, one can obtain more new solutions in $`D=10`$ or in lower dimensions. As an interesting application of this, we note that it is possible to untwist a U(1) bundle by a T-duality transformation along the coordinate parametrizing U(1) fiber. As for the usual brane solutions, one may try to identify low energy theories defined on the world-volumes. The number and the supermultiplet structure of collective coordinates depend on the realizations of unbroken supersymmetries, the singularities of the solutions and the isometries of the internal manifolds. A classification of spaces having Killing spinors in diverse dimensions is the key ingredient for such an analysis. After identifying the low energy field theories, the next important question is to learn how to take decoupling or near horizon limits. We note that, brane solutions having transverse Einstein spaces give rise to generalizations of the original AdS/CFT dualities which correspond to compactifications on arbitrary Einstein manifolds. By defining sensible decoupling limits for other type of solutions found in this paper, one may discover interesting realizations of holographic principle.
warning/0004/math0004171.html	ar5iv	text	# Combinatorics and Quotients of Toric Varieties ## 0. Introduction This paper studies two related subjects. One is some combinatorics arising from linear projections of polytopes and fans of cones. The other is quotient varieties of toric varieties. The relation is that projections of polytopes are related to quotients of projective toric varieties and projection of fans are related to quotients of general toric varieties. Despite its relation to geometry the first part is purely combinatorial and should be of interest in its own right. For the combinatorial part, consider a linear projection $`\pi :VW`$ between two real vector spaces. Let $`P`$ be a polytope in $`V`$ of full dimension and $`Q=\pi (P)`$. The projections of all faces of $`P`$ under $`\pi `$ induce a polytopal subdivision of $`Q`$. This gives rise to the poset $`\mathrm{\Gamma }=\mathrm{\Gamma }(P,Q)`$ of the so-called chambers in $`Q`$ (Definition 1.1). The sequence $`0\mathrm{ker}\pi V\mathrm{@}>\pi >>W0`$ induces the dual sequence $$0W^{}V^{}\mathrm{@}>\pi ^{}>>(\mathrm{ker}\pi )^{}0.$$ The projections of all cones of the normal fan $`\mathrm{\Delta }(P)`$ of $`P`$ under $`\pi ^{}`$ induce a fan in $`(\mathrm{ker}\pi )^{}`$. This is the normal fan $`\mathrm{\Delta }(\mathrm{\Sigma }(P,Q))`$ of the Billera-Sturmfels fiber polytope $`\mathrm{\Sigma }(P,Q)`$. Let $`\mathrm{\Gamma }^{}=\mathrm{\Gamma }^{}(P,Q)`$ be the poset of cones in this fan (Definition 1.2). The above descriptions clearly put $`\mathrm{\Gamma }`$ and $`\mathrm{\Gamma }^{}`$ in a linear dual situation. Indeed, they are part of a larger picture which goes beyond just dual analogy. First, it has been known that $`\mathrm{\Gamma }^{}`$ is in one-to-one correspondence with the poset $`𝒯_{coh}=𝒯_{coh}(P,Q)`$ of the so-called coherent strings of the projection $`\pi :PQ`$ (Definition 2.1). The poset $`𝒯=𝒯(P,Q)`$ of locally coherent strings, the generalizations of coherent strings, has also been around for quite a while (Definition 3.1). Investigating into this, it is quite natural to ask the following questions: 1. what are the generalizations of elements in $`\mathrm{\Gamma }^{}`$ so that they correspond to locally coherent strings in $`𝒯`$? 2. are there coherent costrings linear dual to coherent strings? 3. are there locally coherent costrings linear dual to locally coherent strings? 4. are there perfect correspondences for these new objects similar to the correspondence between $`\mathrm{\Gamma }^{}`$ and $`𝒯_{coh}`$? Indeed, perfect correspondences do occur and are the main theme of this paper. To this end, some new objects are to be introduced in response to each and every of the above questions: 1. the poset $`\mathrm{\Gamma }_{vir}^{}=\mathrm{\Gamma }_{vir}^{}(P,Q)`$ of virtual cones of the projection $`\pi ^{}`$. A virtual cone is a combinatorial abstraction that generalizes a real polyhedral cone (Definition 7.1). It corresponds to a locally coherent string of $`\pi :PQ`$. 2. the poset $`𝒯_{coh}^{}=𝒯_{coh}^{}(P,Q)`$ of coherent costrings of the projection $`\pi ^{}`$. This notion is linear dual to that of coherent strings. A coherent costring is roughly speaking a special lifting of a fan via the map $`\pi ^{}`$ (Definition 4.2). It corresponds to a chamber in $`\mathrm{\Gamma }`$. 3. the poset of $`𝒯^{}=𝒯^{}(P,Q)`$ of locally coherent costrings of the projection $`\pi ^{}`$. This notion is linear dual to that of locally coherent strings. A locally coherent costring is roughly speaking a lifting of a fan via the map $`\pi ^{}`$ (Definition 5.1). 4. the poset $`\mathrm{\Gamma }_{vir}=\mathrm{\Gamma }_{vir}(P,Q)`$ of virtual chambers (cells) of the projection $`\pi `$. A virtual chamber (cell) is a combinatorial abstraction that generalizes a real chamber (cell) in $`\mathrm{\Gamma }`$ (Definition 6.1). It corresponds to a locally coherent costring. Locally coherent strings first appeared in when $`dimQ=1`$ and were generalized and studied afterwards, especially in association with the generalized Baues conjecture (see for a recent survey). Virtual chambers for a simplex was introduced and studied in . Our first main result in this paper is ###### Theorem A (Theorem 8.1.) Let the notations be as explained in the above. Then 1. $`𝒯(P,Q)`$ is canonically anti-isomorphic to $`\mathrm{\Gamma }_{vir}^{}(P,Q)`$ which extends the anti-isomorphism between $`𝒯_{coh}(P,Q)`$ and $`\mathrm{\Gamma }^{}(P,Q)`$; 2. $`𝒯^{}(P,Q)`$ is canonically anti-isomorphic to $`\mathrm{\Gamma }_{vir}(P,Q)`$ which extends the anti-isomorphism between $`𝒯_{coh}^{}(P,Q)`$ and $`\mathrm{\Gamma }(P,Q)`$. As hinted in the beginning, the above combinatorial objects underlie some quotient toric geometry. We now elaborate. Let $`X_{\mathrm{\Delta }_0}=_{\sigma \mathrm{\Delta }_0}A_\sigma `$ be the toric variety over the complex number field defined by a fan $`\mathrm{\Delta }_0`$ with a lattice $`N`$, where $`A_\sigma `$ is the affine variety defined by the cone $`\sigma `$ in $`\mathrm{\Delta }_0`$. Let $`M^2`$ be any sublattice of the dual lattice $`M`$ of $`N`$. Consider the exact sequence $$0M^2M\mathrm{@}>\pi >>M^10.$$ If we apply $`\mathrm{Hom}_{}(,^{})`$ to it, then we get the exact sequence $$1\mathrm{Hom}_{}(M^1,^{})\mathrm{Hom}_{}(M,^{})\mathrm{Hom}_{}(M^2,^{})1.$$ The group $`G=\mathrm{Hom}_{}(M^1,^{})`$ is the product of the torus $`\mathrm{Hom}_{}(M_{free}^1,^{})`$ and the finite group $`\mathrm{Hom}_{}(M_{tor}^1,^{})`$. Let $$0N^1N\mathrm{@}>\pi ^{}>>N^20$$ be the dual sequence. The projective quotients of projective toric varieties by subtorus have been studied by Kapranov, Sturmfels and Zelevinsky in . See also, . For the arbitrary toric variety $`X_{\mathrm{\Delta }_0}`$ (not necessarily complete, not necessarily quasi-projective), we have the following ###### Theorem B (Theorem 11.5) Let $`\mathrm{\Delta }\mathrm{\Delta }_0`$ be a subset of cones and $`U=_{\sigma \mathrm{\Delta }}A_\sigma `$ be an open subset of $`X_{\mathrm{\Delta }_0}`$ where $`A_\sigma `$ is the affine open subset defined by $`\sigma `$. Then 1. $`U`$ has a categorical quotient by the action of $`G`$ if $`\mathrm{\Delta }`$ is a locally coherent costring of the projection $`\pi ^{}`$ and every cone in $`\pi ^{}(\mathrm{\Delta })`$ is strongly convex. In this case the quotient $`U//G`$ is isomorphic to $`X_{\pi ^{}(\mathrm{\Delta })}`$, the toric variety defined by the induced fan $`\pi ^{}(\mathrm{\Delta })`$; 2. $`U//G`$ is geometric if and only if $`dim\pi ^{}(\sigma )=dim\sigma `$ for all $`\sigma \mathrm{\Delta }`$ (i.e., $`\mathrm{\Delta }`$ is tight by definition, or equivalently, minimal in $`𝒯^{}`$). When $`X_{\mathrm{\Delta }_0}`$ is projective and equipped with a $`T`$-linearized ample line bundle $`L`$, there corresponds to a polytope $`P`$ in the lattice $`M`$. Let $`Q=\pi (P)`$. Then by the equivalence of (2) in Theorem A, we have ###### Corollary C Every virtual chamber (cell) in $`\mathrm{\Gamma }_{vir}(P,Q)`$ defines a (not-necessarily projective) toric quotient variety of $`X_{\mathrm{\Delta }_0}`$ whose fan is induced from the corresponding locally coherent costring in $`𝒯^{}(P,Q)`$. Note that the real chambers (cells) correspond to projective toric quotients whose fans are those given by global coherent costrings. In the end of this paper, we try to relate the combinatorial bistellar flips to the geometric Mori-type flips among quotients of affine spaces. Also we introduce a new notion, virtual (or pseudo) oriented matroid, which is somewhat related to the rest of the paper. Finally, we mention that there are some recent works that study the quotients of toric varieties by subtori , , and . Note that the group $`G`$ that we use to take quotient is typically disconnected. Also, I learned first from Reiner that Santos () has obtained some important results on Baues conjectures which are related to our combinatorial constructions. Acknowledgments. My thanks go to Bernd Sturmfels for suggesting that I investigate the perfect correspondences that appeared in this paper. This paper grew out of our conversation at Utah Algebraic Geometry Conference in 1995 and thereafter and I thank him for his suggestions, comments and encouragements. I owe thanks to Victor Reiner for saving me from an embarrassment. I thank Michel Brion and Günter Ziegler for some useful conversations or correspondence regarding the materials in this paper, and also Andrzej Bialynicki-Birula for sending his papers to me. ## 1. Chambers and cones Let $`V`$ and $`W`$ be vector spaces over $``$ of finite dimensions, and $$\pi :VW$$ be a linear projection. Assume that $`PV`$ is a polytope of full dimension and $`Q=\pi (P)`$. ###### Definition 1.1. The (polytopal) cell in $`Q`$ containing a point $`qQ`$ is the intersection of $`\pi `$-images of all the faces of $`P`$ that contain the given point $`q`$, that is $$𝐜(q):=[q]:=\{\pi (F)\|q\pi (F),FL(P)\}.$$ The set of all cells is denoted by $`\mathrm{\Gamma }=\mathrm{\Gamma }(P,Q)`$, partially ordered by inclusion. Maximal cells are called chambers<sup>1</sup><sup>1</sup>1Cells are often called chambers in other literature. Then, chambers in our sense would have to be called maximal chambers.. Dually, the sequence $`0\mathrm{ker}\pi V\mathrm{@}>\pi >>W0`$ induces the sequence $$0W^{}V^{}\mathrm{@}>\pi ^{}>>(\mathrm{ker}\pi )^{}0.$$ Recall that given any linear function $`\psi V^{}`$, the set $`P^\psi `$ of all points in $`P`$ on which $`\psi `$ is maximal is a face of $`P`$, and all nonempty faces of $`P`$ arise in this way. In fact, given any $`F`$ in the face lattice $`L(P)`$ of $`P`$, the set of $`\psi V^{}`$ such that $`P^\psi =F`$ is a cone, called the normal cone of $`F`$, and denoted by $`N(P,F)=N(F)`$. The set $`\mathrm{\Delta }(P)`$ of the normal cones of all faces of $`P`$ is the normal fan of $`P`$. ###### Definition 1.2. The (polytopal) cone in $`(\mathrm{ker}\pi )^{}`$ containing a point $`\psi (\mathrm{ker}\pi )^{}`$ is the intersection of all $`\pi ^{}`$-images of the cones in $`\mathrm{\Delta }(P)`$ that contains $`\psi `$, that is, $$\sigma (\psi ):=[\psi ]:=\{\pi ^{}N(P,F)\|FL(P),\psi \pi ^{}N(P,F)\}.$$ We use $`\mathrm{\Gamma }^{}=\mathrm{\Gamma }^{}(P,Q)`$ to denote the set of all those cones, partially ordered by inclusion. (These cones form the normal fan of the fiber polytope $`\mathrm{\Sigma }(P,Q)`$ ()). Consider the polytope $`P_q=\pi ^1(q)`$ for any $`qQ`$. For any linear function $`\psi (\mathrm{ker}\pi )^{}`$, let $`P_q^\psi `$ be the nonempty face of $`P_q`$ defined by $`\psi `$. The normal cone $`N(P_q,P_q^\psi )=N(P_q^\psi )`$ depends only on the cell $`𝐜=𝐜(q)`$. So, we may use $`\mathrm{\Delta }(𝐜)`$ to be the normal fan $`\mathrm{\Delta }(P_q)`$ for all $`q\mathrm{relint}𝐜`$, where “relint” denotes “relative interior”. For each $`\psi `$, let $`[\psi ]_𝐜`$ denote the minimal cone in $`\mathrm{\Delta }(𝐜)`$ that contains $`\psi `$. Note that $`\mathrm{\Delta }(𝐜)`$ is a common refinement of $`\mathrm{\Delta }(𝐜^{})`$ for all $`𝐜^{}𝐜`$. In addition, for each face $`P_q^\psi `$ of $`P_q`$ there exists a unique minimal face of $`P`$ that contains $`P_q^\psi `$: the intersection of all faces that contain $`P_q^\psi `$. We may use $`F_{q,\psi }`$ to denote this face of $`P`$. Since it only depends on the cell $`𝐜=𝐜(\psi )`$, we may write it as $`F_{𝐜,\psi }`$. Moreover, $`F_{𝐜,\psi }`$ only depends on the cone $`[\psi ]_𝐜`$. So, we oftentimes denote it by $`F_{𝐜,[\psi ]_𝐜}`$. Note that $`F_{𝐜,\psi }=F_{𝐜,\psi ^{}}`$ if and only if $`[\psi ]_𝐜=[\psi ^{}]_𝐜`$. Note also that $`N(P,F_{q,\psi })`$ projects onto $`N(P_q,P_q^\psi )`$ under the projection $`\pi ^{}`$. Having the above explained, we can conclude that the fan of cones in the definition 1.2 is the common refinement of all the normal fans $`\mathrm{\Delta }(𝐜)`$ ($`𝐜\mathrm{\Gamma }`$). ## 2. Coherent strings Given any linear function $`\psi (\mathrm{ker}\pi )^{}`$, we have a collection of faces in $`L(P)`$ defined as follows: $$(\psi )=\{F_{𝐜,\psi }\|𝐜\mathrm{\Gamma }\}.$$ ###### Definition 2.1. (See ) The above collection $`(\psi )=\{F_{𝐜,\psi }\|𝐜\mathrm{\Gamma }\}`$ of faces of $`P`$ is called a (global) coherent string of the projection $`\pi :PQ`$. We set $`𝒯_{coh}=𝒯_{coh}(P,Q)`$ to be the set of all coherent strings of $`\pi :PQ`$, partially ordered by: $`(\psi )(\psi ^{})\text{if}(\psi )(\psi ^{}).`$ ###### Proposition 2.2. Given any linear function $`\psi (\mathrm{ker}\pi )^{}`$, the coherent string $`(\psi )`$ satisfies the following properties 1. $`\{\pi (F_{𝐜,\psi })\|𝐜\mathrm{\Gamma }\}`$ is a polytopal subdivision of $`Q`$ without repetitions; 2. $`\pi (F_{𝐜^{},\psi })\pi (F_{𝐜,\psi })`$ if and only if $`F_{𝐜^{},\psi }=\pi ^1(\pi (F_{𝐜^{},\psi }))F_{𝐜,\psi }`$. Proof. Obviously, we have $$\{\pi (F_{𝐜,\psi })\|𝐜\mathrm{\Gamma }\}=𝐜=Q.$$ But, $`q\mathrm{relint}\pi (F_{𝐜,\psi })\mathrm{relint}\pi (F_{𝐜^{},\psi })\mathrm{}`$ if and only if $`F_{𝐜,\psi }=_{FP_q^\psi }F=F_{𝐜^{},\psi }`$. This proves (1). (2) follows easily. ###### Lemma 2.3. $`(\psi )=(\psi ^{})`$ if and only if $`\sigma (\psi )=\sigma (\psi ^{})`$. Proof. $`(\psi )=(\psi ^{})`$ if and only if $`F_{𝐜,\psi }=F_{𝐜,\psi ^{}}`$ for all $`𝐜\mathrm{\Gamma }`$ if and only if $`[\psi ]_𝐜=[\psi ^{}]_𝐜`$ for all $`𝐜\mathrm{\Gamma }`$ if and only if $`\sigma (\psi )=\sigma (\psi ^{})`$. So, we may denote a coherent string by $`(\sigma )`$ for some $`\sigma \mathrm{\Gamma }^{}`$. Given any point $`qQ`$, we use $`[q]_\sigma `$ to denote the minimal face in $`(\sigma )`$ such that its $`\pi `$-image contains $`q`$. Observe that $`F_{q,\sigma }=F_{q^{},\sigma }`$ if and only if $`[q]_\sigma =[q^{}]_\sigma `$. ###### Remark 2.4. Notice that two different coherent strings may induce identical subdivision of $`Q`$. This will not be the case if the projection is non-degenerate, i.e., sends distinct vertices of $`P`$ to distinct points. In general, the issue can typically be dealt with by keeping the label of each vertex. ###### Proposition 2.5. (Billera-Sturmfels ) The poset $`𝒯_{coh}`$ of all coherent strings and the poset $`\mathrm{\Gamma }^{}`$ are anti-isomorphic. That is, $`𝒯_{coh}(\mathrm{\Gamma }^{})^{opp}`$ and $`\mathrm{\Gamma }^{}(𝒯_{coh})^{opp}`$. Proof. Clearly the function $`:\mathrm{\Gamma }^{}𝒯_{coh}`$ is a bijection by the definition and Lemma 2.3. Let $`\sigma `$ be a cone and $`(\sigma )`$ be its corresponding coherent string. Assume that $`\sigma \sigma `$. Pick $`\psi \sigma `$ and $`\psi ^{}\sigma `$. Then $`P_q^\psi P_q^\psi ^{}`$. That is, $`F_{q,\psi }F_{q,\psi ^{}}`$. Since $`(\sigma )=\{F_{q,\psi }\|qQ\}`$ and $`(\sigma )=\{F_{q,\psi ^{}}\|qQ\}`$, we obtain that $`(\sigma )(\sigma )`$. That is, the function $`:\mathrm{\Gamma }^{}𝒯_{coh}`$ is an order-reversing bijection. ## 3. Locally coherent strings Locally coherent strings are natural generalization of coherent strings. Geometrically, they are “continuous liftings of polytopal subdivisions of $`Q`$” to the face lattice of $`P`$. It turns out that all of these things can be constructed as follows. Consider an arbitrary map $$\stackrel{~}{\mathrm{\Psi }}:\mathrm{\Gamma }(\mathrm{ker}\pi )^{}.$$ It gives rise to a collection of faces of $`P`$ as follows $$(\stackrel{~}{\mathrm{\Psi }})=\{F_{𝐜,\stackrel{~}{\mathrm{\Psi }}(𝐜)}\|𝐜\mathrm{\Gamma }\}.$$ ###### Definition 3.1. (See ) The collection $`(\stackrel{~}{\mathrm{\Psi }})`$ of faces of $`P`$ is called a locally coherent string if $`\stackrel{~}{\mathrm{\Psi }}`$ satisfies the so-called locally coherent condition: that is, for any $`𝐜𝐜^{}`$ $$\mathrm{relint}[\stackrel{~}{\mathrm{\Psi }}(𝐜)]_𝐜\mathrm{relint}[\stackrel{~}{\mathrm{\Psi }}(𝐜^{})]_𝐜^{}.$$ We set $`𝒯=𝒯(P,Q)`$ to denote the set of all locally coherent strings, partially ordered by: $`\stackrel{~}{\mathrm{\Psi }}\stackrel{~}{\mathrm{\Psi }^{}}`$ if and only if $`(\stackrel{~}{\mathrm{\Psi }})(\stackrel{~}{\mathrm{\Psi }^{}})`$. One can check that the locally coherent condition is equivalent to $$\mathrm{relint}\pi ^{}(N(F_{𝐜,\stackrel{~}{\mathrm{\Psi }}(𝐜)}))\mathrm{relint}\pi ^{}(N(F_{𝐜^{},\stackrel{~}{\mathrm{\Psi }}(𝐜^{})})).$$ ###### Proposition 3.2. (, ) Assume that a map $`\stackrel{~}{\mathrm{\Psi }}:\mathrm{\Gamma }(\mathrm{ker}\pi )^{}`$ satisfies the locally coherent condition. Then the locally coherent string $`(\stackrel{~}{\mathrm{\Psi }})`$ enjoys 1. $`\{\pi (F)\|F(\stackrel{~}{\mathrm{\Psi }})\}`$ is a subdivision of $`Q`$ without repetition. 2. $`\pi (F)\pi (F^{})`$ if and only if $`F=F^{}\pi ^1(\pi (F))`$ for any $`F,F^{}(\stackrel{~}{\mathrm{\Psi }})`$. Furthermore, every collection of the faces of $`P`$ satisfying (1) and (2) arises in this way. Proof. The theorem was proved in for the case when $`dimQ=1`$ and was explicitly pointed out in for the general cases. A key observation in the proof is that when $`𝐜^{}<𝐜`$, the locally coherent condition implies that $`F_{𝐜^{},\stackrel{~}{\mathrm{\Psi }}(𝐜^{})}`$ is either identical to $`F_{𝐜,\stackrel{~}{\mathrm{\Psi }}(𝐜)}`$ or is a proper face of $`F_{𝐜,\stackrel{~}{\mathrm{\Psi }}(𝐜)}`$. (1) and (2) basically follow from this observation. Conversely, given any collection $``$ satisfying (1) and (2). For any $`𝐜\mathrm{\Gamma }`$, let $`F(𝐜)`$ be the unique minimal face in $``$ such that $`\mathrm{relint}𝐜\mathrm{relint}\pi (F(𝐜))`$. Let $`\psi (𝐜)\mathrm{relint}\pi ^{}(N(F(𝐜)))(\mathrm{ker}\pi )^{}`$ such that $`F(𝐜)=F_{𝐜,\psi (𝐜)}`$. Define $$\stackrel{~}{\mathrm{\Psi }}:\mathrm{\Gamma }(\mathrm{ker}\pi )^{}$$ as $$\stackrel{~}{\mathrm{\Psi }}:𝐜\psi (𝐜).$$ Now, if $`𝐜^{}𝐜`$ (we may well assume that $`𝐜`$ covers $`𝐜^{}`$), either $`F(𝐜^{})=F(𝐜)`$ (in this case, $`\mathrm{relint}𝐜^{}\mathrm{relint}\pi (F(𝐜))`$) or $`F(𝐜^{})`$ is a proper face of $`F(𝐜)`$ (in this case, $`𝐜^{}`$ lies on the boundary of $`\pi (F(𝐜))`$). In the first case, we have $`[\psi (𝐜)]_𝐜`$ equals $`[\psi (𝐜^{})]_𝐜^{}`$ which obviously implies that $`\mathrm{relint}[\psi (𝐜)]_𝐜\mathrm{relint}[\psi (𝐜^{})]_𝐜^{}`$. Assume we are in the second case. We then have that $`[\psi (𝐜)]_𝐜`$ divides $`[\psi (𝐜^{})]_𝐜^{}`$ from its relative interior which also implies that $`\mathrm{relint}[\psi (𝐜)]_𝐜\mathrm{relint}[\psi (𝐜^{})]_𝐜^{}`$. This completes the proof. The minimal elements of $`𝒯`$ tie to the following ###### Definition 3.3. $`(\stackrel{~}{\mathrm{\Psi }})`$ is said to be tight if $$dim\pi (F)=dimF$$ for all $`F(\stackrel{~}{\mathrm{\Psi }})`$. (That is, $`F`$ does not drop dimension under the projection.) ###### Proposition 3.4. $`(\stackrel{~}{\mathrm{\Psi }})`$ is minimal if and only if $`(\stackrel{~}{\mathrm{\Psi }})`$ is tight. Proof. This is Lemma 9.5 of whose proof for the coherent case extends to the general case. The proof goes as follows. Assume that $`F(\stackrel{~}{\mathrm{\Psi }})`$ is minimal but there exists a face $`FF(\stackrel{~}{\mathrm{\Psi }})`$ such that the projection $`F\pi (F)`$ drops dimension. Now consider the polytopal projection $`F\pi (F)`$, take any non-trivial coherent string $`\mathrm{\Psi }_F`$ of this projection, and substitute $`F`$ from $`(\stackrel{~}{\mathrm{\Psi }})`$ by $`(\mathrm{\Psi }_F)`$, we will get a new locally coherent string of the original projection which is certainly smaller than $`(\stackrel{~}{\mathrm{\Psi }})`$. Hence $`(\stackrel{~}{\mathrm{\Psi }})`$ has to be tight. The converse is easy. We omit further details. ## 4. Coherent costrings Let $`qQ`$ be any interior point. The point $`q`$ selects the following collection of cones in $`\mathrm{\Delta }(P)`$ $$\mathrm{\Delta }(q)=\{N(P,F_{q,\psi })\|\psi (\mathrm{ker}\pi )^{}\}.$$ The collection only depends on the cell $`𝐜(q)`$. That is, ###### Lemma 4.1. $`\mathrm{\Delta }(q)=\mathrm{\Delta }(q^{})`$ if and only if $`𝐜(q)=𝐜(q^{})`$. Proof. $`\mathrm{\Delta }(q)=\mathrm{\Delta }(q^{})`$ if and only if $`F_{q,\sigma }=F_{q^{},\sigma }`$ for all $`\sigma \mathrm{\Gamma }^{}`$ if and only if $`[q]_\sigma =[q^{}]_\sigma `$ for all $`\sigma \mathrm{\Gamma }^{}`$ if and only if $`𝐜(q)=𝐜(q^{})`$. Thus we may denote $`\mathrm{\Delta }(q)`$ by $`\mathrm{\Delta }(𝐜)`$ if $`q\mathrm{relint}𝐜`$ ($`𝐜\mathrm{\Gamma }`$). ###### Definition 4.2. $`\mathrm{\Delta }(𝐜)`$ is called a coherent costring. Let $`𝒯_{coh}^{}=\{\mathrm{\Delta }(𝐜)\|𝐜\mathrm{\Gamma }\}`$ be the set of all coherent costrings, partially ordered by: $$\mathrm{\Delta }(𝐜)\mathrm{\Delta }(𝐜^{})\text{if}\mathrm{\Delta }(𝐜)\mathrm{\Delta }(𝐜^{}).$$ ###### Proposition 4.3. Given any $`𝐜\mathrm{\Gamma }`$. The collection $`\mathrm{\Delta }(𝐜)`$ satisfies the following properties 1. $`\{\pi ^{}(N(P,F_{𝐜,\psi }))\|\psi (\mathrm{ker}\pi )^{}\}`$ is a cone subdivision of $`(\mathrm{ker}\pi )^{}`$ without repetitions. 2. $`\pi ^{}(N(P,F_{𝐜,\psi }))\pi ^{}(N(P,F_{𝐜^{},\psi }))`$ if and only if $$N(P,F_{𝐜,\psi })=(\pi ^{})^1(\pi ^{}(N(P,F_{𝐜,\psi })))\pi ^{}(N(P,F_{𝐜^{},\psi })).$$ Moreover, $`\{\pi ^{}(\mathrm{\Delta }(q))\}`$ equals the normal fan $`\mathrm{\Delta }(P_q)`$ if $`qQ`$. Proof. The proof of (1) and (2) is analogous to that of Proposition 2.2 and is a careful check using definitions. In section 5, we will prove a general version (i.e., for locally coherent costrings) which includes this as a particular case. For the last statement, just observe that $`\pi ^{}(N(P,F_{q,\psi }))=N(P_q,P_q^\psi )`$. We omit further details. ###### Remark 4.4. Like in Remark 2.4, two different coherent costrings may induce identical cone subdivision of $`(\mathrm{ker}\pi )^{}`$. ###### Proposition 4.5. The poset $`𝒯_{coh}^{}`$ of all coherent costrings and the poset $`\mathrm{\Gamma }`$ of all cells are anti-isomorphic. That is, $`𝒯_{coh}^{}=\mathrm{\Gamma }^{opp}`$ and $`\mathrm{\Gamma }=(𝒯_{coh}^{})^{opp}`$. Proof. Treat $`\mathrm{\Delta }`$ as a function from $`\mathrm{\Gamma }`$ to $`𝒯_{coh}^{}`$. This is then a bijection by Lemma 4.1. Assume that $`𝐜^{}𝐜`$ are two cells. Pick $`q_𝐜^{}\mathrm{relint}𝐜^{}`$ and $`q_𝐜\mathrm{relint}𝐜`$. We have that $`F_{q_𝐜^{},\psi }F_{q_𝐜,\psi }`$ for all $`\psi (\mathrm{ker}\pi )^{}`$. That is, $`N(P,F_{q_𝐜^{},\psi })N(P,F_{q_𝐜,\psi })`$ for all $`\psi `$. By definition, we obtain $`\mathrm{\Delta }(𝐜^{})\mathrm{\Delta }(𝐜)`$. That is, $`\mathrm{\Delta }`$ is an order-reversing bijection. ## 5. Locally coherent costrings Locally coherent costrings are generalizations of (global) coherent costrings. Geometrically, they are $`\pi ^{}`$-liftings of conical subdivisions of $`(\mathrm{ker}p)^{}`$. ###### Definition 5.1. Let $`\mathrm{\Delta }_0`$ be an arbitrary fan in $`V^{}`$. A subcollection $`\mathrm{\Delta }`$ of cones in $`\mathrm{\Delta }_0`$ is called a locally coherent costring of $`\pi ^{}:\|\mathrm{\Delta }_0\|V^{}\pi ^{}(\|\mathrm{\Delta }_0\|)(\mathrm{ker}\pi )^{}`$ if it satisfies that 1. $`\{\pi ^{}(\sigma )\|\sigma \mathrm{\Delta }\}`$ form a cone subdivision of $`\pi ^{}(\|\mathrm{\Delta }_0\|)`$ without repetition. 2. $`\pi ^{}(\sigma )\pi ^{}(\sigma ^{})`$ if and only if $`\sigma =(\pi ^{})^1(\pi ^{}(\sigma ))\sigma ^{}`$. It follows from the definition that if $`\sigma \sigma ^{}\mathrm{\Delta }`$, then $`\pi ^{}(\sigma )`$ and $`\pi ^{}(\sigma ^{})`$ must intersect along a proper face. In particular, $`\pi ^{}(\mathrm{\Delta })=\{\pi ^{}(\sigma )\|\sigma \mathrm{\Delta }\}`$ is a fan in $`W`$. (Warning: a fan in our sense may not be strongly convex.) The following picture illustrates an example of a (locally) coherent costring $`\mathrm{\Delta }`$ in $`^3`$ with the projection $`\pi ^{}(\mathrm{\Delta })`$ in $`^2`$. (See for pictures of locally coherent strings.) Figure 1 The most interesting case where rich structures and duality emerge is the case when $`\mathrm{\Delta }_0`$ is the normal fan $`\mathrm{\Delta }(P)`$ of the polytope $`P`$. ###### Definition 5.2. The set of locally coherent costrings of $`\pi ^{}:\|\mathrm{\Delta }(P)\|(\mathrm{ker}\pi )^{}`$ is denoted by $`𝒯^{}=𝒯^{}(P,Q)`$. This set is equipped with a natural partial order by inclusion. Note that $`\mathrm{\Gamma }^{}`$ is a common refinement of $`\pi ^{}(\mathrm{\Delta })`$ for all $`\mathrm{\Delta }𝒯^{}`$. As before, we will actually construct all locally coherent costrings. So, again we consider a map $$\stackrel{~}{\mathrm{\Delta }}:\mathrm{\Gamma }^{}Q$$ satisfying the locally coherent condition: for any $`\sigma >\sigma ^{}`$, $$\mathrm{relint}[\stackrel{~}{\mathrm{\Delta }}(\sigma )]_\sigma \mathrm{relint}[\stackrel{~}{\mathrm{\Delta }}(\sigma ^{})]_\sigma ^{}.$$ $`\stackrel{~}{\mathrm{\Delta }}`$ gives rise to a collection $`𝒞(\stackrel{~}{\mathrm{\Delta }})`$ of cones in $`\mathrm{\Delta }(P)`$ as follows $$𝒞(\stackrel{~}{\mathrm{\Delta }})=\{N(P,F_{\stackrel{~}{\mathrm{\Delta }}(\sigma ),\sigma })\|\sigma \mathrm{\Gamma }^{}\}.$$ ###### Theorem 5.3. Let $`\stackrel{~}{\mathrm{\Delta }}:\mathrm{\Gamma }^{}Q`$ satisfy the locally coherent condition. $`𝒞(\stackrel{~}{\mathrm{\Delta }})`$ is a locally coherent costring of $`\pi ^{}:\|\mathrm{\Delta }(P)\|(\mathrm{ker}\pi )^{}`$. Furthermore, every locally coherent costring of $`\pi ^{}:\|\mathrm{\Delta }(P)\|(\mathrm{ker}\pi )^{}`$ arises in this way. Proof. The proof of this theorem is in a dual manner analogous to the case of locally coherent strings. We see immediately that $$_{\sigma \mathrm{\Gamma }^{}}\pi ^{}(N(F_{\stackrel{~}{\mathrm{\Delta }}(\sigma ),\sigma }))_{\sigma \mathrm{\Gamma }^{}}\sigma =(\mathrm{ker}\pi )^{}$$ which implies the equality. To see that we have a subdivision without repetition, just observe that if $`\sigma >\sigma ^{}`$, either $`F_{\stackrel{~}{\mathrm{\Delta }}(\sigma ),\sigma }=F_{\stackrel{~}{\mathrm{\Delta }}(\sigma ^{}),\sigma ^{}})`$ or $`F_{\stackrel{~}{\mathrm{\Delta }}(\sigma ),\sigma }`$ is a face of $`F_{\stackrel{~}{\mathrm{\Delta }}(\sigma ^{}),\sigma ^{}})`$ by the locally coherent condition. Condition (2) of Definition 5.1 follows from the same observation. Conversely, for any $`\sigma \mathrm{\Gamma }^{}`$, let $`N(F(\sigma ))`$ be the minimal cone in $`𝒞`$ such that $`\pi ^{}(N(F(\sigma )))`$ contains $`\sigma `$. Let $`q(\sigma )\mathrm{relint}F(\sigma )`$ be a point such that $`F(\sigma )=F_{q(\sigma ),\sigma }`$. Define $$\stackrel{~}{\mathrm{\Delta }}:\mathrm{\Gamma }^{}Q$$ by $$\stackrel{~}{\mathrm{\Delta }}:\sigma q(\sigma ).$$ If $`\sigma >\sigma ^{}`$ (it suffices to assume that $`\sigma `$ covers $`\sigma ^{}`$, that is, $`dim\sigma =dim\sigma ^{}+1`$), since the subdivision of $`Q`$ induced by the coherent string defined by $`\sigma `$ refines the subdivision of $`Q`$ induced by the coherent string defined by $`\sigma ^{}`$, we have either $`F(\sigma )=F(\sigma ^{})`$ or $`F(\sigma )`$ divides $`F(\sigma ^{})`$ from its relative interiors. In either case, it implies that $`\mathrm{relint}[q(\sigma )]_\sigma \mathrm{relint}[q(\sigma ^{})]_\sigma ^{}`$. In other words, the map $`\stackrel{~}{\mathrm{\Delta }}`$ indeed satisfies the locally coherent condition. ###### Definition 5.4. A locally coherent costring $`𝒞(\stackrel{~}{\mathrm{\Delta }})`$ is tight if $`dim\pi ^{}(\sigma )=dim\sigma `$ for every $`\sigma 𝒞(\stackrel{~}{\mathrm{\Delta }})`$. Analogous to Proposition 3.4 for locally coherent strings, we have the following proposition for locally coherent costrings whose proof is similar to that of Proposition 3.4 and is left to the reader. ###### Proposition 5.5. A locally coherent costring $`𝒞(\stackrel{~}{\mathrm{\Delta }})`$ is tight if and only if it is minimal. ###### Remark 5.6. In view that $`\mathrm{\Gamma }^{}=𝒯_{coh}`$, every locally coherent costring $`𝒞(\mathrm{\Delta })`$ singles out a unique minimal element from any given coherent string $`(\sigma )𝒯_{coh}`$: i.e., the minimal face in $`(\sigma )`$ such that its image contains the cell $`\stackrel{~}{\mathrm{\Delta }}(\sigma )`$. ## 6. Virtual chambers Virtual cells are generalizations of the real cells. In toric geometry, real cells define projective quotient varieties, while virtual cells define (not-necessarily projective) quotient varieties. We refer to for a treatment of virtual chambers in the case of triangulations of a point set configuration. Any map $`\stackrel{~}{𝐜}:\mathrm{\Gamma }^{}Q`$ gives rise to a collection of faces of $`P`$ as follows $$(\stackrel{~}{𝐜})=\{F_{\stackrel{~}{𝐜}(\sigma ),\sigma }\|\sigma \mathrm{\Gamma }^{}\}.$$ ###### Definition 6.1. $`(\stackrel{~}{𝐜})`$ is called a virtual cell if $`\stackrel{~}{𝐜}`$ satisfies the following locally coherent condition: for $`\sigma \sigma ^{}`$ $$\mathrm{relint}[\stackrel{~}{𝐜}(\sigma )]_\sigma \mathrm{relint}[\stackrel{~}{𝐜}(\sigma ^{})]_\sigma ^{}.$$ We use $`\mathrm{\Gamma }_{vir}=\mathrm{\Gamma }_{vir}(P,Q)`$ to denote the set of all virtual cells, partially ordered by inclusion: $`\stackrel{~}{𝐜}\stackrel{~}{𝐜}^{}`$ if $`(\stackrel{~}{𝐜})(\stackrel{~}{𝐜}^{})`$. Maximal virtual cells are called virtual chambers. ###### Theorem 6.2. Assume that a map $`\stackrel{~}{𝐜}:\mathrm{\Gamma }^{}Q`$ satisfies the locally coherent condition. Then the virtual cell $`(\stackrel{~}{𝐜})`$ satisfies the following property: for any coherent string $`\mathrm{\Psi }`$, $`𝒞(\stackrel{~}{𝐜})(\mathrm{\Psi })`$ is nonempty and contains exactly one minimal element. Furthermore, every collection of the faces of $`P`$ satisfying the above property arises in this way. Proof. Let $`=(\sigma )𝒯_{coh}`$ with $`\sigma \mathrm{\Gamma }^{}`$. Clearly both $`(\sigma )`$ and $`(\stackrel{~}{𝐜})`$ contain the element $`F_{\stackrel{~}{𝐜}(\sigma ),\sigma }`$. That is, the intersection is not empty. Now given any two faces $`F_{p,\sigma }`$ and $`F_{q,\sigma }(\sigma )(\stackrel{~}{𝐜})`$, either they have a common face or they do not meet at all. If they do not meet at all, $`\pi ^{}(N(F_{p,\sigma }))`$ and $`\pi ^{}(N(F_{q,\sigma }))`$ will not have a non-trivial common face. Hence they meet in relative interiors (both contain $`\sigma `$). But as members of the locally coherent costring $`𝒞(\mathrm{\Delta })`$ (transported from $`\mathrm{\Delta }=\stackrel{~}{c}`$), they can not meet in relative interiors. This shows that $`F_{p,\sigma }F_{q,\sigma }`$ is a common face of each. Clearly, $`(\sigma )(\stackrel{~}{𝐜})`$ must, as a subset of a coherent string $`(\sigma )`$ and with property that every two members have a common face, have a unique minimal element. On the other hand, given a collection $``$ that meets every coherent string $`(\sigma )`$ in exactly one minimal element. Let $`F(\sigma )(\sigma )`$ be this face. Pick a point $`q_\sigma \mathrm{relint}\pi (F(\sigma ))`$. We obtain a map $$\stackrel{~}{𝐜}:\mathrm{\Gamma }^{}Q$$ by sending $`\sigma `$ to $`q_\sigma `$. We now check that this map satisfies the locally coherent condition. For any $`\sigma ^{}\sigma `$, we have that the subdivision $`\{\pi (F)\|F(\sigma )\}`$ refines the subdivision $`\{\pi (F)\|F(\sigma ^{})\}`$. It follows that $`\mathrm{relint}\pi (F(\sigma ))\mathrm{relint}\pi (F(\sigma ^{}))`$ because either $`F(\sigma )=F(\sigma ^{})`$ or $`\pi (F(\sigma ))`$ divides $`\pi (F(\sigma ))`$ in its relative interiors. This checks the locally coherent condition. ###### Proposition 6.3. The following two are equivalent. 1. For any coherent string $`\mathrm{\Psi }`$, $`(\stackrel{~}{𝐜})(\mathrm{\Psi })`$ is nonempty and contains exactly one minimal element. 2. For any locally coherent string $`\stackrel{~}{\mathrm{\Psi }}`$, $`(\stackrel{~}{𝐜})(\stackrel{~}{\mathrm{\Psi }})`$ is nonempty and contains exactly one minimal element. Proof. For point configurations in general position, this is proved in for virtual chambers. (2) obviously implies (1). To see the other direction, assume $`(\stackrel{~}{𝐜})`$ satisfies (1) but not (2), then there will be two faces $`F`$ and $`F^{}`$ in some locally coherent string $`\stackrel{~}{\mathrm{\Psi }}`$ such that they are both contained in $`(\stackrel{~}{𝐜})`$ and are two distinct minimals in $`(\stackrel{~}{𝐜})(\stackrel{~}{\mathrm{\Psi }})`$. Since $`\pi (F)\pi (F^{})`$ must be empty, we can find a coherent string $`\mathrm{\Psi }`$ to contain both $`F`$ and $`F^{}`$. This contradicts to that $`(\stackrel{~}{𝐜})`$ satisfies (1). This equivalence will not, however, be used later. ## 7. Virtual cones Virtual cones are simply dual versions of virtual cells. An arbitrary map $`\stackrel{~}{\sigma }:\mathrm{\Gamma }(\mathrm{ker}\pi )^{}`$ defines a collection $`𝒞(\stackrel{~}{\sigma })`$ of cones of $`\mathrm{\Delta }(P)`$ $$𝒞(\stackrel{~}{\sigma })=\{N(F_{𝐜,\stackrel{~}{\sigma }(𝐜)},P)\|𝐜\mathrm{\Gamma }\}.$$ ###### Definition 7.1. $`𝒞(\stackrel{~}{\sigma })`$ of cones of $`\mathrm{\Delta }(P)`$ is called a virtual cone if $`\stackrel{~}{\sigma }`$ satisfies the locally coherent condition. The set of all virtual cones is denoted by $`\mathrm{\Gamma }_{vir}^{}=\mathrm{\Gamma }_{vir}^{}(P,Q)`$, partially ordered by inclusion. ###### Theorem 7.2. Let $`\stackrel{~}{\sigma }:\mathrm{\Gamma }(\mathrm{ker}\pi )^{}`$ satisfy the locally coherent condition. Then the virtual cone $`𝒞(\stackrel{~}{\sigma })`$ satisfies the following property. For every coherent costring $`𝒞(𝐜)`$, $`𝒞(\stackrel{~}{\sigma })𝒞(𝐜)`$ is non-empty and contains exactly one maximal member. Furthermore, every collection of cones of $`\mathrm{\Delta }(P)`$ that satisfies above property arises in this way. Proof. Transporting $`𝒞(𝐜)`$ and $`𝒞(\stackrel{~}{\sigma })`$ to their corresponding face collections $`(𝐜)`$ and $`(\stackrel{~}{\mathrm{\Psi }})`$ where $`\stackrel{~}{\mathrm{\Psi }}=\stackrel{~}{\sigma }`$, we see immediately that $`(\stackrel{~}{\mathrm{\Psi }})(𝐜)`$ contains a unique minimal element (i.e., the intersection of all the faces of $`(\stackrel{~}{\mathrm{\Psi }})`$ that contain $`𝐜`$). Translating back, this means that $`𝒞(\stackrel{~}{\sigma })𝒞(𝐜)`$ contains exactly one maximal member. ###### Remark 7.3. Again, the following should be equivalent. 1. $`𝒞(\stackrel{~}{\sigma })`$ meets every coherent costring in exactly one maximal member. 2. $`𝒞(\stackrel{~}{\sigma })`$ meets every locally coherent costring in exactly one maximal member. ## 8. Summary of dualities It may be worthwhile to summarize some of the results in the previous section as the following dualities. ###### Theorem 8.1. The following always holds. 1. A collection $``$ of the faces of $`P`$ is a locally coherent string of $`\pi :PQ`$ if and only if $`𝒞=\{N(P,F)\|F\}`$ is a virtual cone of $`\mathrm{\Delta }(P)(\mathrm{ker}\pi )^{}`$. 2. A collection $``$ of the faces of $`P`$ is a virtual cell of $`\pi :PQ`$ if and only if $`𝒞=\{N(P,F)\|F\}`$ is a locally coherent costring of $`\mathrm{\Delta }(P)(\mathrm{ker}\pi )^{}`$. 3. A collection $`𝒞`$ of the cones of $`\mathrm{\Delta }(P)`$ is a locally coherent costring of $`\mathrm{\Delta }(P)(\mathrm{ker}\pi )^{}`$ if and only if $`=\{F\|N(P,F)𝒞\}`$ is a virtual cell of $`PQ`$. 4. A collection $`𝒞`$ of the cones of $`\mathrm{\Delta }(P)`$ is a virtual cone of $`\mathrm{\Delta }(P)(\mathrm{ker}\pi )^{}`$ if and only if $`=\{F\|N(P,F)𝒞\}`$ is a locally coherent string of $`\pi :PQ`$. Furthermore, the correspondence in each of (1), (2), (3), (4) reverses the inclusion-induced partial orders. Proof. This theorem is a corollary to the combination of the results in the previous sections. Let $`\mathrm{Mor}^{}(S,S^{})`$ be the poset of order-reversing maps between two posets $`S`$ and $`S^{}`$. The we have 1. $`𝒯\mathrm{Mor}^{}(\mathrm{\Gamma },\mathrm{\Gamma }^{})`$ and elements in $`𝒯_{coh}`$ correspond to constant maps; 2. $`𝒯^{}\mathrm{Mor}^{}(\mathrm{\Gamma }^{},\mathrm{\Gamma })`$ and elements in $`𝒯_{coh}`$ correspond to constant maps; 3. $`\mathrm{\Gamma }_{vir}^{}\mathrm{Mor}^{}(\mathrm{\Gamma },\mathrm{\Gamma }^{})=\mathrm{Mor}^{}(𝒯_{coh}^{},𝒯_{coh})`$ and elements in $`\mathrm{\Gamma }^{}`$ correspond to constant maps; 4. $`\mathrm{\Gamma }_{vir}\mathrm{Mor}^{}(\mathrm{\Gamma }^{},\mathrm{\Gamma })=\mathrm{Mor}^{}(𝒯_{coh},𝒯_{coh}^{})`$ and elements in $`\mathrm{\Gamma }`$ correspond to constant maps; Note the natural correspondences amongst the bigger posets. These relations seem to worth some investigation. In particular, it would be interesting to know if an arbitrary element in each of the bigger posets corresponds to some meaningful combinatorial construction. ## 9. Realizable virtual chamber. Examples indicate that most of virtual cells can be obtained by topologically deforming the original polytopal cell complex. We now explain this. Let $`P_t`$ be a topological deformation of $`P`$, while the lattice structure on $`L(P)`$ or all of the face relations in $`L(P)`$ are memorized. Then the image $`Q_t=\pi (P_t)`$ will be called a tamed deformation of $`Q=Q_0`$. Note that the deformation also induces a (tamed) deformation of walls, and in particular, a (tamed) deformation $`\mathrm{\Gamma }_t`$ of the chamber complex $`\mathrm{\Gamma }`$. Another possible way to think of a tamed deformation is to interpret it as a deformation $`\pi _t`$ of the affine projection $`\pi =\pi _0`$. In that way, all the twisted projections $`\pi _t`$ ($`t0`$) may be called virtual affine projections. Thus a tamed deformation $`Q_t`$ and its induced chamber complex $`\mathrm{\Gamma }_t`$ will just be the image of $`P`$ under the virtual affine projection $`\pi _t`$. ###### Example 9.1. Figure 2 is a typical tamed deformation of a hexagon and its internal wall structures. Note that by the deformation, we gain one more (deformed) chamber, three more deformed 1-dimensional cells, and two more 0-dimensional cells (gain three but loose one). The next picture is not a tamed deformation because an interior wall and a face of the hexagon intersect at a wrong place. (Notice the gain of an extra chamber in the upper right corner.) Figure 2 Figure 3 ###### Remark 9.2. In terms of geometry of the torus action, the above means that one can not alert the intersection relations among fixed point varieties (for various subtori) under the deformation. ###### Proposition 9.3. Every tamed deformation of $`P`$ induces a tamed deformation of the $`\pi `$-induced polytopal subdivision $`\pi (T)`$ of $`Q`$ for any given $`T𝒯`$. In particular, the combinatorial type of $`\pi (T)`$ remains the same under the deformation. It follows that if $`\mathrm{\Gamma }_t`$ is the cell complex induced by $`P_t`$, then $`\mathrm{\Gamma }_t`$ is the common refinement of $`\pi (T_t)`$ for all $`T𝒯`$. ###### Corollary 9.4. A deformed cell in $`\mathrm{\Gamma }_t`$ corresponds canonically to a virtual cell in $`\mathrm{\Gamma }`$. Proof. Each cell of $`\mathrm{\Gamma }_t`$ specifies a unique minimal deformed subpolytope in $`\pi (T_t)`$ (thus a unique minimal face $`F_t`$ in $`T_t`$) for every $`TT`$. This leads to a unique minimal face $`F`$ in the original $`T`$. By Theorem 6.2, we conclude that every cell of $`\mathrm{\Gamma }_t`$ gives rise to a virtual cell of $`\pi :PQ`$. ###### Definition 9.5. A virtual cell (chamber) that corresponds to some deformed cell (chamber) is called a realizable virtual cell (chamber). ###### Definition 9.6. A deformed wall in $`Q_t`$ will be called a virtual wall. Of Figure 2, the three bended walls are typical virtual walls. Some realizable virtual chambers deserve special attention. ###### Definition 9.7. () A cell in $`\mathrm{\Gamma }_t`$ ($`t=0`$ or otherwise) that contains an extremal vertex is called lexicographic. It is easy to see that the realizable virtual cell that corresponds to a lexicographic cell is independent of the deformation. From this, we get ###### Corollary 9.8. The subpost of realizable virtual cells is connected. ###### Remark 9.9. It is not clear if all virtual cells are realizable. By the work of , it looks that it is the case when $`dimQ=1`$ and perhaps also when $`dimQ=2`$. In general, results on the generalized Baues conjecture suggests that it should fail in general. ## 10. Toric varieties Let $`X^d`$ be an arbitrary toric variety over the field of complex numbers, not necessarily quasi-projective, acted upon by the open dense torus $`(^{})^d`$. We will study categorical (good) quotients of $`X`$ by a subgroup $`G`$ of $`(^{})^d`$. The method renders it useful the combinatorial concepts of §2. Fix a lattice $`M^d`$ and its dual lattice $`N`$. A fan $`\mathrm{\Delta }_0`$ in $`N_{}=N`$ consists of a finite collection of strongly convex rational polyhedral cones $`\sigma N_{}`$ which is closed under intersection and taking faces. Given (and fix) the fan $`\mathrm{\Delta }_0`$ and a cone $`\sigma `$ in the fan, the dual cone $$\sigma ^{}=\{vM_{}:v,\sigma 0\}$$ determines the semigroup algebra $`[\sigma ^{}M]`$ which is finitely generated and the affine variety $$A_\sigma =\mathrm{Spec}([\sigma ^{}M]).$$ It turns out that these affine varieties can be glued together to get a toric variety $$X_{\mathrm{\Delta }_0}=\underset{\sigma \mathrm{\Delta }_0}{}\mathrm{Spec}([\sigma ^{}M]).$$ Recall that $`T_N=N^{}=\mathrm{Hom}_{}(M,^{})`$ acts on $`[M]`$ by the formula $$\gamma _n(t)\chi ^m=t^{n,m}\chi ^m$$ where $`\gamma _n`$ is the one-parametric subgroup generated by $`nN`$. $`X_{\mathrm{\Delta }_0}`$ can also be constructed as a quotient of certain open subset of some affine space by a diagonalizable group. We refer the reader to or for more details. ## 11. Locally coherent costrings and quotients Let $`M^2`$ be any sublattice of $`M`$. Consider the exact sequence $$0M^2M\mathrm{@}>\pi >>M^10.$$ If we apply $`\mathrm{Hom}_{}(,^{})`$ to it, then we get the exact sequence $$1\mathrm{Hom}_{}(M^1,^{})\mathrm{Hom}_{}(M,^{})\mathrm{Hom}_{}(M^2,^{})1.$$ The group $`G=\mathrm{Hom}_{}(M^1,^{})`$ is the product of the torus $`\mathrm{Hom}_{}(M_{free}^1,^{})`$ and the finite group $`\mathrm{Hom}_{}(M_{tor}^1,^{})`$. Let $$0N^1N\mathrm{@}>\pi ^{}>>N^20$$ be the dual sequence. (We point out that we also use $`\pi `$ ($`\pi ^{}`$) to denote its $``$-linear extension.) We need a couple of lemmas to proceed ###### Lemma 11.1. Let $`\sigma _0\mathrm{\Delta }_0`$ be such that $`\pi ^{}(\sigma _0)=\sigma `$ is strongly convex. Then 1. $`[\sigma _0^{}M]^G[\sigma ^{}M^2]`$; 2. $`A_{\sigma _0}//G`$ exists and is naturally identified with $`A_\sigma `$. Proof. Assertion (2) obviously follows from (1). We only need to prove (1). On one hand, it follows from the definitions that $`[\sigma ^{}M^2][\sigma _0^{}M]^G`$. On the other hand, any point of $`[M]`$ that is left invariant under the action of $`G=\mathrm{Hom}_{}(M^1,^{})`$ must belong to $`M_2`$ (one can check this by the action formula given in §11). This leads to $`[\sigma _0^{}M]^G\sigma ^{}M^2`$. The lemma thus follows. The lemma below gives a criterion for the separateness of a quotient ###### Lemma 11.2. Let $`\sigma `$ and $`\sigma ^{}`$ be two cones in $`\mathrm{\Delta }_0`$. The open subset $`A_\sigma A_\sigma ^{}`$ has the separated categorical quotient $`(A_\sigma A_\sigma ^{})//G`$ if and only if $`\pi ^{}(\sigma )\pi ^{}(\sigma ^{})`$ is a face of each. Proof. By Lemma 11.1, the (possibly non-separated) quotient variety $`A_\sigma A_\sigma ^{}//G`$ is naturally identified with the variety $`A_{\pi ^{}(\sigma )}A_{\pi ^{}(\sigma ^{})}`$. By a standard fact from the theory of toric varieties (see Lemma 1.4 of ), this is separated if and only if $`\pi ^{}(\sigma )\pi ^{}(\sigma ^{})`$ is a common face of both. ###### Definition 11.3. A $`G`$-equivariant map $`XY`$ is called a good quotient if $`X`$ admits a covering by invariant affine open subsets $`U_i=\mathrm{Spec}(S_i)`$ such that $`Y=_iU_i//G`$ where $`U_i//G=\mathrm{Spec}(S_i^G)`$. ###### Definition 11.4. A locally coherent costring of $`\pi ^{}:NN^2`$ is called strongly convex if every cone $`\pi ^{}(\sigma )`$ ($`\sigma \mathrm{\Delta }`$) is strongly convex. Recall that a locally coherent costring is tight (minimal) if the dimension of any cone in the locally coherent costring does not drop under the projection. Here comes the main theorem of this section. ###### Theorem 11.5. Let $`\mathrm{\Delta }`$ be a strongly convex locally coherent costing of $`NN^2`$. Then $`U(\mathrm{\Delta })=_{\sigma \mathrm{\Delta }}A_\sigma X`$ is an invariant Zariski open subset such that 1. the separated categorical good quotient $`U(\mathrm{\Delta })U(\mathrm{\Delta })//G`$ exists. 2. $`U(\mathrm{\Delta })//G`$ is the toric variety $`X_{\pi ^{}(\mathrm{\Delta })}`$ defined by the induced fan $`\pi ^{}(\mathrm{\Delta })`$. 3. $`U(\mathrm{\Delta })U(\mathrm{\Delta })//G`$ is a geometric quotient if and only if $`\mathrm{\Delta }`$ is tight (minimal). Proof. That $`U(\mathrm{\Delta })`$ is an invariant Zariski open subset is obvious. (1) and (2) follow from the combination of Lemmas 11.1 and 11.2 and the construction of toric varieties. For much of it, it suffices to observe that $`U(\mathrm{\Delta })//G=_{\sigma \mathrm{\Delta }}A_\sigma //G=_{\sigma \mathrm{\Delta }}A_{\pi ^{}(\sigma )}`$ (by Lemma 12.1 (2)). Statement (3) is local. So to show it, it suffices to prove that $`A_\sigma //G`$ is a geometric quotient for each $`\sigma \mathrm{\Delta }`$. $`A_\sigma //G`$ being geometric is equivalent to that for every $`xA_\sigma `$, $`G_x`$ is finite. Let $`N_\sigma `$ be the lattice generated by $`\sigma `$. Then by §3.1 (or §2.1) of Fulton , $`T_{N_\sigma }=N_\sigma ^{}`$ is the identity component of the isotropy subgroup of any point in the $`T`$-orbit $`𝒪_\sigma A_\sigma `$ determined by $`\sigma `$. Moreover, the identity component of the isotropy subgroup of any point in $`A_\sigma `$ is contained in $`T_{N_\sigma }`$. So a point in $`A_\sigma `$ having isotropy subgroup in $`G=T_{N^1}`$ of positive dimension if and only if $`dimT_{N^1}T_{N_\sigma }>0`$ if and only if $`N^1N_\sigma \{0\}`$ if and only if $`\sigma `$ drops dimension under the projection to $`N^2`$. This proves the theorem. As an immediate corollary of Theorem 11.5 (3), we have ###### Corollary 11.6. If $`\mathrm{\Delta }`$ is tight, then $`U(\mathrm{\Delta })`$ is the toric variety which is a principal $`G`$-bundle over $`U(\mathrm{\Delta })//G`$, assuming that $`G`$ acts on $`U(\mathrm{\Delta })`$ freely. Moreover, all principal $`G`$-bundle among toric varieties arise this way. ###### Remark 11.7. As in the very general case of a quotient theory, degenerate quotients exist. These are the quotients whose dimensions are less than expected (that is, smaller than $`dimXdimG`$). In our toric situation, this will be the case if the induced fan $`\pi ^{}(\mathrm{\Delta })`$ is not strongly convex (that is, some projection $`\pi ^{}(\sigma )`$ contains a non-trivial vector space, or equivalently, $`\sigma N^1`$ is infinite). In this case, $`U(\mathrm{\Delta })//G`$ is the toric variety defined by a reduction fan $`\pi ^{}(\mathrm{\Delta })_{red}`$ and is of dimension less than $`dimXdimG`$. ## 12. Virtual chambers and quotients Assume now that $`X`$ is projective and equipped with a $`T`$-linearized ample line bundle $`L`$. This corresponds to a polytope $`P`$ in the dual lattice $`M`$ of $`N`$. Let $`Q=\pi (P)`$ ###### Corollary 12.1. Let $`\stackrel{~}{𝐜}`$ be a virtual (real) cell of the projection $`\pi :PQ`$. Then it defines a quotient toric variety $`X_{\stackrel{~}{𝐜}}`$ whose defining fan is isomorphic to the fan in $`(\mathrm{ker}\pi )^{}`$ induced by the images of the locally (global) coherent costing $`\mathrm{\Delta }(\stackrel{~}{𝐜})`$. Proof. This is a special case of Theorem 11.5. ## 13. Cox’s quotient construction of toric varieties Let $`\mathrm{\Delta }`$ be any fan in $`N`$ and $`X_\mathrm{\Delta }`$ be the corresponding toric variety. Cox has given a construction of $`X_\mathrm{\Delta }`$ as the quotient of an open subset of $`^{\|\mathrm{\Delta }(1)\|}`$ where $`\mathrm{\Delta }(1)`$ is the set of 1-dimensional cones of $`\mathrm{\Delta }`$ (). $`^{\|\mathrm{\Delta }(1)}\|`$ comes equipped with a natural basis $`\{e_\rho \}`$ where $`\rho `$ is the first lattice point on an edge of $`\mathrm{\Delta }(1)`$. There is canonical projection $`\pi ^{}`$ from $`^{\|\mathrm{\Delta }(1)\|}`$ to $`N`$ defined by sending each $`e_\rho `$ to $`\rho `$. Let $`G=\mathrm{Hom}_{}(A_{n1}(X_\mathrm{\Delta }),^{})`$. ###### Corollary 13.1. (Cox ) Let the notations be as above. 1. $`\mathrm{\Delta }`$ canonically corresponds to a locally coherent costing of $`\pi :^{\|\mathrm{\Delta }(1)\|}N`$ and the toric variety $`X_\mathrm{\Delta }`$ is the quotient of $`^{\|\mathrm{\Delta }(1)\|}`$ by the group $`G`$ defined by the above locally coherent costing. 2. $`X_\mathrm{\Delta }`$ is a geometric quotient if and only if $`\mathrm{\Delta }`$ is simplicial. Proof. 1. This is tautological. Given $`\sigma \mathrm{\Delta }`$, let $`\sigma (1)`$ be the set of edges of $`\sigma `$. Clearly $`\{\mathrm{span}_{\rho \sigma (1)}\{e_\rho \}\|\sigma \mathrm{\Delta }\}`$ is the desired locally coherent string. 2. By Theorem 11.5 (3), $`X_\mathrm{\Delta }`$ is geometric if and only if $`\mathrm{\Delta }`$ is tight. It is tight if and only if the projection $`\mathrm{span}_{\rho \sigma (1)}\{e_\rho \}\mathrm{span}_{\rho \sigma (1)}\{\rho \}`$ does not drop dimension. This can happen if and only if $`\mathrm{span}_{\rho \sigma }\{\rho \}`$ is itself simplicial (i.e., $`\{\rho \}_{\rho \sigma }`$ are linearly independent). ## 14. Quotients defined by realizable virtual chambers Realizable virtual cells (i.e., deformed cells) are special. In this section, we shall explain that their corresponding quotients share some characteristics of projective quotients defined by real cells. We will freely adopt the notations from §9. Let $`𝐜_t`$ and $`𝐜_t^{}`$ be two adjacent cells in $`\mathrm{\Gamma }_t`$ such that the intersection $$𝐜_t𝐜_t^{}=𝐜_t^0$$ is a face of each. Then we have $$U(𝐜_t)U(𝐜_t^0)U(𝐜_t^{})$$ and the inclusions induce a diagram of morphisms of algebraic varieties $$U(𝐜_t)//GU(𝐜_t^{})//G$$ $$ff^{}$$ $$U(𝐜_t^0)//G.$$ Set $`\mathrm{\Sigma }_0`$ to be $`X^{\mathrm{ss}}(𝐜_t^0)//GX^\mathrm{s}(𝐜_t^0)//G`$. Then $`\mathrm{\Sigma }_0`$ admits a stratification by the so-called orbit types. Two points of $`X`$ have the same orbit type if their stabilizers in $`G`$ are identical. This induces a $`G`$-invariant stratification of $`X`$ which in turn induces a stratification of $`\mathrm{\Sigma }_0`$. As in the projective case (e.g., $`t=0`$), we have ###### Theorem 14.1. Assume that $`𝐜_t^0`$ is of codimension 1 and intersects with the interior of $`\mathrm{\Gamma }_t`$. Then (i) $`f`$ and $`f^{}`$ are isomorphisms over the complement to $`\mathrm{\Sigma }_0`$; (ii) over each connected component $`\mathrm{\Sigma }_0^{}`$ of a stratum of $`\mathrm{\Sigma }_0`$, each fiber of $`f`$ ($`f^{}`$) is isomorphic to a quotient of a weighted projective space of dimension $`d`$ ($`d^{}`$) by the finite group $`\pi _0(G_z)`$ where $`z`$ is some point in $`X^{\mathrm{ss}}(𝐜_t^0)X^\mathrm{s}(𝐜_t^0)`$; (iii) $`d+d^{}+1=\mathrm{codim}\mathrm{\Sigma }_0^{}`$. Proof. It is identical to the proof of Theorem 2.2 of . ###### Remark 14.2. For simplicity, we shall call the above birational transformation a Mori-type flips. Here we firmly give a few warnings about the use of the term Mori-type flips. 1. Strictly, Mori flips are defined in the category of projective varieties. We borrow the term in the non-projective case as well. 2. In a some special cases, some of our Mori-type flips may just correspond to blowups, while strictly Mori flips are birational transformations that are isomorphic in codimension 1. ###### Proposition 14.3. Assume that $`𝐜_t^0`$ is of codimension 1, lies on the boundary of $`\mathrm{\Gamma }_t`$, and $`𝐜_t`$ is a chamber containing $`𝐜_t^0`$ as a face. Then the map $`f:X^{\mathrm{ss}}(𝐜_t)//GX^{\mathrm{ss}}(𝐜_t^0)//G`$ is fiber bundle whose typical fiber is isomorphic to a quotient of a weighted projective space by a finite abelian group. Proof. It is identical to the proof of Theorem 2.1 of (the finite abelian group action was unfortunately overlooked there). ###### Corollary 14.4. Every lexicographic (virtual) cell corresponds to a projective quotient which is, modulo finite abelian group actions, a tower of weighted projective bundles over a fixed point variety. It follows then ###### Corollary 14.5. Every pair of quotient algebraic varieties defined by realizable virtual cells are connected by a sequence of Mori-type flips. Proof. Assume that the two quotients are defined by a cell in $`\mathrm{\Gamma }_t`$ and a cell $`\mathrm{\Gamma }_s`$, respectively. Both quotients can be flipped to the same projective quotient defined by a lexicographic cell by crossing (virtual) walls. ## 15. Some notes related to Chow quotient Besides the invariant theoretic quotients, there is a canonical “quotient” space, the Chow quotient. To recall the construction (from ), note that the closure $`\overline{Gx}`$ is a projective subvariety, and as $`x`$ ranges within an $`G`$-invariant open subset $`UX`$ of generic points, these varieties will have the same dimension and degree. Let $`\mathrm{Chow}(X)`$ be the Chow variety of all algebraic cycles in $`X`$. The assignment $`Gx\overline{Gx}`$ defines an embedding of $`U/G`$ into $`\mathrm{Chow}(X)`$. The closure of the image of $`U/G`$ in $`\mathrm{Chow}(X)`$ is the Chow quotient and denoted by $`X//^{ch}G`$. The Chow quotient is a toric variety of the residue torus and its corresponding fan is the normal fan of the fiber polytope $`\mathrm{\Sigma }(P,Q)`$ (). Assume now that we are in the situation as in §13. Then the $`T`$-linearized an ample line bundle over $`X`$ determines a unique moment map $`\mu _T`$ for the $`T`$-action whose total image is the polytope $`P`$. It is a standard fact that for the image of $`\overline{𝒪}_x=\overline{Tx}`$ under $`\mu _T`$ is the face of $`P`$. Also, with respect to the same ample line bundle $`L`$, a moment map of the subtorus $`G`$-action is the composition $$\mu _G:XPQ$$ where the last map is the natural projection we began with. Let $`C=_ia_iC_iX//^{ch}G`$ where $`C_i=\overline{Gx_i}`$ for some $`x_iX`$. Then ###### Proposition 15.1. $`\{\mu _T(\overline{𝒪}_{x_i})\}_i`$ is a coherent string of $`PQ`$. Moreover, every coherent string can be realized this way. Proof. This is a reformulation of Proposition 3.6 of . It is known that the Chow quotient maps to every projective GIT quotient. This fact extends to non-projective cases as well. ###### Corollary 15.2. Given any locally coherent costing $`\mathrm{\Delta }`$, there is a natural projection from the Chow quotient $`X//^{ch}G`$ to the quotient variety $`U(\mathrm{\Delta })//G`$ that is a birational $`T/G`$-toric morphism. Proof. In the projective category, see the proofs in and . In our special situation, the following proof is quick: just observe from the definitions that the fan of the Chow quotient refines the fan induced by any locally coherent costing. The birationality and the $`T/G`$-equivariancy is obvious. ## 16. Bistellar flips and Mori-type flips Let $`𝒜=\{a_0,\mathrm{},a_n\}`$ be a finite lattice point set in $`^d^d`$ and $`Q=conv(𝒜)`$. We may consider $`Q`$ as the projection of the standard n-simplex $`P`$ by a suitable linear map $`\pi `$. In this case, the minimal elements of ($`𝒯_{coh}`$) $`𝒯`$ are the (coherent) triangulations of $`Q`$ using solely the vertices from $`𝒜`$. Coherent triangulations correspond to the vertices of the secondary polytope $`\mathrm{\Sigma }(𝒜)`$ of $`𝒜`$. Two vertices are joint by an edges of $`\mathrm{\Sigma }(𝒜)`$ if and only if the corresponding triangulations are related by a bistellar flips. Here a bistellar flip is a local re-arrangement of a triangulation without introducing new vertices (see ). Embed $`𝒜^d^{d+1}`$ in the affine hyperplane of height 1 ($`x_{d+1}=1`$). Given any triangulation $`T𝒯`$, taking the cone over $`T`$, we get a fan $`\mathrm{\Delta }_T^0`$ with the support $`cone(𝒜)`$. To get a complete fan $`\mathrm{\Delta }_T`$, we add the ray $`R`$ generated by the lattice point $`0(1)^d`$. The cones in $`\mathrm{\Delta }_T\mathrm{\Delta }_T^0`$ are generated by cones in $`(cone(𝒜))`$ and the ray $`R`$. ###### Proposition 16.1. $`T`$ is coherent if and only if $`\mathrm{\Delta }_T`$ is projective. Proof. The proof is almost tautological. Note that ($`T`$ is coherent) $`\mathrm{\Delta }_T`$ is projective if ($`T`$) $`\mathrm{\Delta }_T`$ possesses a strictly convex piece-wise (affine) linear function. Thus if $`\mathrm{\Delta }_T`$ possesses such a function, restricting it to $`T`$ will give a desired function for $`T`$. If $`T`$ has such a function, because $`\mathrm{\Delta }_T`$ is obtained from the cone over $`T`$ be adding one edge $`(0(1))`$, this function can be extended to be a strictly convex piece-wise linear function on $`\mathrm{\Delta }_T`$ by assigning to $`0(1)`$ a generic number (to assure that the extended linear functions are all different). ###### Example 16.2. Figure 4 is a non-coherent triangulation $`T`$ () and its corresponding non-projective fan $`\mathrm{\Delta }_T`$. (cf. for another example of non-projective fans.) Figure 4 Let $`𝒯_{coh}^0`$ be the subset of coherent triangulations that actually use all vertices of $`𝒜`$. Set $$G_{\widehat{𝒜}}=\mathrm{Hom}_{}(A_d(X_{\mathrm{\Delta }_T}),^{})$$ for any $`T𝒯_{coh}^0`$. One checks that the group does not depend on the choice of $`T`$. The group $`G_{\widehat{𝒜}}`$ acts on the complex vector space $`()^{\|𝒜\|+1}`$ via the inclusion $`G_{\widehat{𝒜}}(^{})^{\|𝒜\|+1}.`$ Now by our earlier results, we have ###### Proposition 16.3. 1. There is canonical correspondence between the triangulations of $`Q`$ and quotient varieties of $`^{\|𝒜\|+1}`$ by the group $`G_{\widehat{𝒜}}`$; 2. under this correspondence, coherent triangulations correspond to projective quotient varieties, not-necessarily-coherent triangulations correspond not-necessarily-projective quotient varieties; 3. a bistellar flip between two triangulations transports to a Mori-type flip<sup>2</sup><sup>2</sup>2See Remark 14.2. between their corresponding quotients. ## 17. Notes on virtual (pseudo) oriented matroids This section is loosely related to the rest of the paper, especially §§4–7. It raises more questions than it attempts to solve. An oriented matroid of rank $`n`$ correspond to a (tamed) topological deformation of an arrangement of hyperplanes in $`^n`$ with that a realizable one corresponds to the hyperplane arrangement itself. Given a hyperplane arrangement, it induces a (special) fan $`\mathrm{\Delta }`$ (a normal fan of a zonotope) in $`^n`$. For each hyperplane $`e`$, we assign an orientation: naming one side positive, the other negative. Thus each cocell $`\sigma `$ (i.e., a cone in $`\mathrm{\Delta }`$) corresponds to a sign vector $`v^\sigma `$ in $`\{+,0,\}^N`$ according to the location of the cocell relative to each and every hyperplane $`e`$, where $`N`$ is the number of hyperplanes. So, $`v_e^\sigma `$ assumes one of the values $`+,0,`$ according to one of the three positions of the cocell $`\sigma `$ relative to the hyperplane $`e`$. The oriented matroid defined by $`\mathrm{\Delta }`$ is a combinatorial abstraction of properties of those sign vectors. We refer the reader to for a list of axioms that characterizes the system of the sign vectors $`\{v^\sigma \|\sigma \mathrm{\Delta }\}`$, which are formal (and equivalent) definitions of an oriented matroid. We use $`(\mathrm{\Delta })`$ to denote this oriented matroid. Such a $`(\mathrm{\Delta })`$ is called a realizable oriented matroid. If we deform $`\mathrm{\Delta }`$ to get $`\mathrm{\Delta }_t`$, the same strategy will still give us a sign vector for every deformed cocell. Hence it also leads to a (possibly non-realizable) oriented matroid $`(\mathrm{\Delta }_t)`$. In the above, we only consider the normal fan of a zonotope, a linear projection of a standard cube. What about the normal fan of any polytope? Or more generally, what about an arbitrary fan? For an arbitrary complete fan $`\mathrm{\Delta }`$, it can always be extended to be the fan $`\widehat{\mathrm{\Delta }}`$ induced by an arrangement of hyperplanes: just take the linear spans of the cones of codimension one. This extension is canonical. So we still get a set of hyperplanes and still can give orientations of these hyperplanes. Given a cone $`\sigma `$ in $`\mathrm{\Delta }`$, it may take a definite position relative to some hyperplane $`e`$: + side, $``$ side, or 0 (lying on the hyperplane). For some other hyperplane $`e`$, $`e`$ runs though $`\sigma `$ in the relative interiors so that $`\sigma `$ does not take any definite position relative to $`e`$. We may say that for $`\sigma `$, $`e`$ is a “ghost” hyperplane. Should this be problematic? What if we take this situation as a “position” too? So, let us use the first letter in “undecided” to indicate this situation. Then a cone $`\sigma `$ gives rise to a generalized sign vector $`v^\sigma `$ in $$\{+,0,,u\}^N\{(u,\mathrm{},u)\}.$$ That is, for any hyperplane $`e`$, $`v_e^\sigma `$ assumes one of the values $`+,0,,u`$ according to one of the above four positions. Since $`\mathrm{\Delta }`$ is not empty, for every given $`\sigma `$, $`v_e^\sigma u`$ for some $`e`$, so that $`v^\sigma `$ does not assume the value $`\{(u,\mathrm{},u)\}.`$ Now we may call the combinatorial abstraction of the system of these generalized sign vectors a virtual (or pseudo) oriented matroid, still denoted by $`(\mathrm{\Delta })`$. Take any (tamed) topological deformation $`\mathrm{\Delta }_t`$ of $`\mathrm{\Delta }`$, we can assign a virtual oriented matroid $`(\mathrm{\Delta }_t)`$ to $`\mathrm{\Delta }_t`$ in the same way as we did for $`\mathrm{\Delta }`$ itself. We may call the virtual oriented matroid $`(\mathrm{\Delta })`$ a realizable one in accordance with the case of oriented matroids. Each generalized sign vector $`v`$ in $`\{+,0,,u\}^N`$ splits into many (real) sign vectors: whenever $`v_e=u`$ we can let it to split into to $`\widehat{v}_e=+,0,`$. So, $`v_e`$ splits into $`3^{\chi (v)}`$ many (real) sign vectors where $`\chi (v)`$ is the number of hyperplanes such that $`v_e=u`$. The system of $`\{\widehat{v}_e\}`$ is nothing but the system of the sign vectors for $`(\widehat{\mathrm{\Delta }})`$, where $`\widehat{\mathrm{\Delta }}`$ is the canonical extension of $`\mathrm{\Delta }`$ as explained earlier. We call $`(\widehat{\mathrm{\Delta }})`$ the oriented matroid extension of $`(\mathrm{\Delta })`$. This extension is canonical. Hence given any system $``$ of generalized sign vectors, the above splitting procedure leads to a canonical system $`\widehat{}`$ of (real) sign vectors. Again, $`\widehat{}`$ is called the canonical extension of $``$. From this, we can give a practical definition of virtual oriented matroid as follows. ###### Definition 17.1. A set $``$ of generalized sign vectors in $$\{+,0,,u\}^N\{(u,\mathrm{},u)\}$$ is called a virtual (or pseudo) oriented matroid of rank $`n`$ if its canonical extension $`\widehat{}`$ is an oriented matroid of rank $`n`$. An oriented matroid admits an involutive automorphism $$I:vv.$$ A virtual oriented matroid does not have this property. In fact, if it does, it must be equal to its canonical extension, hence itself an oriented matroid. In this sense, an oriented matroid is a virtual oriented matroid equipped with an involution. Put it slightly differently, a virtual oriented matroid is an “oriented matroid without involution”. The following diagram should be interesting, which would extend the topological representation theorem of oriented matroids to the virtual situation. \| {deformed fans} \| {virtual oriented matroids} \| \| --- \| --- \| \| $``$ \| $``$ \| \| {fans} \| {realizable virtual oriented matroids} \| \| $``$ \| $``$ \| \| {hyperplane arrangements} \| {realizable oriented matroids} \| \| $``$ \| $``$ \| \| {deformed hyperplane arrangements} \| {oriented matroids} \| \| $``$ \| $``$ \| \| {deformed fans} \| {virtual oriented matroids}. \| We end our venture to virtual oriented matroids as follows. ###### Conjecture 17.2. With the help of the canonical oriented matroid extension, we may list a system of axioms which characterizes a virtual oriented matroid in a way analogous to the case where a system of axioms characterizes an oriented matroid. ###### Problem 17.3. Modify as many results of as possible so that the revised results are valid in the virtual situation. Moreover, find some interesting applications of the new theory. Finally, we would also like to ask ###### Question 17.4. Do virtual oriented matroids provide geometric structures for some combinatorial spaces in analogous to that oriented matroids provide differentiable structures for MacPherson’s combinatorial differential manifolds ?
warning/0004/math0004145.html	ar5iv	text	# Compact Range Property and Operators on 𝑪^∗-algebras ## 1. INTRODUCTION A Banach space $`E`$ is said to have the compact range property (CRP) if every $`E`$-valued countably additive measure of bounded variation has compact range. It is well known that every Banach space with the Radon-Nikodym property (RNP) has the (CRP) and for dual Banach spaces, the (CRP) were completely characterized as those whose predual do not contain any copies $`\mathrm{}^1`$. For more in depth discussions on Banach spaces with the (CRP), we refer to . The following characterization can be found in : A Banach space $`E`$ has the (CRP) if and only if every $`1`$-summing operator from $`C[0,1]`$ into $`E`$ is compact. Since $`C[0,1]`$ is a (commutative) $`C^{}`$-algebra, it is a natural question whether $`C[0,1]`$ can be replaced by any $`C^{}`$-algebras. Let us recall that in , it was shown that if $`X`$ is a Banach space that does not contain any copies of $`\mathrm{}^1`$ then any $`1`$-summing operators from any given $`C^{}`$-algebra into $`X^{}`$ is compact; hinting that, as in commutative case, the (CRP) is the right condition to provide compactness. The present note is an improvement of . Our main result confirms that, if $`𝒜`$ is a $`C^{}`$-algebra and $`E`$ is a Banach space that has the (CRP) then every $`1`$-summing from $`𝒜`$ into $`E`$ is compact. Our proof relies on factorizations of summing operators used in and properties of integral operators. There is another well kown characterization of spaces with the (CRP) in terms of operators defined on $`L^1[0,1]`$: a Banach space $`E`$ has the $`(CRP)`$ if and only every operator $`T`$ from $`L^1[0,1]`$ into $`E`$ is Dunford Pettis (completely continuous) thus the (CRP) is also referred to as the complete continuity property (CCP). Unlike the $`1`$-summing operators on $`C^{}`$-algebras, operators defined on non-commutative $`L^1`$-spaces do not behave the same way as those defined on $`L^1[0,1]`$ do. In the last section of this note, we will discuss these operators along with $`C^{}`$-summing operators studied by Pisier in . Our terminology and notation are standard as may be found in and for Banach spaces, and for $`C^{}`$-algebras and operator algebras. ## 2. PRELIMINARIES In this section, we recall some definitions. ###### Definition 1. Let $`X`$ and $`Y`$ be Banach spaces and $`0<p<\mathrm{}`$. An operator $`T:XY`$ is said to be $`p`$-summing if there is a constant $`C`$ such that for any finite sequence $`(x_1,x_2,...,x_n)`$ of $`X`$, one has $$(\underset{i=1}{\overset{n}{}}Tx_i^p)^{\frac{1}{p}}Csup\{(\underset{i=1}{\overset{n}{}}\|x_i,x^{}\|^p)^{\frac{1}{p}};x^{}X^{},x^{}1\}.$$ The smallest constant $`C`$ for which the above inequality holds is denoted by $`\pi _p(T)`$ and is called the $`p`$-summing norm of $`T`$. ###### Definition 2. We say that an operator $`T:XY`$ is an integral operator if it admits a factorization: where $`i`$ is the natural inclusion from $`Y`$ into $`Y^{},\mu `$ is a probability measure on a compact space $`K`$, $`J`$ is the natural inclusion and $`\alpha `$ and $`\beta `$ are bounded linear operators. We define the integral norm $`i(T):=inf\left\{\alpha \beta \right\}`$ where the infinum is taken over all such factorizations. Similarily, we shall say that $`T`$ is strictly integral if $`T`$ is integral and on the factorization above $`\beta `$ takes its values in $`Y`$. It is well known that integral operators are $`1`$-summing but the converse is not true. If $`X=C(K)`$ where $`K`$ is a compact Hausdorff space then it is well known that every $`1`$-summing operator from $`X`$ into $`Y`$ is integral. For more details on the different properties of the classes of operators involved, we refer to . The following simple fact will be needed in the sequel. ###### Proposition 3. Let $`T:XY`$ be a strictly integral operator. If $`Y`$ has the (CRP) then $`T`$ is compact. ###### Proof. The operator $`T`$ has a factorization $`T=\beta J\alpha `$ where $`\alpha :XL^{\mathrm{}}(\mu )`$, $`J:L^{\mathrm{}}(\mu )L^1(\mu )`$ and $`\beta :L^1(\mu )Y`$ are as in the above definition. Note that $`J`$ is $`1`$-summing so $`\beta J:L^{\mathrm{}}(\mu )Y`$ is $`1`$-summing and since $`L^{\mathrm{}}(\mu )`$ is a $`C(K)`$-space and $`Y`$ has the (CRP), $`\beta J`$ (and hence $`T`$) is compact. ∎ We recall that a von Neumann algebra $``$ is said to be $`\sigma `$-finite if the identity is countably decomposable equivalently if there exist a faithful state $`\phi _{}`$. As is customary, for every functional $`\phi _{}`$ and $`x`$, $`x\phi `$ (resp. $`\phi x)`$ denotes the normal functional $`y\phi (yx)`$ (resp. $`y\phi (xy))`$. ## 3. MAIN RESULT ###### Theorem 4. For a Banach space $`E`$, the following are equivalent: * $`E`$ has the CRP; * Every $`1`$-summing operator $`T:C[0,1]E`$ is compact; * For any given $`C^{}`$-algebra $`𝒜`$, every $`1`$-summing operator $`T:𝒜E`$ is compact. The equivalence $`(1)(2)`$ is well known, we refer to , for more details. Clearly (3) $``$ (2) so what we need to show is (1) $``$ (3). For this, it is enough to consider the following particular case (see for this reduction). ###### Proposition 5. Let $`E`$ be a Banach space with the (CRP) and $``$ be a $`\sigma `$-finite von Neumann algebra. If $`T:E`$ is $`1`$-summing and is weak to weakly continuous then $`T`$ is compact. ###### Proof. The proof is a refinement of the argument used in Proposition 3.2 of . We will include most of the details for completeness. Without loss of generality, we can and do assume that $`E`$ is separable. Let $`\delta >0`$. From Lemma 2.3 of , $$Tx2(1+\delta )\pi _1(T)xf+fx_{_{}}\text{for every}x,$$ where $`f`$ is a faithful normal state in $`_{}`$. If $`L^2(f)`$ is completion of the prehilbertian space $`(,,)`$ where $`x,y=f({\displaystyle \frac{xy^{}+y^{}x}{2}})`$ then we have the following factorization: where $`J`$ is the inclusion map, $`\theta (Jx)=,J(x^{})`$ for every $`x`$ and $`L(\frac{xf+fx}{2})=Tx`$. We recall that $`L`$ is a well defined bounded linear map since $`\{xf+fx;x\}`$ is dense in $`_{}`$ and $`L(xf+fx)4(1+\delta )\pi _1(T)xf+fx_{_{}}`$. Let $`S:=J^{}\theta J`$. Claim: $`JL^{}:E^{}L^2(f)`$ is compact. For this, let us consider $`L^{}:E^{}`$. Since $`E`$ is separable, it is isometric to a subspace of $`C[0,1]`$. Let $`I_E`$ be the isometric embedding of $`E`$ in $`C[0,1]`$ and $`i`$ be the natural inclusion of $`C[0,1]`$ into $`C[0,1]^{}`$. Define the following map $`\stackrel{~}{T}`$ from $``$ into $`C[0,1]`$ by setting $`\stackrel{~}{T}=\overline{I_ET(x^{})}`$ for every $`x`$. (Here, $`\overline{f}`$ is the map $`t\overline{f(t)}`$ for $`fC[0,1]`$ with $`\overline{f(t)}`$ being the conjugate of the complex number $`f(t)`$). Clearly, $`\stackrel{~}{T}`$ is linear and bounded and it can be shown that $`\stackrel{~}{T}`$ is $`1`$-summing and is weak to weakly continuous. In fact, if $`(x_1,x_2,...,x_n)`$ is a finite sequence in $``$ then $$\begin{array}{cc}\hfill \underset{i=1}{\overset{n}{}}\stackrel{~}{T}x_i& =\underset{i=1}{\overset{n}{}}\overline{I_ET(x_i^{})}\hfill \\ & =\underset{i=1}{\overset{n}{}}I_ET(x_i^{})\hfill \\ & =\underset{i=1}{\overset{n}{}}T(x_i^{})\hfill \\ & \pi _1(T)sup\{\underset{i=1}{\overset{n}{}}\|x_1^{},\phi \|,\phi ^{},\phi 1\}\hfill \\ & \pi _1(T)sup\{\underset{i=1}{\overset{n}{}}\|x_i,\phi ^{}\|,\phi ^{},\phi 1\}\hfill \end{array}$$ so $`\stackrel{~}{T}`$ is $`1`$-summing with $`\pi _1(\stackrel{~}{T})\pi _1(T)`$. Moreover if $`(x_\alpha )_\alpha `$ is a net that converges to zero weak in $``$ so does the net $`(x_\alpha ^{})_\alpha `$ and since $`T`$ is weak to weakly continuous, $`(T(x_\alpha ^{}))_\alpha `$ converges to zero weakly in $`E`$ and hence $`(\stackrel{~}{T}(x_\alpha ))_\alpha `$ is weakly null which shows that $`\stackrel{~}{T}`$ is weak to weakly continuous. To complete the proof, consider Since $`C[0,1]^{}`$ has the Hahn-Banach extension property and $`i\stackrel{~}{T}`$ is $`1`$-summing, $`i\stackrel{~}{T}`$ is an integral operator. Let $`K:C[0,1]^{}_{}`$ such that $`K^{}=i\stackrel{~}{T}`$ (such operator exists since $`i\stackrel{~}{T}`$ is weak to weakly continuous); $`K`$ is integral () and since $`_{}`$ is a complemented subspace of its bidual $`^{}`$ (see for instance ), $`K`$ is strictly integral and therefore $`LK:C[0,1]^{}E`$ is strictly integral and by Proposition 3, $`LK`$ (and hence $`(LK)^{}=i\stackrel{~}{T}L^{}`$) is compact. Let $`(U_n)`$ be a bounded sequence in $`E^{}`$. There exists a subsequence $`(U_{n_k})`$ so that $`(i\stackrel{~}{T}L^{}(U_{n_k}))_k`$ is norm convergent in $`C[0,1]^{}`$. Since $`i`$ and $`I_E`$ are isometries, we get that $`(TL^{}(U_{n_k}))_k`$ is norm convergent so $`lim_{k,m}T(L^{}(U_{n_k})^{})T(L^{}(U_{n_m})^{})=0`$. As in , we get $$\underset{k,m}{lim}T(L^{}(U_{n_k})^{})T(L^{}(U_{n_m})^{}),U_{n_k}U_{n_m}=\underset{k,m}{lim}JL^{}(U_{n_k}U_{n_m})_{L^2(f)}^2=0$$ which proves that $`\left(JL^{}(U_{n_k})\right)_k`$ is norm-convergent in $`L^2(f)`$. The proof is complete. ∎ ###### Theorem 6. Let $`𝒜`$ be a $`C^{}`$-algebra, $`E`$ be a Banach space and $`0<p<1`$. Every $`p`$-summing operator from $`𝒜`$ into $`E`$ is compact. ###### Proof. Let $`T:𝒜E`$ be an operator with $`\pi _p(T)<\mathrm{}`$. One can choose, by the Pietsch Factorization Theorem, a probability space $`(\mathrm{\Omega },\mathrm{\Sigma },\mu )`$ such that where $`S`$ is a subspace of $`L^{\mathrm{}}(\mu )`$, $`S_p`$ is the closure of $`S`$ in $`L^p(\mu )`$ and $`i_p`$ is the restriction of the natural inclusion $`j_p`$. Denote by $`S_1`$ the closure of $`S`$ in $`L^1(\mu )`$, by $`i_1`$ the restriction of the natural inclusion and $`i_{1,p}`$ the natural inclusion of $`S_1`$ into $`S_p`$. Claim: $`\stackrel{~}{T}i_{1,p}:S_1E`$ is weakly compact. To see this, let $`(f_n)_n`$ be a bounded sequence in $`S_1L^1(\mu )`$. By Komlòs’s Theorem, there exists a subsequence $`(f_{n_k})_k`$ and a function $`fL^1(\mu )`$ such that $`lim_m\mathrm{}\frac{1}{m}_{k=1}^mf_{n_k}(\omega )=f(\omega )`$ for a.e. $`\omega \mathrm{\Omega }`$. Since $`0<p<1`$, $$\underset{m\mathrm{}}{lim}\frac{1}{m}\underset{k=1}{\overset{m}{}}f_{n_k}f_p=0.$$ This shows that $`fS_p`$ and $`\left(\stackrel{~}{T}i_{1,p}\left(\frac{1}{m}_{k=1}^mf_{n_k}\right)\right)_m`$ converges to $`\stackrel{~}{T}(f)`$ in $`E`$ and the claim follows. Using the factorization of weakly compact operator , $`i_{1,p}\stackrel{~}{T}`$ factors through a reflexive space and since $`i_1J`$ is $`1`$-summing, the theorem follows from Theorem 4. ∎ ## 4. Concluding remarks Let us recall some definitions ###### Definition 7. Let $`X`$ and $`Y`$ be Banach spaces. An operator $`T:XY`$ is called Dunford-Pettis if $`T`$ sends weakly compact sets into norm compact sets. The following class of operators was introduced by Pisier in as extension of the $`p`$-summing operators in the setting of $`C^{}`$-algebras. ###### Definition 8. Let $`𝒜`$ be a $`C^{}`$-algebra and $`E`$ be a Banach space, $`0<p<\mathrm{}`$. An operator $`T:𝒜E`$ is said to be $`p`$-$`C^{}`$-summing if there exists a constant $`C`$ such that for any finite sequence $`(A_1,\mathrm{},A_n)`$ of Hermitian elements of $`𝒜`$, one has $$\left(\underset{i=1}{\overset{n}{}}T(A_i)^p\right)^{\frac{1}{p}}C\left(\underset{i=1}{\overset{n}{}}\|A_i\|^p\right)^{1/p}_𝒜.$$ Let $``$ be a finite von Neumann algebra with a faithful tracial state $`\tau `$ and let $`J`$ be the canonical inclusion map from $``$ into $`L^1(,\tau )`$. As in the commutative case, we have the following : ###### Proposition 9. Let $`E`$ be a Banach space and $`T:L^1(,\tau )E`$ a bounded linear map. Then the following are equivalent: * $`T`$ is Dunford-Pettis; * $`TJ`$ is compact. ###### Proof. $`(i)(ii)`$ is trivial. For the converse, let $`(a_n)_n`$ be a weakly null sequence in the unit ball of $`L^1(,\tau )`$. It is clear that $`(a_n^{})_n`$ is also weakly null so without loss of generality, we can assume that $`(a_n)_n`$ is a sequence of self-adjoint operators. For each $`n1`$, set $`a_n=_{\mathrm{}}^{\mathrm{}}t𝑑e_t^{(n)}`$ the spectral decomposition of $`a_n`$ and for every $`N1`$, let $$p_{n,N}=_N^N1𝑑e_t^{(n)}.$$ It is clear that for every $`n1`$ and $`N1`$, $$\tau (\mathrm{𝟏}p_{n,N})=\tau (_{\{\|t\|>N\}}1𝑑e_t^{(n)})\frac{1}{N}\tau (\|a_n\|).$$ By the Akeman’s characterization of relatively weakly compact subset in $`L^1(,\tau )`$ (see for instance Theorem 5.4 p.149), we conclude that for any given $`ϵ>0`$, there is $`N_01`$ such that for every $`n1`$, $`a_n(\mathrm{𝟏}p_{n,N_0})ϵ`$. Moreover $`(a_np_{n,N_0})_n`$ is a bounded sequence in $``$ and since $`TJ`$ is compact, there is a compact subset $`K_ϵ`$ of $`E`$ such that $`\{T(a_n);n\}K_ϵ+ϵB_E`$. The proof is complete. Fix a type $`II_1`$ von Neumann algebra $``$ such that $``$ contains a complemented copy of a Hilbert space $`H`$. The space $`H`$ is reflexive (and therefore has (CRP) ) but the projection map $`P`$ from $`L^1(,\tau )`$ onto $`H`$ can not be Dunford-Pettis. A very well known property of $`p`$-summing operators is that they are Dunford-Pettis. This in not the case for $`C^{}`$-summing operators in general. By Proposition 9, $`PJ`$ is not compact. We remark that the argument used in requires only that the operator is $`C^{}`$-summing and Dunford-Pettis hence since $`J`$ is clearly $`C^{}`$-summing and $`PJ`$ is not compact, $`PJ`$ should not be Dunford-Pettis. Acknowledgments. The author wishes to express his thanks to Patrick Dowling for helpful discussions regarding this note.
warning/0004/hep-th0004180.html	ar5iv	text	# References A two-dimensional integrable axionic $`\sigma `$-model and T-duality János Balog Research Institute for Particle and Nuclear Physics, Hungarian Academy of Sciences, H-1525 Budapest 114, P.O.B. 49, Hungary Péter Forgács Laboratoire de Mathématiques et Physique Théorique Université de Tours Parc de Grandmont, 37200 Tours, France László Palla Institute for Theoretical Physics Eötvös University H-1117 Budapest, Pázmány P. sétány 1A, Hungary Abstract An $`S`$-matrix is proposed for the two dimensional O(3) $`\sigma `$-model with a dynamical $`\theta `$-term (axion model). Exploiting an Abelian T-duality transformation connecting the axion model to an integrable SU(2)$`\times `$U(1) symmetric principal $`\sigma `$-model, strong evidence is presented for the correctness of the proposed $`S`$-matrix by comparing the perturbatively calculated free energies with the ones based on the Thermodynamical Bethe Ansatz. This T-duality transformation also leads to a new Lax-pair for both models. The quantum non-integrability of the O(3) $`\sigma `$-model with a constant $`\theta `$-term, in contradistinction to the axion model, is illustrated by calculating the $`23`$ particle production amplitude to lowest order in $`\theta `$. In this paper we shall study the following two dimensional $`\sigma `$-model described by the Lagrangian<sup>1</sup><sup>1</sup>1We use the following conventions. For a vector $`v`$ in two-dimensional Minkowski space $`v^\mu v_\mu =v_+v_{}`$ where $`v_\pm =v_0\pm v_1`$. The antisymmetric tensor is defined by $`ϵ^{01}=1`$, $`\tau ^a=\sigma ^a/2`$ with $`\sigma ^a`$ being the standard Pauli-matrices satisfying $`\sigma ^a\sigma ^b=\delta ^{ab}+iϵ^{abc}\sigma ^c`$.: $$=\frac{1}{2\stackrel{~}{\lambda }}_\mu n^a^\mu n^a+\frac{\stackrel{~}{\lambda }}{32\pi ^2(1+\stackrel{~}{g})}_\mu \theta ^\mu \theta +\frac{\theta }{8\pi }ϵ^{\mu \nu }ϵ^{abc}n^a_\mu n^b_\nu n^c,$$ (1) where $`n^an^a=1`$. In fact (1) is an O(3) non-linear $`\sigma `$-model, coupled to a scalar field, $`\theta `$, (whose normalization has been chosen for later convenience) through the Hopf term. This latter is proportional to the topological current of the O(3) model and with the normalization chosen in (1) its space-time integral (after a Wick rotation) yields the topological charge, which can take integer values only. This implies that the variable $`\theta `$ is actually an angle, taking its values between $`0`$ and $`2\pi `$. This observation will play an important role in all our considerations. We shall refer to (1) as the ‘axion model’ since it can be thought of as the O(3) nonlinear $`\sigma `$-model with a dynamical $`\theta `$-term which can be regarded as a two-dimensional analogue of its phenomenologically important four-dimensional counterpart . The following (heuristic) consideration might be useful to gain some insight into the physics of the axion model. Let us integrate out the O(3) fields, $`n^a`$, in some generating functional of the theory (1). This way one would obtain a non-vanishing effective potential for the $`\theta `$ field. Because $`\theta `$ is $`2\pi `$-periodic the effective potential must be also periodic. The effective theory of the $`\theta `$ field is therefore expected to be similar to the Sine-Gordon model, with a periodic potential and corresponding topological current $`K_\mu =ϵ_{\mu \nu }^\nu \theta /2\pi `$. Note that $`\theta `$ being an angular variable makes it difficult to integrate it out in the functional integral in spite of $``$ being only quadractic in $`\theta `$. The axion model belongs to a family of (classically) integrable two-dimensional non-linear $`\sigma `$-models with an O(3) symmetry discovered in Ref. . All of the models of Ref. have target spaces with non-vanishing torsion (in addition to the metric tensor field) but the axion model is especially simple, as its torsion is constant. It has already been observed in Ref. that the axion model (1) can be mapped to a one parametric deformation of the SU(2) principal $`\sigma `$-model by an Abelian T-duality transformation. This latter (torsionless) model has an SU$`{}_{\mathrm{L}}{}^{}(2)`$U<sub>R</sub>(1) symmetry and recently there has been some revival of interest in it . Its Lagrangian can be written as: $$_\mathrm{\Sigma }=\frac{1}{2\lambda }\left\{L_\mu ^aL^{a\mu }+gL_\mu ^3L^{3\mu }\right\},$$ (2) where $$L_\mu =G^1_\mu G=\tau ^aL_\mu ^a,$$ (3) and $`g`$ is the parameter of deformation. The theory (2) is known to be integrable classically and there is little doubt that it is also quantum-integrable . Its spectrum contains two massive doublets (kinks) whose scattering is described by the tensor product of an SU(2)$`\times `$U(1) symmetric solution of the bootstrap $`S`$-matrix equations: $$S(\theta )=S^{(\mathrm{})}(\theta )S^{(p)}(\theta ),$$ (4) where $`S^{(p)}(\theta )`$ denotes the Sine-Gordon (SG) $`S`$-matrix. Depending on the value of the parameter, $`p`$, in addition to the kinks there are also some bound states (breathers) in the spectrum transforming as $`3+1`$ under SU(2). In the following we shall show that the somewhat unexpected identification between the axion model and the deformed principal $`\sigma `$-model through a T-duality transformation allows us to learn more about both of them. Assuming the validity of the duality transformation at the quantum level between the two theories implies the absence of particle production in the axion model (1), and also that its factorized scattering theory is given by the $`S`$-matrix of Eq. (4). The proposed quantum integrability of the axion model might seem somewhat surprising, as it is generally believed that the O(3) model with a constant $`\theta `$-term is not quantum integrable, except for the special value $`\theta =\pi `$ (despite the fact that the $`\theta `$-term, being a total derivative, does not change the classical physics of the model). We now show that in the framework of the form-factor bootstrap approach the $`\theta `$ term mediates particle production in the O(3) $`\sigma `$-model, indeed. To lowest order in $`\theta `$ the $`23`$ particle production amplitude can be written as $`p,b;p^{},b^{};p^{\prime \prime },b^{\prime \prime }\|q,a;q^{},a^{}_{(\theta )}=(2\pi )^2i\theta \delta ^{(2)}(p+p^{}+p^{\prime \prime }qq^{})`$ $`p,b;p^{},b^{};p^{\prime \prime },b^{\prime \prime }\|T(0)\|q,a;q^{},a^{}_{(0)}+𝒪(\theta ^2),`$ (5) where in the first line the amplitude is in the O(3) model with a $`\theta `$-term, while in the second line the matrix element of the topological charge density operator $`T`$ is to be calculated in the original O(3) $`\sigma `$-model (with $`\theta =0`$). In other words we simply apply perturbation theory in $`\theta `$. Let us now consider the following simplified kinematical configuration: the incoming particles have momenta $`q_1=Q`$ and $`q_1^{}=Q`$, whereas the produced (outgoing) three particles have momenta $`p_1=Q^{}`$, $`p_1^{}=0`$ and $`p_1^{\prime \prime }=Q^{}`$ respectively. Here $`Q^{}`$ can easily be expressed in terms of $`Q`$ and the kink mass $`M`$ using energy conservation. For large $`Q`$, using the results of Ref. , we find $$p,b;p^{},b^{};p^{\prime \prime },b^{\prime \prime }\|T(0)\|q,a;q^{},a^{}_{(0)}\pi ^{\frac{5}{2}}\frac{Q^2}{\mathrm{ln}^3Q/M}\left(ϵ^{a^{}ba}\delta ^{b^{}b^{\prime \prime }}ϵ^{b^{\prime \prime }ba}\delta ^{b^{}a^{}}\right).$$ (6) Eq. (6) shows that already to first order in $`\theta `$, the $`23`$ particle production amplitude is different from zero. Thus at least for small values of $`\theta `$, the introduction of this term destroys quantum integrability of the O(3) $`\sigma `$-model, indeed. To exhibit now the classical T-duality transformation between the two models , we introduce the parametrization $$n^1=\mathrm{sin}\vartheta \mathrm{sin}\phi ,n^2=\mathrm{sin}\vartheta \mathrm{cos}\phi ,n^3=\mathrm{cos}\vartheta ,\theta =\frac{4\pi }{\stackrel{~}{\lambda }}\sqrt{1+\stackrel{~}{g}}\chi ,$$ (7) in terms of which the Lagrangian (1) (after an integration by parts) becomes $$=\frac{1}{2\stackrel{~}{\lambda }}\left\{_\mu \vartheta ^\mu \vartheta +\mathrm{sin}^2\vartheta _\mu \phi ^\mu \phi +_\mu \chi ^\mu \chi +2\sqrt{1+\stackrel{~}{g}}\mathrm{cos}\vartheta ϵ^{\mu \nu }_\mu \chi _\nu \phi \right\}.$$ (8) We now perform an Abelian T-duality transformation with respect to the variable $`\chi `$, which corresponds to the canonical transformation $$\chi ^{}=\frac{\stackrel{~}{\lambda }}{\sqrt{1+\stackrel{~}{g}}}p_\alpha p_\chi =\frac{\sqrt{1+\stackrel{~}{g}}}{\stackrel{~}{\lambda }}\alpha ^{},$$ (9) where (and in the following) $`p_\chi `$ resp. $`p_\alpha `$ denote the canonical momenta conjugate to $`\chi `$ resp. to its ‘dual’ $`\alpha `$. In terms of these new variables the dual Lagrangian turns out to be: $$_\mathrm{\Sigma }=\frac{1}{2\stackrel{~}{\lambda }}\left\{_\mu \vartheta ^\mu \vartheta +(1+\stackrel{~}{g}\mathrm{cos}^2\vartheta )_\mu \phi ^\mu \phi +(1+\stackrel{~}{g})\left[_\mu \alpha ^\mu \alpha +2\mathrm{cos}\vartheta _\mu \alpha ^\mu \phi \right]\right\},$$ (10) which is nothing but the Lagrangian (2), when parametrizing the SU(2) valued field, $`G`$, by the Euler angles $$G=e^{i\phi \tau ^3}e^{i\vartheta \tau ^1}e^{i\alpha \tau ^3},$$ (11) and taking into account the relations at the classical level between the couplings: $$\stackrel{~}{\lambda }=\lambda ,\stackrel{~}{g}=g.$$ (12) The observation, that the axionic model is the T dual of $`_\mathrm{\Sigma }`$ also explains why $`\theta `$ is an angular variable. Indeed it has been shown in Ref. that in case of the principal $`\sigma `$-model ($`g=0`$) the Abelian T duality (9) maps its target space ($`S^3`$) into $`S^2\times S^1`$. The arguments of Ref. can be easily applied to the present case with $`g>1`$, and it is clear that in Eq. (1) $`n^a`$ parametrize the $`S^2`$ and $`\theta `$ parametrizes the $`S^1`$. In fact the equations of motion of both models (1) and (2) are known to admit a Lax representation indicating their (classical) integrability . Indeed, introducing the matrix valued current $$I_\mu =\frac{\stackrel{~}{\lambda }}{8\pi }nϵ_{\mu \nu }^\nu \theta \frac{\sqrt{\stackrel{~}{g}}}{2}ϵ_{\mu \nu }^\nu n+\frac{1}{2}n_\mu n,$$ (13) where $`n=in^a\sigma ^a`$, the equations of motion of (1) can be written as: $$^\mu I_\mu =0,_\mu I_\nu _\nu I_\mu =[I_\mu ,I_\nu ].$$ (14) The standard form (14) of the equations of motion allows for the introduction of a Lax pair $$U_\pm =\frac{1}{1\pm \omega }I_\pm ,$$ (15) satisfying the zero curvature equation $$_\mu U_\nu _\nu U_\mu =[U_\mu ,U_\nu ],$$ (16) for all values of the spectral parameter $`\omega `$. The current, $`I_\mu `$, is closely related to the matrix valued Noether current, $`𝒩_\mu =i\tau ^a𝒩_\mu ^a`$, defined by $`\delta =^\mu \epsilon ^a𝒩_\mu ^a`$ corresponding to the symmetry transformation $`\delta n^a=ϵ^{abc}\epsilon ^bn^c`$: $$I_\mu =\stackrel{~}{\lambda }𝒩_\mu +ϵ_{\mu \nu }^\nu T,T=\left(\frac{\stackrel{~}{\lambda }}{8\pi }\theta \frac{\sqrt{\stackrel{~}{g}}}{2}\right)n.$$ (17) The fact that, apart from a trivially conserved piece, $`I_\mu `$ can be identified with the Noether current of the manifest O(3) symmetry of the Lagrangian explains only the first equation in Eq. (14). The trivially conserved part of $`I_\mu `$ is essential that the zero curvature equation be also satisfied. The equations of motion of the deformed principal model (2) can be written entirely in terms of the current $`L_\mu `$ as $$^\mu L_\mu ^3=0,^\mu L_\mu ^1=igL^{2\mu }L_\mu ^3,^\mu L_\mu ^2=igL^{1\mu }L_\mu ^3.$$ (18) It is known that this system can be put to the Lax form , i.e. there is a spectral parameter dependent current, $`V_\mu =\tau ^aV_\mu ^a`$, satisfying the zero curvature equation (16). This current can be written as: $$V_\pm ^{1,2}=\alpha _\pm L_\pm ^{1,2},V_\pm ^3=a_\pm L_\pm ^3,$$ (19) where $$\alpha _\pm =\frac{4+g\omega ^2}{4g\omega ^2\pm 4\omega },a_\pm =\frac{4g\omega ^24g\omega }{4g\omega ^2\pm 4\omega }.$$ (20) We can now use the classical T-duality transformation (9) to map the linear system of the axion model (14) to a new Lax pair for the deformed $`\sigma `$-model (2). It is given by Eq. (15), where the current, $`I_\mu `$, has to be replaced by $$\widehat{I}_\mu =_\mu GG^1+g\left(G\tau ^3G^1\right)L_\mu ^3i\sqrt{g}ϵ_{\mu \nu }^\nu \left(G\tau ^3G^1\right).$$ (21) Eq. (21) is obtained from (13) by the T-duality transformation (9). $`\widehat{I}_\mu `$ is related to the Noether current $`\widehat{𝒩}_\mu `$, corresponding to the manifest symmetry $`\delta G=i\epsilon ^a\tau ^aG`$ of (2) and can be written analogously to $`𝒩_\mu `$: $$\widehat{I}_\mu =\lambda \widehat{𝒩}_\mu +ϵ_{\mu \nu }^\nu \widehat{T},\widehat{T}=i\sqrt{g}G\tau ^3G^1.$$ (22) It is clear that the new Lax pair (15) and the ‘old’ one, (19), cannot be related by a gauge transformation since they have different pole structures as functions of the spectral variable, $`\omega `$. In the $`g1`$ limit, the axion model reduces to the original O(3) $`\sigma `$-model (decoupled from the $`\theta `$ field), and the Lax pair (15) becomes equivalent to that of Ref. , where it has been pointed out that the corresponding $`\widehat{I}_\mu `$’s are ultralocal currents. We note that the Lax pairs (15) and (19) correspond to (different) deformations of the usual Lax pairs of the principal chiral $`\sigma `$-model, linear in $`_\mu GG^1`$ respectively $`G^1_\mu G`$. Next we carry out a standard test on the proposed $`S`$-matrix (4) for the axion model by comparing its (zero temperature) free energy obtained from the Thermodynamical Bethe Ansatz (TBA) and in weak coupling perturbation theory (PT) . For the deformed $`\sigma `$-model (2) this comparison has been done in Ref. where complete consistency has been found between the results of PT and of the TBA. For the axion model it is sufficient to compute the free energy in PT as the results of the TBA can be literally taken over from Ref. . In the present case one obtains as a bonus, a further nontrivial check on the quantum equivalence between the axion and the deformed $`\sigma `$-model, hence also on the validity of the T-duality transformation at the quantum level. Up to now when quantum equivalence between dually related models has been tested, mostly $`\beta `$-functions have been compared. The fact that the higher coefficients of the $`\beta `$-functions are scheme dependent makes such a comparison more difficult and less conclusive. The equivalence of the $`\beta `$-functions is certainly a necessary condition for the validity of quantum T-duality. At one loop order the $`\beta `$-functions of the couplings, $`\beta _\lambda `$, $`\beta _g`$ and $`\beta _{\stackrel{~}{\lambda }}`$, $`\beta _{\stackrel{~}{g}}`$ are simply obtained from each other by the classical relation (12). At two loops, however, it has been found in that using the background field method and dimensional regularization the following perturbative redefinition of the couplings $$\stackrel{~}{\lambda }=\lambda +\frac{\lambda ^2}{4\pi }(1+g),\stackrel{~}{g}=g+\frac{\lambda }{4\pi }(1+g)^2,$$ (23) (i.e. a change of scheme) is induced by the T-duality transformation. Taking into account Eqs. (23) the two loop $`\beta `$-functions of the two models turn out to be equivalent. Alternatively, introducing a renormalization group (RG) invariant combination of the two couplings: $$p=2\pi \underset{t\mathrm{}}{lim}\frac{1+g(t)}{\lambda (t)},\stackrel{~}{p}=2\pi \underset{t\mathrm{}}{lim}\frac{1+\stackrel{~}{g}(t)}{\stackrel{~}{\lambda }(t)},$$ (24) where $`t\mathrm{ln}h`$, one finds $`p=\stackrel{~}{p}`$ up to two loops . It is important to note that the RG invariant quantity (24) can be consistently identified with the parameter $`p`$ in the $`S`$ matrix (4). Let us introduce an effective $`\beta `$-function for $`\lambda (t)`$ by $`\beta _{\mathrm{eff}}(\lambda ,p)=\beta _\lambda (\lambda ,\mathrm{\Gamma }(\lambda ,p))`$, expressing $`g(t)`$, in terms of the running coupling, $`\lambda (t)`$, and $`p`$ as $`g(t)=\mathrm{\Gamma }(\lambda (t),p)`$. Using the perturbative result for $`\mathrm{\Gamma }(\lambda (t),p)`$ one finds $$\beta _{\mathrm{eff}}(\lambda ,p)=\beta _{\mathrm{eff}}(\stackrel{~}{\lambda },\stackrel{~}{p})=\frac{\lambda ^2}{2\pi }+\frac{p2}{8\pi ^2}\lambda ^3+\mathrm{}.$$ (25) Thus as far as coupling constant renormalization is concerned, the two models are equivalent, both are asymptotically free, and the actual value of $`p`$ effects only the two loop coefficient. The classical free energy density is obtained by minimizing the Legendre transform of the Hamiltonian density coupled to some conserved currents $$\widehat{}=_0h_iJ_0^i,\widehat{H}=𝑑x\widehat{}=Hh_iQ_i.$$ (26) Since the axion field, $`\theta `$, is actually an angle, its winding number (topological charge) can be non trivial. Therefore we present here the Legendre transformation of the modified Hamiltonian (26) for a rather general case. Let us consider a general sigma model with torsion $$_0=\frac{1}{2}g_{AB}^\mu X^A_\mu X^B+\frac{1}{2}b_{AB}ϵ^{\mu \nu }_\mu X^A_\nu X^B,$$ (27) and the following Ansatz for a set of conserved currents $$J_\mu ^i=𝒞_A^i(X)_\mu X^A+ϵ_\mu ^\nu _A^i(X)_\nu X^A,$$ (28) sufficiently general to include topological currents. The Legendre transformation of (26) yields the Lagrangian of the modified model which can be written as $$\widehat{}=_0+h^iJ_0^i+\frac{1}{2}h^ih^j𝒞_A^i𝒞^{Aj}.$$ (29) In fact $`\widehat{}`$ can be obtained by gauging $`_0`$ i.e. by the substitution $$_\mu X^A_\mu X^A+h^i\delta _{\mu 0}𝒞^{iA}$$ (30) when the antisymmetric field $`b_{AB}`$ is invariant (without compensating gauge transformation) under the symmetry transformation generated by the conserved currents (28). Below we also give a class of classical ground states (around which $`\widehat{}`$ is to be expanded) assuming that the metric, the antisymmetric tensor field and the quantities characterizing the currents are independent of a set of coordinates, $`\theta ^\alpha `$, corresponding to the splitting $`X^A=(y^k,\theta ^\alpha )`$: $$g_{AB}=g_{AB}(y),b_{AB}=b_{AB}(y),𝒞_A^i=𝒞_A^i(y),_A^i=_A^i(y).$$ (31) In this case the ground state is characterized by constant $`y^k`$-s and constant $`\theta _{}^{}{}_{}{}^{\alpha }`$-s $`y^ky_0^k`$, $`\theta _{}^{}{}_{}{}^{\alpha }\theta _{}^{}{}_{0}{}^{\alpha }`$, where the $`y_0^k`$-s stand for the extrema of $$H_{\mathrm{eff}}=\frac{1}{2}h^ih^j(g_{AB}𝒞^{Ai}𝒞^{Bj}+_\alpha ^i_\beta ^j(\gamma ^1)^{\alpha \beta }),$$ (32) and $$\theta _{}^{}{}_{0}{}^{\alpha }=(\gamma ^1)^{\alpha \beta }_\beta ^i(y_0)h^i.$$ (33) In (32-33) $`\gamma _{\alpha \beta }`$ denotes the restriction of $`g_{AB}`$ to the submanifold coordinatized by $`\theta ^\alpha `$. The axion model has a ‘manifest’ (i.e. up to a total derivative) SU(2)$`\times \mathrm{U}_\theta (1)`$ symmetry, where the $`\mathrm{U}_\theta (1)`$ subgroup is generated by the shift $`\theta \theta +\mathrm{const}`$. It is very important to note that although the SU(2) symmetry of the Lagrangian (1) corresponds to that of the $`S`$-matrix, the ‘manifest’ $`\mathrm{U}_\theta (1)`$ symmetry cannot be identified with the corresponding one of the $`S`$-matrix (4). Here the duality transformation provides the clue; the corresponding $`\stackrel{~}{\mathrm{U}}_\mathrm{R}`$(1) symmetry of the axion model is actually the image of the manifest U<sub>R</sub>(1) symmetry of (2) under the duality transformation, thus it is generated by the topological current of the axion field. Corresponding to the SU$`{}_{\mathrm{L}}{}^{}(2)`$U<sub>R</sub>(1) symmetry of the deformed $`\sigma `$-model there are two Noether charges, $`Q_\mathrm{L}`$ resp. $`Q_\mathrm{R}`$, associated to the U<sub>L</sub>(1) resp. U<sub>R</sub>(1) subgroups. Introducing two chemical potentials coupled to the $`Q_\mathrm{L}`$ resp. $`Q_\mathrm{R}`$, charges the Hamiltonian (26) takes the form: $`H=H_\mathrm{\Sigma }h_\mathrm{L}Q_\mathrm{L}h_\mathrm{R}Q_\mathrm{R}.`$ Then one can distinguish between three different types of finite density ground states: LEFT with $`h_\mathrm{L}>0,`$ and $`h_\mathrm{R}=0,`$ RIGHT with $`h_\mathrm{L}=0,`$ and $`h_\mathrm{R}>0,`$ and DIAG where $`h_\mathrm{L},h_\mathrm{R}>0`$. As found in the RIGHT case is obtained from DIAG by letting $`h_\mathrm{L}=0`$ in the final results. We compute below the corresponding ground state energies to one loop order in the axion model (1), starting with the LEFT case first. With the Euler angle parametrization of $`G`$ (11) the U$`{}_{\mathrm{L}}{}^{}(1)`$ transformation, $`Ge^{i\kappa \tau ^3}G`$ of the deformed $`\sigma `$-model (2) acts as a simple shift, $`\phi (x)\phi (x)+\kappa `$. The corresponding Noether charge, $`Q_\mathrm{L}`$, and its image under the T-duality transformation, $`\stackrel{~}{Q}_\mathrm{L}`$, are simply $$Q_\mathrm{L}=𝑑xp_\phi ,\stackrel{~}{Q}_\mathrm{L}=𝑑x\stackrel{~}{p}_\phi ,\mathrm{where}p_\phi =\frac{_\mathrm{\Sigma }}{\dot{\phi }},\stackrel{~}{p}_\phi =\frac{}{\dot{\phi }},$$ since the canonical transformation implementing the T-duality mapping (9) effects only $`p_\alpha `$, $`\chi ^{}`$, $`\alpha ^{}`$ and $`p_\chi `$, leaving the other fields, $`\phi `$, $`\vartheta `$, $`p_\phi `$, $`p_\vartheta `$, unchanged. Since in the LEFT case the $`b_{AB}`$ field in Eq. (1) is invariant, one can simply ‘gauge’ the Lagrangian of the axion model in an external ($`h_\mathrm{L}`$) field (see Eq. (30). The classical ground state is found to be $`\phi \chi 0`$, $`\vartheta \pi /2`$. (The corresponding solution of the deformed $`\sigma `$-model is given by $`\phi \alpha 0`$, $`\vartheta \pi /2`$.) Expanding the (Euclidean) Lagrangian (after suitable rescalings, etc.) we obtain $$\overline{}=\frac{2h_\mathrm{L}^2}{\stackrel{~}{\lambda }_0}+\frac{1}{2}mm^T+𝚘(\stackrel{~}{\lambda }),$$ (34) where $$=\left(\begin{array}{ccc}^2+4h_\mathrm{L}^2& 0& 2h_\mathrm{L}\sqrt{1+\stackrel{~}{g}_0}ϵ_{\mu 2}_\mu \\ 0& ^2& 0\\ 2h_\mathrm{L}\sqrt{1+\stackrel{~}{g}_0}ϵ_{\mu 2}_\mu & 0& ^2\end{array}\right),$$ (35) and $`m=(\vartheta ,\phi ,\chi )`$. ($`\stackrel{~}{\lambda }_0`$, $`\stackrel{~}{g}_0`$ denote the bare coupling and parameter of the axion/dual model). In Eq. (35) we kept the $`ϵ`$ tensor explicitly, as it requires a careful definition in $`n=2ϵ`$ dimensions which we use to regularize the momentum integrals. We adopt the definiton of , where this antisymmetric tensor corresponds to an almost complex structure: $`ϵ_{\mu \nu }=ϵ_{\nu \mu }`$, $`ϵ_{\mu \nu }ϵ_{\mu \sigma }=\delta _{\nu \sigma }`$. The one loop quantum corrections to the classical ground state (the first term in Eq. (34)) require the calculation of a functional determinant, leading to $$(h)=\frac{4h_\mathrm{L}^2}{n}\frac{d^np}{(2\pi )^n}\frac{\stackrel{~}{p}_1^2\stackrel{~}{g}_0\stackrel{~}{p}_2^2}{\stackrel{~}{p}^4+4h_\mathrm{L}^2(\stackrel{~}{p}_1^2\stackrel{~}{g}_0\stackrel{~}{p}_2^2)},\stackrel{~}{p}_\mu =ϵ_{\mu \nu }p_\nu .$$ (36) To evaluate (36) we apply the modified dimensional regularization of Ref. as $`\stackrel{~}{p}_2=ϵ_{2\nu }p_\nu `$ plays here a distinguished role and it is kept as a one dimensional variable. In fact for our purposes it is sufficient to calculate the difference $`(h)_\mathrm{\Sigma }(h)`$, where $`_\mathrm{\Sigma }(h)`$ is the corresponding determinant in the deformed $`\sigma `$-model (Eq. (3.12) in Ref. ). Since both $`(h)`$ and $`_\mathrm{\Sigma }(h)`$ are already the first quantum corrections to the classical expressions we may set $`\stackrel{~}{g}=g`$ (and make no distinction between bare and renormalized $`g`$’s) when computing their difference to lowest order and we end up with $$(h)_\mathrm{\Sigma }(h)=\frac{(2h_\mathrm{L})^n}{n}(1+g)\frac{d^nq}{(2\pi )^n}\frac{(q_1^2q_2^2)q^4}{N_1N_2}=\frac{(2h_\mathrm{L})^n}{n}(1+g)w(g),$$ (37) where $`N_1=q^4+q_1^2gq_2^2`$, $`N_2=q^4+q_2^2gq_1^2`$. Although the integrand yielding $`w(g)`$ is antisymmetric under $`q_1q_2`$, the integral is divergent by power counting for $`n=2`$, i.e. it must be computed in $`n=2ϵ`$ dimensions. Its derivative, $`w^{}(g)`$, is, however, convergent by power counting and it has also an antisymmetric integrand, therefore this latter may be evaluated in $`n=2`$ dimensions giving $`w^{}(g)0`$. Then to compute $`w(g)`$ one may choose e.g. the point $`g=1`$: $$w(1)=\frac{d^nq}{(2\pi )^n}\frac{q_1^2q_2^2}{(q^2+1)^2}=\frac{n11}{n}\frac{d^nq}{(2\pi )^n}\frac{q^2}{(q^2+1)^2}=\frac{1}{4\pi },$$ (38) where writing the second equality, we used that $`q_1`$ is $`n1`$ dimensional, while $`q_2`$ is a $`1`$ dimensional variable. From (38) one finds that after taking into account the change of the renormalization scheme (23), in PT the free energy densities of the two models (1) and (2) do indeed coincide for the LEFT case. Recently this calculation has been performed also at the two-loop level . To discuss the RIGHT and DIAG cases we find it more convenient to use the parametrization of for the SU(2) valued field, $`G`$: $$G=\frac{i\sigma ^2}{\sqrt{1+\|\mathrm{\Psi }\|^2}}\left(\begin{array}{cc}1& \mathrm{\Psi }^{}\\ \mathrm{\Psi }& 1\end{array}\right)\left(\begin{array}{cc}e^{i\mathrm{\Phi }}& 0\\ 0& e^{i\mathrm{\Phi }}\end{array}\right),$$ (39) where $`\mathrm{\Psi }`$ resp. $`\mathrm{\Phi }`$ is a complex resp. a real scalar field. Now $`U_\mathrm{R}(1)`$ acts as a shift, $`\mathrm{\Phi }\mathrm{\Phi }+\kappa `$, and then the corresponding Noether charge of the deformed $`\sigma `$-model is simply $`Q_\mathrm{R}=𝑑xp_\mathrm{\Phi }`$. $`Q_\mathrm{L}`$ is slightly more complicated when expressed in terms of the canonical momenta (as $`e^{i\kappa \sigma ^3}i\sigma ^2=i\sigma ^2e^{i\kappa \sigma ^3}`$): $`Q_\mathrm{L}=𝑑x[p_\mathrm{\Phi }+2i(p_\mathrm{\Psi }\mathrm{\Psi }p_\mathrm{\Psi }^{}\mathrm{\Psi }^{})].`$ Using Buscher’s rule , the Lagrangian of the axion (dual) model now takes the form: $$^d=\frac{\stackrel{~}{\lambda }}{8(1+\stackrel{~}{g})}(_\mu f)^2+\frac{2}{\stackrel{~}{\lambda }}\frac{_\mu \mathrm{\Psi }^\mu \mathrm{\Psi }^{}}{N^2}+\frac{1+\stackrel{~}{g}}{\stackrel{~}{\lambda }}\widehat{𝒜}_\mu \widehat{𝒜}^\mu \frac{i}{2}ϵ^{01}(\dot{f}\widehat{𝒜}_1f^{}\widehat{𝒜}_0),$$ (40) where $`f`$ is the dual to $`\mathrm{\Phi }`$, and $$\widehat{𝒜}_\mu =𝒜_\mathrm{\Psi }^{}_\mu \mathrm{\Psi }^{}𝒜_\mathrm{\Psi }_\mu \mathrm{\Psi }=\frac{1}{N}(\mathrm{\Psi }_\mu \mathrm{\Psi }^{}\mathrm{\Psi }^{}_\mu \mathrm{\Psi }),N=1+\|\mathrm{\Psi }\|^2.$$ (41) The canonical transformation connecting $`_\mathrm{\Sigma }`$ and $`^d`$, maps $`Q_\mathrm{R}`$ and $`Q_\mathrm{L}`$ to $$\stackrel{~}{Q}_\mathrm{R}=𝑑xf^{},\stackrel{~}{Q}_\mathrm{L}=𝑑x\left[f^{}+2i(\stackrel{~}{p}_\mathrm{\Psi }\mathrm{\Psi }\stackrel{~}{p}_\mathrm{\Psi }^{}\mathrm{\Psi }^{})\right],$$ (42) i.e. $`\stackrel{~}{Q}_\mathrm{R}`$ and $`\stackrel{~}{Q}_\mathrm{L}`$ do indeed contain the topological charge of the axion field (proportional to $`f`$). Applying now the general framework, Eqs. (32-33) to the present cases; $`i=`$(L,R), $`X^A=(`$$`f`$,$`\mathrm{\Psi }`$,$`\mathrm{\Psi }^{}`$), with $`\theta ^\alpha =(f)`$, $`y^k=(\mathrm{\Psi },\mathrm{\Psi }^{})`$. Using the explicit form of $`\stackrel{~}{p}_\mathrm{\Psi }`$ and $`\stackrel{~}{p}_\mathrm{\Psi }^{}`$ one finds for $`\stackrel{~}{J}_\mu ^{\mathrm{R},\mathrm{L}}`$: $$𝒞_A^\mathrm{R}0,_A^\mathrm{R}=\{\begin{array}{cc}1,A=f\hfill & \\ 0,A=\mathrm{\Psi },\mathrm{\Psi }^{},\hfill & \end{array}$$ (43) $$𝒞_A^\mathrm{L}=\{\begin{array}{ccc}0,A=f\hfill & & \\ 𝒩𝒜_\mathrm{\Psi }/\stackrel{~}{\lambda },A=\mathrm{\Psi }\hfill & & \\ 𝒩𝒜_\mathrm{\Psi }^{}/\stackrel{~}{\lambda },A=\mathrm{\Psi }^{}\hfill & & \end{array}_A^\mathrm{L}=\{\begin{array}{cc}(12\|\mathrm{\Psi }\|^2/N),A=f\hfill & \\ 0,A=\mathrm{\Psi },\mathrm{\Psi }^{},\hfill & \end{array}$$ (44) where $`𝒩=4i(12(1+\stackrel{~}{g})\|\mathrm{\Psi }\|^2)/N`$. Substituting these $`𝒞_A^i`$ and $`_A^i`$ into Eq. (32) reveals that $`H_{\mathrm{eff}}`$ depends only on $`\|\mathrm{\Psi }\|^2`$ and that its extremum is at $`\mathrm{\Psi }=0=\mathrm{\Psi }^{}`$. In the DIAG case the actual value of the ground state energy density at this extremum is given by: $$\widehat{H}\|_{\mathrm{min}}=\frac{2(1+\stackrel{~}{g})}{\stackrel{~}{\lambda }}(h_\mathrm{R}+h_\mathrm{L})^2,$$ (45) and the expectation value of $`f^{}`$ is: $`f_0^{}=4(1+\stackrel{~}{g})(h_\mathrm{R}+h_\mathrm{L})/\stackrel{~}{\lambda }`$. We note that $`\widehat{H}\|_{\mathrm{min}}`$ agrees (as it should) with the corresponding result of the deformed $`\sigma `$-model with $`\stackrel{~}{\lambda }\lambda `$, $`\stackrel{~}{g}g`$ (Eq. (3.20) in ). For the RIGHT case the analogous expressions of the axion model are simply obtained from (45) by setting $`h_\mathrm{L}=0`$. At this point we recall the somewhat unusual feature of the axion model once more, i.e. that the U(1) symmetry of the $`S`$-matrix (4) is realized through a topological current analogously to the Sine-Gordon theory. To emphasize this we quote here the value of the classical free energy density corresponding to the Noether charge of the ‘manifest’ U<sub>θ</sub>(1) symmetry of the Lagrangian (1): $`\widehat{H}\|_{\mathrm{min}}^{(\theta )}=2h_\mathrm{R}^2/\stackrel{~}{\lambda }`$, quite different from Eq. (45) with $`h_\mathrm{L}=0`$. To obtain the one loop correction to the free energy, one has to expand $`\widehat{}^d`$ around the minimum, Eq. (45). Writing $`f=xf_0^{}+\widehat{f}`$, where the expectation value of $`\widehat{f}`$ vanishes, one finds that the quadratic terms containing $`\widehat{f}`$ are independent of $`h_\mathrm{L}`$, $`h_\mathrm{R}`$, so the quadratic pieces of $`\widehat{}^d`$ (hence the one loop correction) are effectively determined by $`\mathrm{\Psi }=\sqrt{\stackrel{~}{\lambda }/2}\psi `$ only: $$_2=_\mu \psi ^\mu \psi ^{}m^2\|\psi \|^2+iq(\psi \dot{\psi }^{}\psi ^{}\dot{\psi }),$$ (46) where $`m^2=4h_\mathrm{L}(h_\mathrm{L}\stackrel{~}{g}_0+h_\mathrm{R}(1+\stackrel{~}{g}_0))`$ and $`q=h_\mathrm{L}(1\stackrel{~}{g}_0)h_\mathrm{R}(1+\stackrel{~}{g}_0)`$. After continuation to Euclidean space and writing $`\psi =(\varphi _1+i\varphi _2)/\sqrt{2}`$, Eq. (46) becomes identical to the corresponding pieces for the deformed $`\sigma `$-model, Eq. (C.11) in (with $`\stackrel{~}{\lambda }_0\lambda _0`$, $`\stackrel{~}{g}_0g_0`$). From this it follows that the free energy densities fully agree also for the DIAG and RIGHT cases in both models. Now the results of the comparison between the free energies computed by the TBA (based on the proposed $`S`$ matrix (4)) and in PT for the deformed $`\sigma `$-model (2) in Ref. can be simply taken over for the axion model. The conclusion is that there is complete consistency between the TBA and the perturbative calculations, providing good evidence for the validity of the proposed $`S`$ matrix (4) for the axion model. Since the effective coupling (25) is identical in the two models the $`m/\mathrm{\Lambda }_{\overline{\mathrm{MS}}}`$ ratio found in stays unchanged. Finally we would like to point out that it would be interesting to study the axion model by lattice Monte-Carlo simulations. This would provide us with a completely non-perturbative way of testing quantum T-duality. Technical difficulties arising from the non-reality of the Euclidean action in the context of the lattice Monte-Carlo study of the $`O(3)`$ model with a constant $`\theta `$ term and a suggestion how to circumvent them is discussed in . Acknowledgements J. B. gratefully acknowledges a CNRS grant. This investigation was supported in part by the Hungarian National Science Fund (OTKA) under T 030099 and T 029802 and by the Hungarian Ministry of Education under FKFP 0178/1999.
warning/0004/gr-qc0004035.html	ar5iv	text	# 1 Introduction ## 1 Introduction In the last decade wormhole configurations have been extensively studied and became more and more “respectable”, as they were proved to be almost ubiquitous. There are two main kind of wormholes. In 1988 Coleman , , proposed that the existence of microscopic Euclidean wormholes could explain why the cosmological constant is so small in our universe, and eventually yield a probability distribution for all fundamental constants of nature. This launched a flurry of interest on these kind of wormholes,(e.g. , ), although further results,(e.g. , ), have cast doubts on the initial claims. On the other hand, macroscopic wormholes connecting large Lorentzian universes have also been a subject of intense study, e.g. , . In , Hochberg considers a wormhole model that results from taking two copies of (Lorentzian) Friedman-Robertson-Walker universes, $$ds^2=dt^2+R^2(t)\left\{\frac{dr^2}{1kr^2}+r^2(d\theta ^2+\mathrm{sin}\theta ^2d\varphi ^2)\right\}$$ (1) and remove from each of them a $`4`$dimensional region of the form $`\mathrm{\Omega }_{1,2}=\{r_{1,2}a\}`$. The resulting $`4`$spaces will have identical boundaries $`\mathrm{\Omega }_{1,2}=\{r_{1,2}=a\}`$. By identifying these two boundaries one obtains two FRW spacetimes connected by a wormhole, whose throat is located on their common $`\mathrm{\Omega }`$ boundaries. There is no need to invoke any arbitrary $`3+1`$ decomposition of spacetime, usually ADM, as in the case of continuum models. This means that global issues like topology can still be addressed in the simplicial minisuperspace, which is not possible in the continuum versions. The simplicial framework deals with the universe in a fully $`4`$D way, not relying on any specific $`3+1`$ decomposition of spacetime, usually ADM in the case of continuum models. This means that global issues like topology can still be addressed in the simplicial minisuperspace, which is not possible in the continuum versions. Additionally, it will allow us to use the same simplicial minisuperspace model to study both kinds of wormholes. Simplicial minisuperspace models have proved very reliable in the past by confirming results obtained with the continuum formalism , , . However, they are at their best when applied to situations which the continuum theory cannot easily handle. Namely, in the study of the effects of generalising the definition of history in quantum gravity, to include conifolds and not just manifolds. See ,, . Furthermore, they can sometimes deliver some unexpected results, like the existence of classical Euclidean solutions for all sizes of the boundary $`3`$universe, when we consider an arbitrary scalar coupling $`\eta R\varphi ^2`$, . ## 2 Simplicial Quantum Gravity The simplicial framework in QG is usually seen as an approximation to the more fundamental continuum framework. However, it has been argued that at the scales where quantum gravity effects become relevant, it is reasonable to assume that not just energy but even space and time will be discrete. In this way the simplicial framework could even be the more fundamental one. The basic building block of any simplicial complex is the $`n`$simplex. ###### Definition 2.1 Given $`n+1`$ points in $`R^{n+1}`$, labelled $`v_1,v_2,\mathrm{},v_{n+1}`$, if these points are affinely independent, then an $`n`$dimensional simplex $`\sigma ^n`$ is the convex hull of these points $$\sigma ^n=\{xR^{n+1}:x=\underset{i=1}{\overset{n+1}{}}\lambda _iv_iwhere\lambda _i0and\underset{i=1}{\overset{n+1}{}}\lambda _i=1\}$$ (2) Thus, a $`0`$simplex is a vertex, a $`1`$simplex is an edge, a $`2`$simplex a triangle, etc. Any subset of these vertices also spans a simplex, of lower dimension, which is called a face of $`\sigma ^n`$. From these building blocks we can construct a whole variety of spaces in particular some which by obeying certain natural conditions offer the widest reasonable framework in which to investigate the concept of simplicial geometry. These are called simplicial complexes. ###### Definition 2.2 A simplicial complex $`(K,K)`$ is a topological space $`K`$ and a collection of simplices $`K`$, such that * $`K`$ is a closed subset of some finite dimensional Euclidean space. * If $`\sigma `$ is a face of a simplex in $`K`$, then $`\sigma `$ is also contained in $`K`$. * If $`\sigma _a`$ and $`\sigma _b`$ are simplices in $`K`$, then $`\sigma _a\sigma _b`$ is a face of both $`\sigma _a`$ and $`\sigma _b`$. * The topological space $`K`$ is the union of all simplices in $`K`$. By definition the dimension of the simplicial complex is the maximum dimension of any of its simplices. Each simplex in $`K`$ is totally described by its vertices. Since a simplicial complex describes both the simplices of the space $`K`$ and gives the rules for how these building blocks are connected, then the simplicial complex itself is uniquely determined by the vertices and the rules stating in which simplices they are contained in. There are several standard topological definitions that extend naturally to simplicial complexes, like ###### Definition 2.3 A simplicial complex is compact it it contains a finite number of simplices. ###### Definition 2.4 A simplicial complex is connected if any two of its vertices are connected by a sequence of edges. We shall see that in order for it to be possible to define curvature in a complex, (and consequently an action functional for QG), it is essential that this complex be pure, i.e., ###### Definition 2.5 A simplicial $`n`$complex, $`(K^n,K^n)`$, is pure when every lower dimensional simplex in it is contained in at least one $`n`$simplex of $`K^n`$. A natural definition of boundary is also essential. However, it can be shown that the notion of boundary of a complex is meaningful for only a special subset of pure simplicial complexes, which is defined by ###### Definition 2.6 A pure simplicial $`n`$complex $`(K^n,K^n)`$, is said to be non-branching if every $`(n1)`$simplex is contained in either one or two $`nsimplices`$. Only now can we define the boundary of a complex as: ###### Definition 2.7 The boundary, $`K^{n1}`$, of a pure non-branching simplicial $`n`$complex $`K^n`$ is the complex made up of all $`(n1)`$simplices that are faces of one and only one $`n`$simplex of $`K^n`$. An essential concept for the study of the relationship between continuum spaces and their simplicial representations is that of triangulation ###### Definition 2.8 A simplicial complex $`(K,K)`$ is said to triangulate a smooth space $`𝒦`$ if there is an homeomorphism $`\mathrm{\Phi }`$ between $`K`$ and $`𝒦`$. Conventionally, if $`𝒦`$ is an $`n`$sphere then it is said that $`(K,K)`$ is a simplicial $`n`$sphere, the same for an $`n`$ball, etc. In the topological sense we can say that $`K`$ is the simplicial counterpart of $`𝒦`$. What characterises a continuum topological manifold is the fact that all its points have neighbourhoods homeomorphic to open sets in $`R^n`$, like an open ball $`B^n`$. So the definition of the simplicial version of a continuum manifold is dependent on the definition of the simplicial equivalent of boundary of a point. Only then can the local topology of the complexes be studied. ###### Definition 2.9 Given a simplicial complex $`(K^n,K^n)`$, the combinatorial star of a vertex $`vK^n`$, denoted $`St(v)`$, is the complex consisting of all the $`n`$simplices of $`K^n`$ that contain $`v`$. ###### Definition 2.10 Given a simplicial complex $`(K^n,K^n)`$, the combinatorial link of a vertex $`vK^n`$, denoted $`L(v)`$, is the complex consisting of all the simplices that are in $`St(v)`$ but do not contain $`v`$ itself. The star of a vertex is very much like the neighbourhood of a point in a continuum framework. So we expect that the defining characteristic of the simplicial analogues of manifolds to be that the star of each vertex be homeomorphic to an open ball $`B^n`$, which is equivalent to say that its star is a simplicial open $`n`$ball. However it is customary to reformulate these conditions in terms of links and not stars. It is easy to see that to say that $`St(v)`$ is a simplicial open $`n`$ball is equivalent to say that $`L(v)`$ is a simplicial $`(n1)`$-sphere. And so we will restate the “manifold condition” as that the combinatorial link of every vertex be homeomorphic to an $`(n1)`$-sphere, i.e., $`L(v)`$ be a simplicial $`(n1)`$sphere. For technical reasons the previous conditions are still not enough for the resulting simplicial complexes to have the same properties as continuum manifolds. One last condition is needed. ###### Definition 2.11 A simplicial $`n`$complex $`(K^n,K^n)`$ is said to be strongly connected if any two simplices of $`K^n`$ can be connected by a sequence of $`n`$simplices each intersecting along some $`(n1)`$-simplex. We can now define the simplicial counterparts of smooth manifolds: ###### Definition 2.12 A combinatorial $`n`$manifold $`^n`$, is an $`n`$dimensional simplicial complex such that * It is pure. * It is non-branching. * Any two $`n`$simplices can be connected by a sequence of $`n`$simplices, each intersecting along some $`(n1)`$simplex. * The link of every vertex is a simplicial $`(n1)`$sphere. Note that there are simplicial complexes that are homeomorphic to topological manifolds but are not combinatorial manifolds. The definition of combinatorial manifold carries more structure than simply the topology. ### 2.1 Simplicial Geometry We have so far been concerned solely with the topological and combinatorial aspects of simplicial quantum gravity. We shall now focus on the geometric issues. The simplest way to describe a geometry on an any pure non-branching simplicial complex $`K^n`$ is to use the Regge formalism: * 1)We require the metric on the interior of each $`n`$simplex of $`K^n`$ to be flat. * 2) We then assign edge lengths to each edge of $`K^n`$. Not every assignment of edge lengths is consistent with the simplices having flat interiors. The triangle inequalities and their higher dimensional analogues must be satisfied. Necessary and sufficient conditions for this to happen are that the squared volumes of all $`k`$simplices, with $`k=2,3,\mathrm{}n`$, must be positive. The metric information, which in the continuum spaces is contained on the metric tensor $`g_{\mu \nu }`$, is imprinted on a simplicial complex, $`K^n`$, via an assignment of its squared edge lengths. $$g_{\mu \mu }(x)g_{ij}(\{s_k\})=\frac{s_{0i}+s_{0j}s_{ij}}{2}$$ (3) where $`0`$ is just an arbitrary vertex of $`K^n`$, and $`s_{ij}`$ is the square edge length of the edge $`[ij]`$. Note that in the simplicial framework the metric degrees of freedom are now the squared edge lengths $`s_k`$. This means that all geometric quantities should be expressible in terms of them. Since by definition the metric inside each $`n`$simplex is flat the curvature in a simplicial complex cannot reside there. Instead it is located at the $`(n2)`$simplices of the complex. The curvature associated with each $`(n2)`$simplex is measured by what is called its deficit angle, whose expression is slightly different according to whether the $`(n2)`$simplex is located in the interior or on the boundary of the complex: For an interior $`(n2)`$simplex, $`\sigma _i^{n2}`$, the deficit angle is $$\theta (\sigma _i^{n2})=2\pi \underset{\sigma ^nSt(\sigma _i^{n2})}{}\theta _d(\sigma _i^{n2},\sigma ^n)$$ (4) where $`\theta _d(\sigma _i^{n2},\sigma _n)`$ is called the dihedral angle of $`\sigma _i^{n2}`$ associated with $`\sigma _i^n`$, and is the angle between the the two $`(n1)`$-simplices that belong to $`\sigma _i^n`$ and intersect at $`\sigma _i^{n2})`$. For a $`(n2)`$simplex, $`\sigma _b^{n2}`$, on the boundary of the complex, the deficit angle is $$\psi (\sigma _b^{n2})=\pi \underset{\sigma ^nSt(\sigma _b^{n2})}{}\theta _d(\sigma _b^{n2},\sigma ^n)$$ (5) It is easy to see why the previous definitions of curvature are only valid for pure complexes. Only in a pure complex can we be sure that any $`(n2)`$simplex is indeed contained in some $`n`$simplex. Following Regge, the continuum Einstein action can be discretized for any pure non-branching complex as $`I[K^n,\{s_k\}]`$ $`=`$ $`{\displaystyle \frac{2}{16\pi G}}{\displaystyle \underset{\sigma _i^{n2}}{}}V_{n2}(\sigma _i^{n2})\theta (\sigma _i^{n2})+{\displaystyle \frac{2\mathrm{\Lambda }}{16\pi G}}{\displaystyle \underset{\sigma ^n}{}}V_n(\sigma ^n)`$ (6) $``$ $`{\displaystyle \frac{2}{16\pi G}}{\displaystyle \underset{\sigma _b^{n2}}{}}V_{n2}(\sigma _b^{n2})\psi (\sigma _b^{n2})`$ where $`V_p(\sigma ^p)`$ is the $`p`$ volume of the $`p`$simplex, $`\sigma _p`$. It is easy to see that all the volumes as all the deficit angles can be written as functions of the squared edge lengths $`\{s_k\}`$. See for their explicit form. We can now write the wavefunction of the universe in simplicial quantum cosmology as $$\mathrm{\Psi }[K^{n1},\{s_b\}]=\underset{K^n}{}_CD\{s_i\}e^{I[K^n,\{s_i\},\{s_b\}]}$$ (7) where * $`\{s_i\}`$ are the squared lengths of the interior edges * $`\{s_b\}`$ are the squared lengths of the boundary edges * The kind sum over complexes $`K^n`$ depends on the boundary conditions one adheres to. For example the equivalent of Hartle-Hawking‘s no boundary proposal would be to consider a sum over all compact combinatorial $`n`$manifolds $`^n`$ whose only boundary is $`^{n1}`$, and over all simplicial $`n`$geometries $`\{s_k\}`$ that have boundary edges with squared lengths $`\{s_b\}`$ Note that the while in the continuum we had a functional metric integral, we now have a well defined product of integrals over edge lengths. So it seems we have removed the problems related to issues of gauge fixing and renormalisation associated with defining the measure of the space of metrics. However, such problems reappear in any attempt to take the continuum limit of this Regge integral. Similarly the functional integral associated with the the scalar field now becomes a product of well defined simple integrals over the value of the field in each interior vertex. ## 3 Simplicial Minisuperspace A simplicial minisuperspace approximation (in a pure gravity model) consists of imposing two kinds of restrictions on the quantities being summed over in $`(\text{7}`$). Namely, it involves singling out a particular complex or family of complexes and singling out a few edge lengths, by making all others a function of them. We are thus imposing restrictions of a topological and geometrical nature. In our case we will have: $$\underset{K^n}{}W^4$$ $$\{s_i\}s_i$$ $$\{s_b\}a,b$$ Where $`W^4`$ and the other restrictions will be explained in more detail in the next section. ### 3.1 Topology Restrictions Consider two identical copies of the combinatorial manifold $`𝒞=0\alpha _4`$ we used in as a model of our universe. We shall denote them $`𝒞_1`$, and $`𝒞_2`$. As we have shown in the previous chapter these models provide a good approximation to the wavefunction of the universe, by predicting classical Lorentzian spacetime in the late universe. So we can consider them two separate but identical $`4`$D universes. We now proceed to construct a connection between them acting as a wormhole connecting these two universes. To do so note that $`𝒞_1`$ is composed of five $`4`$simplices: $$𝒞_1=[01234],[01345],[01235],[01245],[02345]$$ Remember that the $`4`$simplex $`[12345]`$ does not belong $`𝒞_1`$ because the triangulation $`\alpha _4`$ of the $`3`$sphere, is just the surface of $`[12345]`$. Similarly $`𝒞_2`$ is composed of five equivalent $`4`$simplices, $$𝒞_2=[0^{}1^{}2^{}3^{}4^{}],[0^{}1^{}3^{}4^{}5^{}],[0^{}1^{}2^{}3^{}5^{}],[0^{}1^{}2^{}4^{}5^{}],[0^{}2^{}3^{}4^{}5^{}]$$ These two copies of our model universe are displayed in figure $`5.1`$. If we now remove one of these identical $`4`$simplices from each complex, say, $`[02345]`$ and $`[0^{}2^{}3^{}4^{}5^{}]`$, and then glue the resulting complexes through the identification of the vertices $`0=0^{}`$, $`2=2^{}`$, $`3=3^{}`$, $`4=4^{}`$ and $`5=5^{}`$, then we have connected the two independent $`4`$D universes, through a $`3`$D throat, thus modelling a wormhole. The resulting $`4`$complex shall be denoted as $`W^4`$, and the identified vertices will be denoted as $`0,2,3,4,5`$, see figure $`5.2`$. There are of course two non-identified vertices $`1`$ and $`1^{}`$. Note also that given the identification $`0=0^{}`$, the complex $`W^4`$ has only one interior vertex. $`W^4`$ is composed of eight $`4`$simplices: $$W^4=[01234],[01235],[01345],[01245],[01^{}234],[01^{}235],[01^{}245],[01^{}345]$$ Note that the $`4`$simplex $`[02345]`$ does not exist in $`W^4`$. Furthermore, $`W^4`$ has only one boundary, the $`3`$complex made up of the following $`3`$simplices: $`W^4`$ $`=`$ $`[1234],[1^{}234],[1235],[1^{}235],[1245],[1^{}245],`$ (8) $`[1345],[1^{}345].`$ The $`3`$simplex $`[2345]`$ does not exist in $`W^4`$, since it belongs only to $`[12345]`$, $`[1^{}2345]`$ and $`[02345]`$ and by definition $`W^4`$ does not contain these $`4`$simplices, and since we want $`W^4`$ to be pure and non branching all $`3`$simplices in $`W^4`$ have to belong either to one or two $`4`$simplices. However, the four triangles that form the surface of $`[2345]`$, i.e., $`[234]`$, $`[235]`$, $`[245]`$ and $`[345]`$, do belong to $`W^4`$. We can thus conclude that the boundary $`W^4`$ is a combinatorial $`3`$manifold. Indeed if we compute the links of its vertices all are homeomorphic to $`S^2`$. $$L_W(1)=L(1^{})=[234][235][245][345]$$ (9) So we see that the links of the vertices $`1`$ and $`1^{}`$ are simply the surface of a tetrahedron $`[2345]`$. $$L_W(2)=[134][1^{}34][135][1^{}35][145][1^{}45]$$ (10) Thus, the link of vertex $`2`$ is just the surface of an hexahedron. The same is true for vertices $`3`$, $`4`$ and $`5`$. This is only valid because $`[2345]`$ does not belong to $`W^4`$. If that was not the case the link of vertex $`2`$ would not be homeomorphic to $`S^2`$ and the boundary $`W^4`$ would not be a $`3`$manifold. The throat $`T^3`$ of the wormhole is the $`3`$D surface of the $`4`$simplex $`[02345]`$ without $`[2345]`$ . See figure $`3`$. It is a combinatorial $`3`$manifold constituted by the $`3`$simplices: $$T^3=[0234],[0345],[0235],[0245]$$ ### 3.2 Geometry For simplicity and since there is still only one interior vertex, we shall assume that all interior edge lengths are equal, and their squared values are $`s_i`$. $$s_{01}=s_{01^{}}=s_{02}=s_{03}=s_{04}=s_{05}=s_i$$ However, it is essential that not all boundary edge lengths are the same if we want to separate the evolution of the whole boundary $`3`$universe from that of the wormhole’s throat. So we shall assume that all boundary edges belonging to the throat $`T^3`$ will have squared lengths $`b`$, $$s_{23}=s_{24}=s_{25}=s_{34}=s_{35}=s_{45}=b$$ while all the other boundary edges have square edge length $`a`$. ## 4 Minisuperspace Wavefunction Having defined our simplicial minisuperspace model our objective is now to evaluate the wavefunction associated with it. Although the $`4`$complex is now somewhat more complicated we simplify the model by not considering any matter sector. The general formulae obviously still apply. So the minisuperspace wavefunction will be given by $$\mathrm{\Psi }[a,b]=_CDs_ie^{I[s_i,a,b]}$$ (11) To implement this expression we start by computing the Euclideanized Regge action for the model. Since we have no matter sector we have simply $`I[W^4,a,b]`$ $`=`$ $`{\displaystyle \frac{2}{16\pi G}}{\displaystyle \underset{\sigma _2^i}{}}V_2(\sigma _2^i)\theta (\sigma _2^i)+{\displaystyle \frac{2\mathrm{\Lambda }}{16\pi G}}{\displaystyle \underset{\sigma _4}{}}V_4(\sigma _4)`$ (12) $``$ $`{\displaystyle \frac{2}{16\pi G}}{\displaystyle \underset{\sigma _2^b}{}}V_2(\sigma _2^b)\psi (\sigma _2^b)`$ where: * $`\sigma _k`$ denotes a $`k`$simplex belonging to the set $`\mathrm{\Sigma }_k`$ of all $`k`$simplices in $`W^4`$. * $`\theta (\sigma _2^i)`$, is the deficit angle associated with the interior $`2`$simplex $`\sigma _2^i=[ijk]`$ * $`\psi (\sigma _2^b)`$ is the deficit angle associated with the boundary $`2`$simplex $`\sigma _2^b`$ * $`V_k(\sigma _k)`$ for $`k=2,3,4`$ is the $`k`$volume associated with the $`k`$simplex, $`\sigma _k`$. As before all these quantities can be expressed in terms of the squared edge lengths, $`s_i`$, $`a`$ and $`b`$, according to the expressions in appendix $`1`$. * With our choice of edge lengths, all eight $`4`$simplices that compose $`W^4`$ are of the same type, and their $`4`$volume can be computed as $$V_4(\sigma _4)=\frac{b^2}{48}\sqrt{3\frac{as_i}{b^2}\frac{s_i}{b}\frac{3}{4}\frac{a^2}{b^2}}$$ (13) * There is only one type of boundary $`3`$simplex. The eight of them, can be denoted as $`[aaabbb]`$, i.e., they have three edges of squared length $`a`$ and three edges of squared length $`b`$. They all have the same volume. As an example take $`[1234]`$, its volume is $$V_3(\sigma _3^b)=\frac{\sqrt{3}b^{3/2}}{12}\sqrt{\frac{a}{b}\frac{1}{3}}$$ (14) * As for interior $`3`$simplices there are two types. Type I can be described as $`[aabs_is_is_i]`$, of which $`[0123]`$ is an example. There are twelve of them and they all have the same volume $$V_3(\sigma _3^{Ii})=\frac{b^{3/2}}{12}\sqrt{4\frac{as_i}{b^2}\frac{s_i}{b}\frac{a^2}{b^2}}$$ (15) There are four more interior $`3`$simplices, of type II, which are $`[0234]`$, $`[0235]`$, $`[0345]`$, and $`[0245]`$. They all have the same volume $$V_3(\sigma _3^{IIi})=\frac{\sqrt{3}b^{3/2}}{12}\sqrt{\frac{s_i}{b}\frac{1}{3}}$$ (16) * We have twelve boundary triangles of the kind $`[aab]`$, which we shall call type I, e.g., $`[123]`$ and they all have the same area $$V_2(\sigma _2^{Ib})=\frac{b}{2}\sqrt{\frac{a}{b}\frac{1}{4}}$$ (17) and four boundary triangles of the kind $`[bbb]`$, type II, e.g., $`[234]`$ whose area is $$V_2(\sigma _2^{IIb})=\frac{\sqrt{3}}{4}b$$ (18) * There are also two kinds of interior triangles. Type I includes the ones of the kind $`[as_is_i]`$, e.g., $`[012]`$. They all the same area $$V_2(\sigma _2^{Ii})=\frac{a}{2}\sqrt{\frac{s_i}{a}\frac{1}{4}}$$ (19) There are six more of type II, $`[bs_is_i]`$, like $`[023]`$, whose area is $$V_2(\sigma _2^{IIi})=\frac{b}{2}\sqrt{\frac{s_i}{b}\frac{1}{4}}$$ (20) For simplicity from now on we shall be working with rescaled metric variables: $$\eta =\frac{s_i}{b}$$ (21) $$\alpha =\frac{a}{b}$$ (22) $$T=\frac{H^2b}{l^2}$$ (23) where $`H^2=l^2\mathrm{\Lambda }/3`$, and $`l^2=16\pi G`$ is the Planck length. We shall work in units where $`c=\mathrm{}=1`$. Since we have two types of interior and boundary triangles we can expect at least four different deficit angles. The deficit angle associated with all eight interior triangles of type I is the same: $$\theta (\sigma _2^{Ii})=2\pi 3\mathrm{arccos}\left\{\frac{1}{2}\frac{(4\alpha 2)\eta \alpha ^2}{(4\alpha 1)\eta \alpha ^2}\right\}$$ (24) Similarly, all six interior triangles of type II have the same deficit angle $$\theta (\sigma _2^{IIi})=2\pi 4\mathrm{arccos}\left\{\frac{2\eta \alpha }{2\sqrt{3\eta 1}\sqrt{(4\alpha 1)\eta \alpha ^2}}\right\}$$ (25) As for the deficit angles associated with the boundary triangles, we can show that all twelve boundary triangles of type I have the same deficit angle $$\psi (\sigma _2^{Ib})=\pi 2\mathrm{arccos}\left\{\frac{\alpha }{2\sqrt{3\alpha 1}\sqrt{(4\alpha 1)\eta \alpha ^2}}\right\}$$ (26) the remaining four boundary triangles of type II also have the same deficit angle $$\psi (\sigma _2^{IIb})=\pi 2\mathrm{arccos}\left\{\frac{(3\alpha 2)}{2\sqrt{3\alpha 1}\sqrt{3\eta 1}}\right\}$$ (27) We can now write the explicit expression of the Regge action for this model: $`I[\eta ,T,\alpha ]`$ $`=`$ $`{\displaystyle \frac{T}{H^2}}\{8\sqrt{\alpha }\sqrt{\eta {\displaystyle \frac{\alpha }{4}}}[2\pi 3\mathrm{arccos}\left({\displaystyle \frac{1}{2}}{\displaystyle \frac{(4\alpha 2)\eta \alpha ^2}{(4\alpha 1)\eta \alpha ^2}}\right)]`$ $`+`$ $`6\sqrt{\eta 1/4}\left[2\pi 4\mathrm{arccos}\left({\displaystyle \frac{2\eta \alpha }{2\sqrt{3\eta 1}\sqrt{(4\alpha 1)\eta \alpha ^2}}}\right)\right]`$ $`+`$ $`6\sqrt{4\alpha 1}\left[\pi 2\mathrm{arccos}\left({\displaystyle \frac{\alpha }{2\sqrt{3\alpha 1}\sqrt{(4\alpha 1)\eta \alpha ^2}}}\right)\right]`$ $`+`$ $`2\sqrt{3}\left[\pi 2\mathrm{arccos}\left({\displaystyle \frac{(3\alpha 2)}{2\sqrt{3\alpha 1}\sqrt{3\eta 1}}}\right)\right]`$ $`+`$ $`{\displaystyle \frac{T^2}{H^2}}\left\{\alpha \sqrt{{\displaystyle \frac{(3\alpha 1)\eta }{\alpha ^2}}{\displaystyle \frac{3}{4}}}\right\}`$ ### 4.1 Analytic Study of The Action For obvious physical reasons we require the boundary edge lengths $`a`$ and $`b`$, and thus $`T`$ and $`\alpha `$, to be real and positive. On the other hand since we are only interested in geometries in which the boundary three-metric is positive definite, we must require that the volume of the boundary three-simplices be positive which is equivalent to requiring that $`\alpha >\frac{1}{3}`$. However, as pointed out in , for it to be possible for the wavefunction of the universe to predict a classical Lorentzian Universe the integration contour considered must be over complex metrics. Thus we are lead to consider $`\eta `$ as being a complex variable, and the analytic study of the action as a function of a complex variable $`\eta `$ is essential. The first thing to do is to identify the branch points of the action. Terms like $`\sqrt{zz_0}`$, are double-valued and have a square-root branch point at $`z=z_0`$. In order to make this term continuous we need to cut the complex plane. The most common branch cut is simply $`(\mathrm{},z_0]`$. So we see that the Euclidean action has $`3`$ square-root branch points located at: $$\eta _0=\frac{3}{4}\frac{\alpha ^2}{3\alpha 1}$$ $$\eta _2=\frac{\alpha }{4}$$ $$\eta _3=\frac{1}{4}$$ On the other hand a term like $`\mathrm{arccos}u(z)`$, is infinitely many-valued and has branch points at $`u(z)=+1,1`$, and at $`u(z)=\mathrm{}`$. The associated branch cuts are usually taken to be $`(\mathrm{},1][1,+\mathrm{})`$. Since $$\mathrm{arccos}u(z)=i\mathrm{log}\left(u(z)+\sqrt{u(z)^21}\right)$$ then we see that there are logarithmic singularities when $`u(z)=\mathrm{}`$. The table below shows the logarithmic branching points and infinities associated with the dihedral angles, where $$\eta _1=\frac{\alpha ^2}{4\alpha 1}$$ \| Dihedral angles \| $`+1`$ \| $`1`$ \| $`\mathrm{}`$ \| \| --- \| --- \| --- \| --- \| \| $`\theta (\sigma _2^{Ii})`$ \| $`1/4`$, $`\eta _0`$ \| $`1/4`$, $`\eta _0`$ \| $`1/3`$ \| \| $`\theta (\sigma _2^{IIi})`$ \| $`\alpha /4`$ \| $`\eta _0`$ \| $`\eta _1`$ \| \| $`\theta (\sigma _2^{Ib})`$ \| $`\eta _0`$ \| $`\eta _0`$ \| $`\eta _1`$ \| \| $`\theta (\sigma _2^{IIb})`$ \| $`\eta _0`$ \| $`\eta _0`$ \| $`1/3`$ \| In order to determine which branch cuts to take we need to know the relative values of these critical points, but those are dependent on the value of the boundary parameter $`\alpha `$. There are however some relations that are always valid for all $`\alpha >1/3`$ $$\eta _0>\frac{\alpha }{4},\frac{1}{3}$$ $$\eta _1>\frac{1}{4}$$ $$\eta _0>\eta _1$$ $$\eta _1>\frac{\alpha }{4}$$ So we see that the only indeterminates are $`min(\eta _1,\frac{1}{3})`$, $`min(\frac{\alpha }{4},\frac{1}{4})`$ and $`min(\frac{\alpha }{4},\frac{1}{3})`$. However, it is easy to see that the first two are just the same $`min(\alpha ,1)`$, and the third can of course be rewritten as $`min(\alpha ,\frac{4}{3})`$. Thus, we have three different regions * $`\frac{1}{3}<\alpha <1`$, where $`\frac{\alpha }{4}<\frac{1}{4}<\eta _1<\frac{1}{3}<\eta _0`$ * $`1<\alpha <\frac{4}{3}`$, where $`\frac{1}{4}<\frac{\alpha }{4}<\frac{1}{3}<\eta _1<\eta _0`$ * $`\frac{4}{3}<\alpha <+\mathrm{}`$, where $`\frac{1}{4}<\frac{1}{3}<\frac{\alpha }{4}<\eta _1<\eta _0`$ So when we consider all the cuts associated to all the terms we shall take as total branch cut $`(\mathrm{},\eta _0]`$. This also takes care of the logarithmic singularity at $`\eta =1/3`$. This cut leads to a continuous action function, in the complex plane with a cut $`(\mathrm{},\eta _0]`$. However this function is still infinitely many-valued. As in the previous chapters in order to remove this multivaluedness we redefine the domain where the action is defined, from the complex plane to the Riemann surface associated with $`I`$. The infinite multivaluedness of the action reflects itself in $`I`$ having an infinite number of branches with taking different values. The Riemann surface, $`𝐑`$ is composed by an infinite number of identical sheets, $`\text{ }\text{ }\mathrm{C}(\mathrm{},\frac{3}{8}]`$, one sheet for each branch of $`I`$. The first sheet $`\text{ }\text{ }\mathrm{C}_1`$ of the action is defined as the sheet where the terms in $`\mathrm{arccos}(z)`$ assume their principal values. That is the sheet where the action takes the form (4), and so $`I^I[\eta ,\alpha ,T]=I[\eta ,\alpha ,T]`$. When we continue the action in $`\xi `$ around one or more branch points, we will leave this first sheet and emerge in some other sheet of the Riemann surface. Only a few of these other sheets are relevant to us. When we encircle all five branch points we leave the first sheet and enter what we shall call the second sheet. It is easy to see that if we cross the branch cut $`(\mathrm{},\eta _0]`$, somewhere between $`\mathrm{}`$ and the smallest branch point ($`\alpha /4`$ or $`1/4`$, according to the value of $`\alpha `$) all terms of the action change their sign in this new sheet. Thus, the action in the second sheet differs by an overall minus sign relative to the action in the first sheet, $`I^{II}[\eta ,\alpha ,T]=I^I[\eta ,\alpha ,T]`$. However, we shall see that the steepest descents (SD) contours yielding the desired wavefunctions of the universe cross the branch cuts at other locations, thus emerging onto other sheets. For example, when $`\alpha >4/3`$, if we take a contour that encircles the branch point $`1/4`$, then at the first crossing of the branch cut, somewhere in $`(\mathrm{},1/4]`$, we go from the first sheet $`\text{ }\text{ }\mathrm{C}_1`$ to the second $`\text{ }\text{ }\mathrm{C}_2`$, and the action changes overall sign. Continuing in such a way that we again cross the branch cut, now between $`1/4`$ and the nearest branch point $`1/3`$, the contour enters what we shall call the third sheet. By carefully studying the behaviour of each term at the branch crossing, taking into account the branch cuts associated to each of them, we conclude that in this third sheet, $`\text{ }\text{ }\mathrm{C}_3`$, the action is $`I^{III}[\eta ,T,\alpha ]`$ $`=`$ $`{\displaystyle \frac{T}{H^2}}\{8\sqrt{\alpha }\sqrt{\eta {\displaystyle \frac{\alpha }{4}}}[2\pi 3\mathrm{arccos}\left({\displaystyle \frac{1}{2}}{\displaystyle \frac{(4\alpha 2)\eta \alpha ^2}{(4\alpha 1)\eta \alpha ^2}}\right)]`$ $`+`$ $`6\sqrt{\eta 1/4}\left[2\pi +4\mathrm{arccos}\left({\displaystyle \frac{2\eta \alpha }{2\sqrt{3\eta 1}\sqrt{(4\alpha 1)\eta \alpha ^2}}}\right)\right]`$ $`+`$ $`6\sqrt{4\alpha 1}\left[\pi 2\mathrm{arccos}\left({\displaystyle \frac{\alpha }{2\sqrt{3\alpha 1}\sqrt{(4\alpha 1)\eta \alpha ^2}}}\right)\right]`$ $`+`$ $`2\sqrt{3}\left[\pi 2\mathrm{arccos}\left({\displaystyle \frac{(3\alpha 2)}{2\sqrt{3\alpha 1}\sqrt{3\eta 1}}}\right)\right]`$ $`+`$ $`{\displaystyle \frac{T^2}{H^2}}\left\{\alpha \sqrt{{\displaystyle \frac{(3\alpha 1)\eta }{\alpha ^2}}{\displaystyle \frac{3}{4}}}\right\}`$ There is another important case, (still when $`\alpha >4/3`$), in which the SD contour encircles another branch point, namely $`\alpha /4`$. In this case we are confronted with a contour that starting on the first sheet, crosses the branch cut between $`1/3`$ and $`\alpha /4`$, encircles $`\alpha /4`$ and crosses the cut, once again somewhere between $`\alpha /4`$ and $`\eta _1`$. During the first crossing if we pay attention to the branch cuts in each term of the action, all terms of the action change their signs except the one associated with $`\mathrm{arccos}\theta (\sigma _b^{II})`$, and so we end up at what we shall call the fourth sheet $`\text{ }\text{ }\mathrm{C}_4`$, where the action is $`I^{IV}[\eta ,T,\alpha ]`$ $`=`$ $`{\displaystyle \frac{T}{H^2}}\{8\sqrt{\alpha }\sqrt{\eta {\displaystyle \frac{\alpha }{4}}}[2\pi 3\mathrm{arccos}\left({\displaystyle \frac{1}{2}}{\displaystyle \frac{(4\alpha 2)\eta \alpha ^2}{(4\alpha 1)\eta \alpha ^2}}\right)]`$ $`+`$ $`6\sqrt{\eta 1/4}\left[2\pi 4\mathrm{arccos}\left({\displaystyle \frac{2\eta \alpha }{2\sqrt{3\eta 1}\sqrt{(4\alpha 1)\eta \alpha ^2}}}\right)\right]`$ $`+`$ $`6\sqrt{4\alpha 1}\left[\pi 2\mathrm{arccos}\left({\displaystyle \frac{\alpha }{2\sqrt{3\alpha 1}\sqrt{(4\alpha 1)\eta \alpha ^2}}}\right)\right]`$ $`+`$ $`2\sqrt{3}\left[\pi 2\mathrm{arccos}\left({\displaystyle \frac{(3\alpha 2)}{2\sqrt{3\alpha 1}\sqrt{3\eta 1}}}\right)\right]`$ $``$ $`{\displaystyle \frac{T^2}{H^2}}\left\{\alpha \sqrt{{\displaystyle \frac{(3\alpha 1)\eta }{\alpha ^2}}{\displaystyle \frac{3}{4}}}\right\}`$ We then proceed to the second branch crossing somewhere between $`\alpha /4`$ and $`\eta _1`$, from where we emerge onto yet another sheet, $`\text{ }\text{ }\mathrm{C}_5`$. Once more if we pay attention to the branch cuts in each term of the action we can conclude that the action in this fifth sheet takes the form $`I^V[\eta ,T,\alpha ]`$ $`=`$ $`{\displaystyle \frac{T}{H^2}}\{8\sqrt{\alpha }\sqrt{\eta {\displaystyle \frac{\alpha }{4}}}[2\pi +3\mathrm{arccos}\left({\displaystyle \frac{1}{2}}{\displaystyle \frac{(4\alpha 2)\eta \alpha ^2}{(4\alpha 1)\eta \alpha ^2}}\right)]`$ $`+`$ $`6\sqrt{\eta 1/4}\left[2\pi 4\mathrm{arccos}\left({\displaystyle \frac{2\eta \alpha }{2\sqrt{3\eta 1}\sqrt{(4\alpha 1)\eta \alpha ^2}}}\right)\right]`$ $`+`$ $`6\sqrt{4\alpha 1}\left[\pi 2\mathrm{arccos}\left({\displaystyle \frac{\alpha }{2\sqrt{3\alpha 1}\sqrt{(4\alpha 1)\eta \alpha ^2}}}\right)\right]`$ $`+`$ $`2\sqrt{3}\left[\pi 2\mathrm{arccos}\left({\displaystyle \frac{(3\alpha 2)}{2\sqrt{3\alpha 1}\sqrt{3\eta 1}}}\right)\right]`$ $`+`$ $`{\displaystyle \frac{T^2}{H^2}}\left\{\alpha \sqrt{{\displaystyle \frac{(3\alpha 1)\eta }{\alpha ^2}}{\displaystyle \frac{3}{4}}}\right\}`$ We have just examined two cases that will be relevant when dealing with the SD contours, however, the lesson to be taken is that given the high number of branch points, the Riemann surface of the action is highly non-trivial and we must consider the changing form of the action along its sheets. ### 4.2 Asymptotic Behaviour of the Action Once more the asymptotic behaviour of the action with respect to the integration variable, $`\eta `$, is essential to the study of the convergence of the path integral yielding the wavefunction. Only after we know this asymptotic behaviour can we guarantee that the steepest descents contour is indeed dominated by the correct classical solutions. However from what we have seen in the previous chapter we must study this asymptotic behaviour not just in the first sheet but in all others where the SD contours are liable to go to infinity. It is easy to see that in the first sheet, the asymptotic behaviour of the action is $$I^I[\eta \mathrm{},\alpha ,T]\frac{\sqrt{3\alpha 1}}{H^2}T\left[TT_{crit}^I(\alpha )\right]\sqrt{\eta }$$ (32) where $`T_{crit}^I(\alpha )`$ $`=`$ $`{\displaystyle \frac{8\sqrt{\alpha }}{\sqrt{3\alpha 1}}}\left[2\pi 3\mathrm{arccos}\left({\displaystyle \frac{2\alpha 1}{4\alpha 1}}\right)\right]`$ $`+`$ $`{\displaystyle \frac{6}{\sqrt{3\alpha 1}}}\left[2\pi 4\mathrm{arccos}\left({\displaystyle \frac{1}{\sqrt{3}\sqrt{4\alpha 1}}}\right)\right];`$ for the third sheet we have $$I^{III}[\eta \mathrm{},\alpha ,T]\frac{\sqrt{3\alpha 1}}{H^2}T(TT_{crit}^{III})\sqrt{\eta }$$ (34) where now $`T_{crit}^{III}(\alpha )`$ $`=`$ $`{\displaystyle \frac{8\sqrt{\alpha }}{\sqrt{3\alpha 1}}}\left[2\pi 3\mathrm{arccos}\left({\displaystyle \frac{2\alpha 1}{4\alpha 1}}\right)\right]`$ $``$ $`{\displaystyle \frac{6}{\sqrt{3\alpha 1}}}\left[2\pi +4\mathrm{arccos}\left({\displaystyle \frac{1}{\sqrt{3}\sqrt{4\alpha 1}}}\right)\right]`$ and finally for the fifth sheet $$I^V[\eta \mathrm{},\alpha ,T]\frac{\sqrt{3\alpha 1}}{H^2}T(TT_{crit}^V)\sqrt{\eta }$$ (36) where now $`T_{crit}^V(\alpha )`$ $`=`$ $`{\displaystyle \frac{8\sqrt{\alpha }}{\sqrt{3\alpha 1}}}\left[2\pi +3\mathrm{arccos}\left({\displaystyle \frac{2\alpha 1}{4\alpha 1}}\right)\right]`$ $`+`$ $`{\displaystyle \frac{6}{\sqrt{3\alpha 1}}}\left[2\pi 4\mathrm{arccos}\left({\displaystyle \frac{1}{\sqrt{3}\sqrt{4\alpha 1}}}\right)\right].`$ In figure $`4`$ we see the plot of $`T_{crit}^I(\alpha )`$. It becomes infinite as $`\alpha `$ approaches $`1/3`$, which corresponds to the vanishing of the volume of $`\sigma _3^{Ib}`$. However, for all other values of $`\alpha `$ away from $`1/3`$ it quickly settles down at its asymptotic value $`Tcrit^I(+\mathrm{})=14.51`$. In the case of $`T_{crit}^{III}`$ it becomes $`\mathrm{}`$ when $`\alpha 1/3`$, and as $`\alpha `$ increases it becomes positive, having the same asymptotic value as $`T_{crit}^I`$, which is its upper bound. As for $`T_{crit}^V`$ it is always negative, whatever the value of $`\alpha `$. ## 5 Classical Solutions Since the complex $`W^4`$ is composed by eight equivalent $`4`$simplices, the metric of the complex will coincide with the metric of each simplex. The simplicial metric is then $$g_{ij}(\{s_k\})=\frac{s_{0i}+s_{0j}s_{ij}}{2}$$ (38) where $`ij=1,2,3,4,5`$. If we now calculate the eigenvalues of $`g_{ij}`$, we get $`(\frac{1}{2},\frac{1}{2},\lambda _{},\lambda _+)`$, where $$\lambda _+=2\eta \frac{1}{2}+\frac{1}{2}\sqrt{16\eta ^24\eta +112\eta \alpha +3\alpha ^2}$$ (39) $$\lambda _{}=2\eta \frac{1}{2}\frac{1}{2}\sqrt{16\eta ^24\eta +112\eta \alpha +3\alpha ^2}$$ (40) It is easy to see that since $`\alpha >1/3`$ then $`\lambda _+(\alpha ,\eta )`$ is always positive, whatever the value of $`\eta `$. For $`\lambda _{}`$, the situation is different. * If $`\eta >\eta _0`$ then $`\lambda _{}(\alpha ,\eta )>0`$ * If $`\eta <\eta _0`$ then $`\lambda _{}(\alpha ,\eta )<0`$ So we see that the complex $`W^4`$ will have Euclidean signature $`(++++)`$, if $`\eta >\eta _0`$, and will have Lorentzian signature $`(+++)`$ if $`\eta <\eta _0`$. Since there is only one internal degree of freedom, $`\eta `$, there is only one classical equation $$\frac{I}{\eta }=0$$ (41) which takes the form $`T`$ $`=`$ $`8\sqrt{{\displaystyle \frac{\alpha }{3\alpha 1}}}\sqrt{{\displaystyle \frac{\eta \eta _0}{\eta \alpha /4}}}\left[2\pi 3\mathrm{arccos}\left({\displaystyle \frac{1}{2}}{\displaystyle \frac{(4\alpha 2)\eta \alpha ^2}{(4\alpha 1)\eta \alpha ^2}}\right)\right]`$ $`+`$ $`{\displaystyle \frac{6}{\sqrt{3\alpha 1}}}\sqrt{{\displaystyle \frac{\eta \eta _0}{\eta 1/4}}}\left[2\pi 4\mathrm{arccos}\left({\displaystyle \frac{2\eta \alpha }{2\sqrt{3\eta 1}\sqrt{(4\alpha 1)\eta \alpha ^2}}}\right)\right]`$ Of course the classical solutions are of the form $`\eta =\eta _{cl}(\alpha ,T)`$ where $`T`$ and $`\alpha `$ are boundary data. However, it is obvious that the previous expression does not lead to a closed expression of $`\eta `$ as a function of $`\alpha `$ and $`T`$. So as in the previous chapters we plot the classical solutions as if $`T`$ was the dependent variable and not $`\eta `$. Of course this changes nothing, and it serves only as a visual aid to understanding the solutions. Note that the asymptotic behaviour of $`T(\alpha _{cl},\eta _{cl})`$, when $`\eta _{cl}+\mathrm{}`$ and when $`\eta _{cl}\mathrm{}`$, is the same $`T(\eta _{cl}\mathrm{})=T_{crit}^I(\alpha )`$ $`=`$ $`{\displaystyle \frac{8\sqrt{\alpha }}{\sqrt{3\alpha 1}}}\left[2\pi 3\mathrm{arccos}\left({\displaystyle \frac{2\alpha 1}{4\alpha 1}}\right)\right]`$ $`+`$ $`{\displaystyle \frac{6}{\sqrt{3\alpha 1}}}\left[2\pi 4\mathrm{arccos}\left({\displaystyle \frac{1}{\sqrt{3}\sqrt{4\alpha 1}}}\right)\right]`$ As in the previous chapters, and for exactly the same reasons, all classical solutions occur in pairs. Each element of the pair $`(\eta _{cl}^I,\eta _{cl}^{II})`$, is numerically equal $`\eta _{cl}^I=\eta _{cl}^{II}`$, but is located respectively in the first and second sheets of the action, and so they have Euclidean actions of opposite sign $`I^I[\eta _{cl}^I]=I^{II}[\eta _{cl}^{II}]`$. For physical reasons we only consider solutions for which $`T`$ is real positive and $`\alpha [1/3,+\mathrm{})`$. By plotting the solutions we conclude that there are three different cases according to the value of $`\alpha `$. * CASE 1 corresponds to $`\frac{1}{3}<\alpha <1`$. As we can see in figure $`5`$ , in this interval there are real Lorentzian solutions for $`\eta (\mathrm{},\alpha /4]`$, and real Euclidean solutions for $`\eta [\eta _0,+\mathrm{}]`$. Note however that although it is not apparent in figure $`5`$ the Euclidean solutions only obey the physical condition $`T>0`$ from a value of $`\eta `$ slightly larger than $`\eta _0`$. For $`\eta [\alpha /4,\eta _0]`$ and for all non-real values of $`\eta `$ the solutions have complex geometry. The divergence of $`T(\eta _{cl},\alpha )`$ at $`\alpha /4`$ guarantees that there are Lorentzian solutions for all values of $`T>T_{crit}^I`$. On the other hand, we have classical Euclidean solutions for all values between $`0`$ and $`T_{crit}(\alpha )`$. So the existence of Euclidean solutions is limited by a ceiling $`T=T_{crit}^I(\alpha )`$. However, since $`T_{crit}^I(\alpha 1/3)+\mathrm{}`$, in this range of $`\alpha `$ we can pretty much envisage raising the ceiling as much as we want, and so there will be Euclidean classical solutions for every value of $`T`$, it is just a case of getting $`\alpha `$ ever closer to $`1/3`$. Furthermore, note that as $`\alpha 1/3`$ the volume of the boundary tetrahedra $$V_3[\sigma _3^b]=\frac{\sqrt{3}l^3}{12H^3}T_{crit}^{3/2}(\alpha )\sqrt{\alpha 1/3}$$ does not vanish because $`T_{crit}(\alpha 1/3)(\alpha 1/3)^{1/4}`$. Regarding the behaviour of the action for these classical solutions it is easy to see that the Euclidean action of the Lorentzian solutions is pure imaginary. On the other hand, the Euclidean solutions have real Euclidean action. * CASE 2 corresponds to $`1<\alpha <4/3`$. The situation here is very similar to that of CASE $`1`$, the only difference being that now the relevant critical point where the Lorentzian classical branch peaks is $`1/4`$. As before we have Lorentzian classical solutions for values of $`T`$ from $`T_{crit}`$ all the way up to $`+\mathrm{}`$, while the existence of Euclidean classical solutions is limited to $`T[0,T_{crit}(\alpha )]`$. However, in this case the ceiling $`T_{crit}(\alpha )`$ is much more effective because in the range $`1<\alpha <4/3`$, $`T_{crit}(\alpha )`$ takes low values and there is no possibility of making it as high as we want by taking the limit $`\alpha 1/3`$, as in the previous case. In what concerns the action the situation is also similar to case $`1`$. The Euclidean solutions all have real Euclidean action. As for the Lorentzian solutions they all have pure imaginary actions. * CASE 3 corresponds to $`\alpha >1`$. As we can see from figure $`6`$ the situation here is different. There are two different branches of classical Lorentzian solutions $`\eta _1(T,\alpha )`$ and $`\eta _2(T,\alpha )`$ peaking respectively at the branch points $`\eta =1/4`$ and $`\eta =\alpha /4`$. We also have the usual classical Euclidean branch, for $`\eta [\eta _0,\mathrm{})`$. This time the Euclidean and Lorentzian regimes are not strictly separated at the “boundary” $`T=T_{crit}(\alpha )`$. This is because although the Euclidean solutions are still limited to boundary data $`T<T_{crit}(\alpha )`$, there is now one branch of Lorentzian solutions $`\eta _2(T,\alpha )`$, that contains solutions for all positive values of the boundary data $`T`$, including $`T<T_{crit}(\alpha )`$. Another big difference is that while the Lorentzian solutions in the $`\eta _1(T,\alpha )`$ branch all have real Euclidean actions, diverging to $`+\mathrm{}`$, as $`\eta 1/4`$, the solutions in the other Lorentzian branch $`\eta _2(T,\alpha )`$ have fully complex actions. The most important feature of the Euclidean action of these solutions is its behaviour near the branch point $`\alpha /4`$. $$Im[I(\eta _{cl}\frac{\alpha }{4}^{},T,\alpha )]+\mathrm{},Re[I(\eta _{cl}\frac{\alpha }{4}^{},T,\alpha )]\mathrm{};$$ the behaviour of $`ReI(\eta _{cl})`$ can be seen in figure $`7`$. As for the Euclidean solutions, their behaviour is similar to that of the two previous cases. There are real Euclidean solutions with real Euclidean action for all the range $`[\eta _0,\mathrm{})`$, but they only correspond to positive values of $`T`$ for values of $`\eta `$ slightly larger than $`\eta _0`$, and they are limited to a maximum value, $`T_{crit}(\alpha )`$, of $`T`$ . In the case of $`\alpha =2`$, and so $`\eta _0=0.6`$, the real Euclidean solutions only correspond to $`T>0`$ from approximately $`\eta =0.631`$. Note that we have Lorentzian classical solutions for the late universe for all possible values of $`\alpha `$. To see this remember that by late universe we mean very large boundary $`3`$spaces, which in our model corresponds to $`a=\alpha (Tl^2/H^2)+\mathrm{}`$. However, since we are restricted to $`\alpha >1/3`$, then $`a+\mathrm{}T+\mathrm{}`$. ## 6 Steepest Descents Contour In this minisuperspace model the wavefunction of the Universe is $$\mathrm{\Psi }[T,\alpha ]=_CD\eta e^{I[\eta ,T,\alpha ]}$$ (44) Once more this expression is only heuristic until we choose the integration contour $`C`$ and the integration measure $`D\eta `$. As in the previous models these choices are largely independent, if we restrict our attention to polynomial measures like $$D\eta =\frac{ds_i}{2\pi il^2}=\frac{T}{2\pi iH^2}d\eta $$ (45) In accordance with the previous chapters we shall follow Hartle’s prescription for the integration contour, namely, that the correct integration contour is the steepest descents contour over complex metrics passing through the dominant extrema of the Euclidean action. We know that all classical solutions occur in pairs, $`\eta _{cl}(\alpha ,T)=\eta _{cl}^I(\alpha ,T)=\eta _{cl}^{II}(\alpha ,T)`$, with Euclidean actions of opposite sign since they are located respectively in the first and second sheets of the action. Since by definition all points in a SD path have the same imaginary part of the action we see that no single SD path can pass through both extrema. However, if we consider a SD contour made of two complex conjugate sections one existing on the first sheet the other on the second , since, : $$I[\eta ]=[I[\eta ^{}]]^{}$$ and $$I[\eta _I]=I[\eta _{II}]$$ where $``$ denotes complex conjugation, we see that if one section goes through $`\eta _{cl}^I(\alpha ,T)`$ the other will go through $`\eta _{cl}^{II}(\alpha ,T)`$, and the resulting wavefunction will be real if the actions are purely imaginary or purely real. Another way of seeing this is to remember that we are now working for $`\eta 𝐑`$, and not $`\eta \text{ }\text{ }\mathrm{C}`$. As such, the SD contour passing through a value of $`\eta _{cl}(\alpha ,T)`$ is $$C_{SD}(T,\alpha )=\{\eta 𝐑:Im[I(\eta ,T,\alpha )]=Im[I(\eta _{cl}(T,\alpha ),T,\alpha )]\}$$ (46) it is then easy to conclude from our knowledge of the behaviour of the action in the several sheets of $`𝐑`$ that if $`C_{SD}`$ passes through $`\eta _{cl}^I`$ in $`\text{ }\text{ }\mathrm{C}_1`$ it will pass through $`\eta _{cl}^{II}`$in $`\text{ }\text{ }\mathrm{C}_2`$. For brevity we shall display only the SD contours associated with the classical solutions in the range $`\alpha >4/3`$, ie., case $`3`$, because the contours associated with the cases $`1`$ and $`2`$ are similar to the contours associated with case $`3`$. In figure $`8`$ we present the SD contour passing through the classical Lorentzian solution $`\eta _{cl}^1(T=100,\alpha =2)=0.2401`$, belonging to the first branch of Lorentzian solutions that peaks at $`\eta =1/4`$. Starting off at the real Lorentzian classical solution $`\eta _{cl}(T=100,\alpha =1)=0.2401`$, if we move upwards the SD contour goes to infinity in the first quadrant along the parabola $$\frac{\sqrt{3\alpha 1}}{H^2}T[TT_{crit}^I(\alpha )]\sqrt{\eta }=\stackrel{~}{I}[\eta _{cl}=0.2401]=265.71$$ (47) where $`\stackrel{~}{I}=Im(I)`$. The convergence of the integral along this part of the SD contour is guaranteed by the asymptotic behaviour of the action $$Re[I^I(\eta \mathrm{},\alpha ,T)]\frac{\sqrt{3\alpha 1}}{H^2}T[TT_{crit}^I(\alpha )]\sqrt{\eta }$$ (48) since for this branch of Lorentzian solutions we always have $`T>T_{crit}^I(\alpha )`$, and in particular $`T_{crit}^I(\alpha =2)=17.03`$. If instead we move downwards from the classical solution we immediately cross onto the second sheet of the action, where the SD contour encircles the branch point $`\eta =1/4`$ crossing the branch cut $`(\mathrm{},\eta _0]`$ once more, this time between $`1/4`$ and $`\alpha /4=1/2`$, at $`\eta =0.2632`$. When it does so it emerges onto the third sheet of the action where it goes to infinity in the first quadrant along the parabola $$\frac{\sqrt{3\alpha 1}}{H^2}T[TT_{crit}^{III}(\alpha )]\sqrt{\eta }=\stackrel{~}{I}[\eta _{cl}=0.2401]=265.71$$ (49) Once more the convergence of the integral along this part of the SD contour is guaranteed by the asymptotic behaviour of the action in the third sheet $$Re[I^{III}(\eta \mathrm{},\alpha ,T)]\frac{\sqrt{3\alpha 1}}{H^2}T[TT_{crit}^{III}(\alpha )]\sqrt{\eta }$$ (50) since $`T_{crit}^{III}`$ is always small or negative. In this particular case $`T_{crit}^I(\alpha =2)=16.69`$ The other section of this SD contour passing through the other extremum located at the second sheet is just the complex conjugate of this contour and so we will not show it here. We now study the SD contours associated to the second branch of classical Lorentzian solutions, i.e., the one that peaks at $`\eta =\alpha /4`$. In figure $`9`$ we display the SD contour associated with the Lorentzian classical solution $`\eta _{cl}^2(T=100,\alpha =0.4931)`$. Moving upwards we again go to infinity in the first quadrant of the first sheet along the parabola $$\frac{\sqrt{3\alpha 1}}{H^2}T[TT_{crit}^I(\alpha )]\sqrt{\eta }=\stackrel{~}{I}[\eta _{cl}=0.4931]=165.53$$ (51) The convergence of the integral is guaranteed by the asymptotic behaviour of the action in the first sheet $$Re[I^I(\eta \mathrm{},\alpha ,T)]\frac{\sqrt{3\alpha 1}}{H^2}T[TT_{crit}^I(\alpha )]\sqrt{\eta }$$ (52) given that $`T_{crit}(\alpha =2)=17.03`$. However, if we move downwards we cross the branch cut and emerge onto what we have defined as the fourth sheet of the action. There the SD contour encircles the branch point $`\alpha /4=0.5`$, and moves upward to cross the branch cut again between the branch points $`\alpha /4=0.5`$ and $`\eta _1(\alpha =2)=0.571`$, more specifically at $`\eta =0.5175`$. By doing this it moves onto the fifth sheet, where once more it goes to infinity in the first quadrant along the parabola $$\frac{\sqrt{3\alpha 1}}{H^2}T[TT_{crit}^V(\alpha )]\sqrt{\eta }=\stackrel{~}{I}[\eta _{cl}=0.4931]=165.53$$ (53) The convergence of the SD contour along this section is assured by the asymptotic behaviour of the action in this fifth sheet $$Re[I^V(\eta \mathrm{},\alpha ,T)]\frac{\sqrt{3\alpha 1}}{H^2}T[TT_{crit}^V(\alpha )]\sqrt{\eta }$$ (54) This is because $`T_{crit}^V(\alpha )`$ is always negative. In this particular case $`T_{crit}^V(2)=46.55`$. In the case of the Euclidean classical solutions they also occur in pairs but since they both have real Euclidean action they can both lie in the same SD path because the imaginary part of the action is equal for both of them, i.e., zero. In figure $`10`$, we present the SD contour associated with the pair of Euclidean classical solutions, numerically equal to $$\eta _{cl}^{sheetI}(T=13,\alpha =2)=\eta _{cl}^{sheetII}(T=13,\alpha =2)=1.180$$ but located in the first and second sheets, and so with real Euclidean actions of opposite sign, $`I[\eta _{cl}^{sheetI}(T=13,\alpha =2)]=I[\eta _{cl}^{sheetII}(T=13,\alpha =2)]=4.033`$. We see that starting at the classical solution $`\eta _{cl}^{sheetI}(T=13,\alpha =2)=1.180`$, if we move upwards in the first sheet the SD contour encircles all branch points and crosses the branch cut $`(\mathrm{},\eta _0=0.6]`$ at $`\eta =0.2638`$, emerging onto the second sheet where it encircles all branch points in the opposite direction and arrives at the second classical solution $`\eta _{cl}^{sheetII}(T=13,\alpha =2)=1.180`$. If we do another loop we will cross back onto the first sheet and will arrive back where we started, i.e., at the classical solution in the first sheet $`\eta _{cl}^{sheetI}(T=13,\alpha =2)=1.180`$. So the SD contour is clearly closed. Since there are no other critical points of the action in this SD contour we do not need to worry about the contribution (to the path integral) of any more points other than the two classical solutions. ## 7 Semiclassical Wavefunctions We have proved that for any classical solution there is always an SD contour passing through it. Furthermore by analysis of the asymptotic behaviour of the action we have been able to establish that any path integral along such SD contours will always be convergent, since none of them crosses any singularities and the contribution from the infinities is vanishing. However this is not enough to prove that the semiclassical approximation is always a good approximation of the SD wavefunction. In order to do that we must prove that the contributions from the regions about the classical solutions are clearly dominant, when compared with the contribution from the rest of the SD contour. For this to happen the values of $`T`$, $`\alpha `$ and $`H`$ must be such that, locally, the integrand, i.e., $`\mathrm{exp}I`$ is sharply peaked about the extrema. In our case the classical solutions always occur in pairs with actions of opposite sign. If these solutions are Lorentzian and have pure imaginary Euclidean actions, $`I[\eta _{cl},T,\alpha ]=ImI[[\eta _{cl},T,\alpha ]`$, then for the linear CPT symmetric wavefunction of the no boundary proposal we expect the real combination of these two contributions. So the semiclassical wavefunction takes the form $$\mathrm{\Psi }_{SC}(T,\alpha )\sqrt{\frac{T^2}{2\pi H^4I^{^{\prime \prime }}[\eta _{cl}(T,\alpha )]}}2\mathrm{cos}\left\{ImI[\eta _{cl}(T,\alpha ),T,\alpha ]\right\}$$ (55) where denotes $`d/d\eta `$. This is what happens for semiclassical approximations based on the the Lorentzian classical solutions in cases $`1`$ and $`2`$, i.e., when $`\alpha (1/3,1]`$, and $`\alpha [1,4/3]`$. It is also what happens in case $`3`$, $`(\alpha >4/3)`$, for the first branch of Lorentzian solutions peaking at $`\eta =1/4`$. In figures $`11`$ and $`12`$ we can see some examples of these wave functions. The semiclassical approximation is good when $`\mathrm{exp}(I)`$ is sharply peaked about the extrema, i.e., when $`ImI[\eta _{cl}(T,\alpha )]`$ is large. This will always be the case for the late universe . If for example we consider $`\alpha =1`$ then for the late universe where $`T+\mathrm{}`$ and $`\eta _{cl}^{Lor}1/4`$, the action takes the form $$ImI_{cl}(T,\alpha =1)\frac{T^2}{2H^2}=\frac{1}{2}(\mathrm{\Lambda }b^2/l^2)$$ (56) So the semiclassical approximation will be very good for large values of $`T`$. It will also be the case over the whole range of $`T`$, (except for $`T_{crit}`$), when $`H^2=\mathrm{\Lambda }l^2/3`$ is sufficiently small as it certainly is in our late universe. However we can see in figures $`11`$ and $`12`$ that the semiclassical wavefunctions diverge as $`TT_{crit}^I`$, but that is only a symptom that the semiclassical approximation breaks down there, which is signalled by the fact that $`I^{^{\prime \prime }}[\eta _{cl}(T,\alpha )0`$ when $`TT_{crit}^I`$. An oscillating wavefunction of the kind in figures $`11`$ and $`12`$ predicts classical Lorentzian spacetime for the late universe, described by the Lorentzian classical solutions computed above. What does it say about the evolution of the wormhole? Well since our late universe corresponds to large $`3`$boundaries, which is the case for our classical solutions when $`T=bH^2/l^2+\mathrm{}`$ and $`a=\alpha b+\mathrm{}`$, then it seems that the wormhole throat will grow with the expansion of the Lorentzian twin universes. So this wavefunction does not predict the collapse of the wormhole. When the classical solutions come in pairs of Euclidean solutions with real valued actions of opposite sign, as in the regions $`\eta [\eta _0,+\mathrm{})`$ where $`0<T<T_{crit}^I`$, the solution with negative real Euclidean action will always dominate over its counterpart. In our case the action of the Euclidean solutions in the first sheet is always negative, and so the solutions on the first sheet will dominate over the ones in the second sheet. However, as discussed above, the semiclassical approximation based on these solutions, will only be valid if the region near them gives the dominant contribution to the full SD contour integration. This only happens when the action is strongly peaked in the SD section around them. Since we are now limited to $`T<T_{crit}^I`$, this is only true when $`H^2`$ is small. And so the asymptotic behaviour of (44) for small $`H^2`$ is given by the semiclassical approximation associated with the Euclidean solution in the first sheet, $`\eta _{cl}^{sheetI}(T,\alpha )`$ $$\mathrm{\Psi }_{SC}^{Eucl}(T,\alpha )\sqrt{\frac{T^2}{2\pi H^4I^{^{\prime \prime }}[\eta _{cl}(T,\alpha )]}}e^{I[\eta _{cl}^{sheetI}(T,\alpha ),T,\alpha ]}$$ (57) The second derivative of the action vanishes as $`TT_{crit}^I`$, and so it will be this term that will dominate the semiclassical wavefunction near the “turning point”, $`T_{crit}^I`$, and lead to its divergence there. However, when $`H`$ is small, the $`exp(I)`$ term is dominant everywhere except when $`T`$ is really close to $`T_{crit}^I`$. For small $`H`$ the real valued Euclidean action peaks around the Euclidean classical solution and so the semiclassical approximation is good for almost all values of $`T`$, except when we get too close to $`T_{crit}^I`$. In figure $`13`$ we see the plot of this semiclassical wavefunction for $`\alpha =2`$ and $`H=4.9`$. There is a clear peak away from $`T_{crit}^I`$, at $`T=13.35`$, where the semiclassical approximation is valid, which seems to indicate a preferred value of the boundary edges $`b=Tl^2/H^2`$ and consequently $`a=b\alpha =2b`$. This peak becomes more pronounced as $`H`$ becomes smaller. We thus seem to have a prediction of a favoured size of the wormhole throat, when we are dealing with wormholes between Euclidean universes. We now discuss a third kind of classical solutions. When $`\alpha >4/3`$, we know that there are two branches of Lorentzian classical solutions, branch $`1`$ peaking at $`\eta =1/4`$ and branch $`2`$ peaking at $`\eta =\alpha /4`$. These solutions in the second branch have fully complex Euclidean actions. $$I[\eta _2^{cl}(T,\alpha )]=ReI[\eta _2^{cl}(T,\alpha )]+iImI[\eta _2^{cl}(T,\alpha )]$$ Their counterparts on the second sheet will also have complex actions, but of opposite sign. So no linear combination of the contribution from these two solutions will ever lead to a real wavefunction. Moreover since the real part of the Euclidean action is always negative for the classical solutions $`\eta _2^{cl}(T,\alpha )`$ on the first sheet, even diverging to $`\mathrm{}`$, as $`\eta \alpha /4`$, as can be seen in figure $`7`$, it seems obvious that the extremum in the first sheet will dominate the path integral. In order to understand a little better this situation we need to do a slight diversion, . The starting point of a semiclassical/WKB approximation is to assume that the full wavefunction is well approximated by the contribution of the regions about some finite number of classical extrema. In our case it is just one and so the wavefunction can always be written as $$\mathrm{\Psi }[h_{ij}]=C\mathrm{exp}\frac{I[h_{ij}]}{\mathrm{}}$$ (58) where $`I=I_r+iI_i`$, where $`I_r`$ and $`I_i`$ are the real and imaginary parts of the Euclidean action, and the pre-factor $`C`$ is a slowly varying amplitude containing all higher-order corrections. At the level of the first approximation in $`\mathrm{}`$ the Wheeler-DeWitt equation on $`\mathrm{\Psi }[h_{ij}]`$ is then reduced to the classical Hamilton-Jacobi equation on $`I`$, . $$H_0\left(\pi ^{ij}=i\frac{\delta I}{\delta h_{ij}},h_{ij}\right)=0$$ (59) $$[G_{ijkl}\frac{\delta I}{\delta h_{ij}}\frac{\delta I}{\delta h_{kl}}+\sqrt{h}(^3R2\mathrm{\Lambda })]=0$$ (60) where $`G_{ijkl}=\frac{1}{2\sqrt{h}}(h_{ik}h_{jl}+h_{il}h_{jk}h_{ij}h_{kl})`$ and the $`\pi ^{ij}`$ are the canonical conjugate momenta of $`h_{ij}`$. The action $`I[h_{ij}]`$ that solves the the Hamilton-Jacobi equation is the classical action of a congruence of solutions, $`[h_{ij}(t)]`$, of the classical equations of motion. These trajectories, parametrised by their coordinate time $`t`$, are defined by $$\pi ^{ij}=i\frac{\delta I}{\delta h_{ij}}$$ (61) If the action $`I`$ is real, the wavefunction has an exponential form and can be formally interpreted as an ensemble of classical Euclidean solutions/trajectories. If the action $`I`$ is pure imaginary, then the wavefunction has oscillatory behaviour, and describes an ensemble of classical Lorentzian solutions. However, if the action $`I`$ is fully complex the Hamilton-Jacobi equation can be decomposed into the real part $$(I_r)^2(I_i)^2+\sqrt{h}(^3R2\mathrm{\Lambda })=0$$ (62) and the imaginary part $$(I_r)\times (I_i)=0$$ (63) From 62, we can see that if the gradient of $`I_i`$ becomes much larger than the gradient of $`I_r`$, then $`I_i`$ will be an approximate solution to the Hamilton-Jacobi equation. This means that the classical evolution is almost Lorentzian. If $`I`$ defines an ensemble of almost Lorentzian solutions, then the real part of the action is also relevant, since it will contribute the term $`\mathrm{exp}(I_r)`$ to the pre-factor. Since it is exponential it will probably be the dominant contribution to the pre-factor. Then the probability measure provided by the pre-factor on the ensemble of classical solutions is, to leading order, of the form $`\mathrm{exp}(2I_r)`$, determining the relative weight of different trajectories. This is exactly what happens in the case of branch $`2`$ of classical solutions, $`\eta _2^{cl}(T,\alpha )`$ when $`\alpha >4/3`$. As can be seen in figure $`14`$, the ratio $`R`$ between the gradients of the imaginary and real parts of the Euclidean action $$R=\frac{ImI[\eta =\eta _2^{cl}]}{ReI[\eta =\eta _2^{cl}]}$$ (64) along the classical solutions of branch $`2`$, becomes infinitely large for the late universe, i.e., for $`a=\alpha (Tl^2/H^2)+\mathrm{}\eta _2^{cl}\frac{\alpha }{4}^{}`$. So for the late universe $`ImI`$ really solves the classical Hamilton-Jacobi equation. And so the wavefunction should describe an ensemble of Lorentzian solutions for the late universe. The contribution of the imaginary part of the action to the wavefunction does indeed yield the characteristic oscillating behaviour associated with the prediction of classical Lorentzian spacetime for the late universe, as can be seen in figure $`15`$. The real part of the action provides the probability measure, $`\mathrm{exp}\{2ReI[\eta _2^{cl}]\}`$, on the ensemble of classical Lorentzian universes, determining the relative weight of these different trajectories. From figure $`7`$, we see that $`\mathrm{exp}\{2ReI[\eta _2^{cl}]\}`$ becomes infinitely large as $`T+\mathrm{}`$, since $`ReI[\eta _2^{cl}]\mathrm{}`$ as $`\eta _2^{cl}\frac{\alpha }{4}^{}`$. It thus seems that in this model, when $`\alpha >4/3`$ if we choose the classical solutions of branch $`2`$ as the basis for the semiclassical approximation of the wavefunction of the universe, then configurations consisting of extremely large Lorentzian classical universes connected by similarly large wormholes will dominate. ## 8 Conclusions The minisuperspace approximation comes very naturally in simplicial quantum gravity, because by discretising spacetime we are in effect substituting the functional $`h_{ij}(𝐱,t)`$, by a set of simple variables, the edge lengths. The minisuperspace approximation consists only in restricting all these edge lengths to being equal to a small number of parameters. So the loss of generality is much smaller than when we impose the minisuperspace approximation to the continuum. Furthermore, in the simplicial minisuperspace, spacetime continues to be treated in a fully $`4`$dimensional way. There is no need to invoke any arbitrary $`3+1`$ decomposition of spacetime, usually ADM, as in the case of continuum models. This means that global issues like topology can still be addressed in the simplicial minisuperspace, which is not possible in the continuum versions. The versatility of our approach is evident in the ease in which one can construct a $`4`$D wormhole configuration. Furthermore, since the signature of the simplicial spacetime is only dependent on the ratio between edge lengths, we can deal with Lorentzian and Euclidean configurations at the same time. Indeed with the same model we were able to study microscopic Euclidean wormholes and large Lorentzian wormholes. The simplicial geometry used is quite simple and so it cannot take into account some relevant dynamics that we would like to model. In particular it would be better if we could have developed a solvable model with an additional interior edge length, that would describe the throat of the wormhole. However, we decided that before we considered that much more complicated model it would be very useful to determine the results of the simplest model. And indeed we found several new and interesting results. For Euclidean wormholes just above the Planck scale, there is a strongly preferred non-zero throat size. As $`H`$ becomes smaller the peak becomes sharper, and we obtain a very stable configuration. Thus the model predicts that Euclidean wormholes of that scale are stable insofar as that their throat do not tend to contract leading to the pinching off of the two universes. For large Lorentzian wormhole configurations the situation is dependent on the value of $`\alpha =a/b`$. If $`\alpha [1/3,4/3]`$, then there is only one family of Lorentzian solutions, all with pure imaginary Euclidean action. The semiclassical wavefunction is a very good approximation of the full wavefunction for the late, large Lorentzian universes. Its oscillating behaviour definitely predicts classical Lorentzian configurations for large universes, with a similarly large wormhole throat connecting them. However, when $`\alpha >4/3`$, a new family of Lorentzian solutions is present. The fact that their Euclidean action is fully complex, with a real part that diverges when the boundary $`3`$space becomes large, (i.e. late universe), seems to indicate that these solutions will dominate the path integral leading to the full wavefunction. A semiclassical interpretation of the resulting wavefunction is that it describes an ensemble of classical Lorentzian spacetimes, weighted by $`\mathrm{exp}(2ReI)`$. Since this exponential diverges to $`+\mathrm{}`$ for increasingly large $`3`$spaces, we can conclude that large wormhole configurations between very large Lorentzian universes are very strongly favoured.
warning/0004/hep-th0004027.html	ar5iv	text	# Evaporation of near-extremal Reissner-Nordström black holes ## Abstract The formation of near-extremal Reissner-Nordström black holes in the S-wave approximation can be described, near the event horizon, by an effective solvable model. The corresponding one-loop quantum theory remains solvable and allows to follow analytically the evaporation process which is shown to require an infinite amount of time. PACS number(s): 04.70.Dy, 04.62.+v Black holes are the most fascinating objects in General Relativity. Since Hawking discovered that they emit thermal radiation , it has been a long standing puzzle to explain their thermodynamical properties, in terms of some microscopic structure, and to understand the dynamical evolution beyond Hawking initial scheme where the gravitational field was treated as a fixed background. Extremal and near-extremal charged black holes have recently played a fundamental role in String Theory, where in some special cases it has been possible to give a statistical explanation of the Bekenstein-Hawking area law for their entropy . Moreover, the scattering of low-energy particles off extremal black holes also provides a convenient setting to study the evaporation process including back-reaction effects. By throwing a long-wavelength particle into an extremal hole, a non-extremal configuration is created and quantum-mechanically one expects it to decay via Hawking emission back to extremality. To render this problem tractable one can boil it down considering large $`q`$ and incoming neutral matter with zero angular momentum. In this context, dilatonic black holes were extensively considered since the scattering particle-hole and the ensuing information loss problem can be analyzed in an analytical framework. This is so because the problem can be reduced to study a two-dimensional effective theory which turns out to be solvable both at the classical and at the one-loop quantum level (see also the reviews for a more detailed description and as useful references for the methods used in this work). Among the ‘nice’ properties of these dilatonic black holes, making the problem under study rather special, the extremal black holes are classically completely regular and, moreover, the temperature near extremality is constant. In contrast, the Reissner-Nordström (RN) black holes always have singularities (extremal case $`lm=q`$ included) and the temperature goes as $`T_H\sqrt{lmq}/q^2`$ near extremality. This case was reconsidered after the improvements in the physical understanding obtained with dilatonic black holes. However, only partial analytic (obtained by means of the adiabatic approximation) and numerical answers have been obtained . The purpose of this letter is to present the first exact results on the evaporation process of a near-extremal RN black hole. If we consider the physics in the S-wave approximation, near the horizon we can describe it by an effective model (the Jackiw-Teitelboim model ) which is solvable even when back-reaction effects are included. Let us consider the solutions of Einstein-Maxwell gravity with null infalling matter $`d\overline{s}^2`$ $`=`$ $`{\displaystyle \frac{(rr_+)(rr_{})}{r^2}}dv^2+2drdv+r^2d\mathrm{\Omega }^2,`$ (1) $`F`$ $`=`$ $`qϵ_2,`$ (2) where $`ϵ_2`$ is the volume element of the unit S<sup>2</sup>. The only nonzero component of the stress tensor is given by $$T_{vv}=\frac{_v(r_+(v)+r_{}(v))}{8\pi l^2r^2},$$ (3) where $`l^2=G`$ is Newton’s gravitational constant. One can describe the formation of a non-extremal black hole by sending a low-energy shock wave $$T_{vv}=\frac{\mathrm{\Delta }m}{4\pi r^2}\delta (vv_0),$$ (4) in the extremal geometry ($`v<v_0`$). This model can be described by an effective two-dimensional theory given by $$I=d^2x\sqrt{g}\left[R\varphi +l^2V(\varphi )\frac{1}{2}\|f\|^2\right],$$ (5) where the field $`f`$ represents the null matter with $`(_vf)^2T_{vv}^f=4\pi r^2T_{vv}`$ and $$\varphi =\frac{r^2}{4l^2},V(\varphi )=(4\varphi )^{\frac{1}{2}}q^2(4\varphi )^{\frac{3}{2}}.$$ (6) The two-dimensional metric is related to the $`rv`$ projection of (1) by the conformal rescaling $$ds^2=\sqrt{\varphi }d\overline{s}^2.$$ (7) The extremal black hole is recovered for the zero of the potential $`V(\varphi _0)=0`$, which corresponds to $`\varphi _0=q^2/4`$. This fact suggests us to consider the effective near-horizon and near-extremal theory defined by the expansion of (5) around $`\varphi _0`$ $$\varphi =\varphi _0+\stackrel{~}{\varphi },m=lq+\mathrm{\Delta }m.$$ (8) We obtain $$I=d^2x\sqrt{g}\left[(R+\frac{4}{l^2q^3})\stackrel{~}{\varphi }\frac{1}{2}\|f\|^2\right]+𝒪(\stackrel{~}{\varphi }^2),$$ (9) where the leading order term is just the Jackiw-Teitelboim (JT) model, which now arises as the effective theory governing the dynamics near extremality and close to the horizon. A $`d=1`$ realization of the AdS$`{}_{d+1}{}^{}/`$CFT<sub>d</sub> correspondence in the JT model exactly accounts for the deviation from extremality of the Bekenstein-Hawking entropy of RN black holes . Therefore one can also expect to obtain an exact picture of the evaporation process near the horizon. The formation of a near-extremal black hole due to a shock wave (4) can be pushed down to the JT metrics of constant negative curvature. For $`v<v_0`$ we have the extremal RN configuration and its near-horizon geometry is given by the Robinson-Bertotti anti-de Sitter geometry $$d\overline{s}^2=\frac{r^2}{r_0^2}dt^2+\frac{r_0^2}{r^2}dr^2+r_0^2d\mathrm{\Omega }^2,$$ (10) where $`r_0`$ is the extremal radius. After the rescaling (7) the two-dimensional metric in null conformal coordinates becomes $$ds^2=\frac{2}{l^2q^3}\stackrel{~}{x}^2dudv,$$ (11) with $`u=v+\frac{l^2q^3}{\stackrel{~}{x}}`$, $`\stackrel{~}{x}=l\stackrel{~}{\varphi }`$. Proceeding in a similar way for the near-extremal configuration $`v>v_0`$, we obtain $$ds^2=(\frac{2}{l^2q^3}\stackrel{~}{x}^2l\mathrm{\Delta }m)d\overline{u}dv,$$ (12) with $$\overline{u}=v+\sqrt{\frac{2lq^3}{\mathrm{\Delta }m}}\mathrm{a}rctanh\sqrt{\frac{2}{l^3q^3\mathrm{\Delta }m}}\stackrel{~}{x}.$$ (13) The coordinates $`(v,u)`$ and $`(v,\overline{u})`$ are the radial null coordinates corresponding to those of RN $`(t+r^{},tr^{})`$ before and after the shock wave. Imposing the continuity of (11) and (12) along $`v=v_0`$ one obtains $$u=v_0+a\mathrm{c}otanh\frac{\overline{u}v_0}{a},$$ (14) where $`a=\sqrt{\frac{2lq^3}{\mathrm{\Delta }m}}`$. From this relation we can work out immediately the outgoing energy flux of Hawking radiation in terms of the Schwarzian derivative between the coordinates $`u`$ and $`\overline{u}`$ $$T_{\overline{u}\overline{u}}^f=\frac{1}{24\pi }\{u,\overline{u}\}=\frac{1}{12\pi a^2}=\frac{\pi }{12}T_H^2.$$ (15) We observe that this flux is constant and coincides with the thermal value of Hawking flux for near-extremal RN black holes, where $`T_H`$ is Hawking’s temperature. This fact can be understood easily since the AdS$`{}_{2}{}^{}\times `$S<sup>2</sup> geometries associated to (11) and (12) represent indeed the near-horizon limit of the RN geometries (1) due to the shock wave (4). So the constant thermal flux for every value of $`\overline{u}`$ corresponds to the flux measured by an inertial observer at future null infinity approaching the event horizon of the RN black hole. In the light of this remark it is interesting to point out that the JT theory also describes the dynamics of extremal and near-extremal RN black holes close to the horizon in the presence of a spherically symmetric Klein-Gordon field. This is so because, as it is well-known, the scalar field propagates freely near the horizon. Our purpose now is to analyze the back-reaction effects in the evaporation process. The one-loop effective theory is obtained by adding the non-local Liouville-Polyakov term to the classical action. We then get ($`\lambda ^2=l^2q^3`$) $`I`$ $`=`$ $`{\displaystyle d^2x\sqrt{g}\left[R\stackrel{~}{\varphi }+4\lambda ^2\stackrel{~}{\varphi }\frac{1}{2}\underset{i=1}{\overset{N}{}}\|f_i\|^2\right]}`$ (16) $``$ $`{\displaystyle \frac{N\mathrm{}}{96\pi }}{\displaystyle d^2x\sqrt{g}R\mathrm{}^1R}+\xi {\displaystyle \frac{N\mathrm{}}{12\pi }}{\displaystyle d^2x\sqrt{g}\lambda ^2},`$ (17) where we have considered the presence of $`N`$ scalar fields. The parameter $`N`$ allows us to consider the theory in the large $`N`$ limit, keeping $`N\mathrm{}`$ fixed. Moreover we have also added a local conterterm (in the form of a 2d cosmological constant which corresponds to the freedom of adding a constant to the 2d conformal anomaly), mimicking the analysis of dilaton gravity theory , to ensure that the extremal geometry remains an exact solution of the one-loop theory at $`\xi =1`$. Our results are the same irrespective of the value of $`\xi `$. We should mention now that for the region we are interested in, the conformal factor $`\sqrt{\varphi }`$ of (7) is almost constant and therefore the semiclassical quantization in terms of the Einstein-Maxwell action and the JT action are equivalent. Far from the horizon this is no more true. The equations of motion derived from (16) in conformal gauge $`ds^2=e^{2\rho }dx^+dx^{}`$ are (from now on we take $`\xi =1`$) $`2_+_{}\rho +\lambda ^2e^{2\rho }`$ $`=`$ $`0,`$ (18) $`_+_{}\stackrel{~}{\varphi }+\lambda ^2\stackrel{~}{\varphi }e^{2\rho }`$ $`=`$ $`0,`$ (19) $`_+_{}f_i`$ $`=`$ $`0,`$ (20) $`2_\pm ^2\stackrel{~}{\varphi }+4_\pm \rho _\pm \stackrel{~}{\varphi }`$ $`=`$ $`T_{\pm \pm }^f{\displaystyle \frac{N\mathrm{}}{12\pi }}t_\pm `$ (22) $`{\displaystyle \frac{N\mathrm{}}{12\pi }}\left((_\pm \rho )^2_\pm ^2\rho \right).`$ The functions $`t_\pm (x^\pm )`$ are related with the boundary conditions of the theory and depend on the quantum state of the system. The Liouville equation (18) has the general solution $$ds^2=\frac{_+A_+_{}A_{}}{(1+\frac{\lambda ^2}{2}A_+A_{})^2}dx^+dx^{},$$ (23) FIG.1. Kruskal diagram of near-extremal black hole. The two timelike boundaries of near-horizon geometry AdS<sub>2</sub> are represented by the vertical line $`x^{}=x^+`$. The infalling shock wave emerges from one boundary (left side of $`x^{}=x^+`$ line) and, crossing the outer and inner horizons, reaches the other boundary (right side of $`x^{}=x^+`$ line). where $`A_\pm (x^\pm )`$ are arbitrary chiral functions. We can choose a particular form of the functions $`A_\pm `$ as a way to fix completely the conformal coordinates. We find convenient to choose $$A_+=x^+,A_{}=\frac{2}{\lambda ^2x^{}}.$$ (24) Before the shock wave these coordinates $`x^\pm `$ correspond to the RN coordinates $`(v,u)`$ and after ($`v>v_0`$) the relation is $$v=x_0^++a\mathrm{a}rctanh\frac{x^+x_0^+}{a},$$ (25) together with $`u=x^{}`$. Then both metrics (11) and (12) are brought into (23) and the physical information is encoded in the field $`\stackrel{~}{\varphi }`$. At the classical level the solution for it is given by $$\stackrel{~}{\varphi }=lq^3\frac{1\mathrm{\Theta }(x^+x_0^+)\frac{\mathrm{\Delta }m}{2lq^3}(x^+x_0^+)(x^{}x_0^+)}{x^{}x^+}.$$ (26) After the shock wave ($`x^+>x_0^+`$) the extremal radius is given by curve $`\stackrel{~}{\varphi }=0`$: $`(x^+x_0^+)(x^{}x_0^+)=a^2`$, and the outer and inner apparent horizons $`r=r_\pm `$ are given by the condition $`_+\stackrel{~}{\varphi }=0`$: $`x^{}=x_0^+\pm a`$. The corresponding Kruskal diagram (in this region the coordinates $`x^\pm `$ are regular at the horizon and therefore they represent a sort of Kruskal frame) is given by Fig.1. At the quantum level we have to solve equations (18-22) and the crucial point is to choose the adequate functions $`t_\pm (x^\pm )`$ for the physical situation. The natural choice is $`t_v=t_{x^+}=0`$ and $`t_u=t_x^{}=0`$ before the shock-wave and $`t_{x^+}=\frac{1}{2}\{v,x^+\}`$ and $`t_x^{}=0`$ after, where now $`v`$ is the light-cone coordinate of the evaporating Vaidya-type metric $$ds^2=(\frac{2\stackrel{~}{x}^2}{l^2q^3}lm(v))dv^2+2d\stackrel{~}{x}dv.$$ (27) The remarkable property of the equations of the near-horizon effective theory is that one can solve them also in conformal gauge. We find that $$\stackrel{~}{\varphi }=\frac{F(x^+)}{x^{}x^+}+\frac{1}{2}F^{}(x^+),$$ (28) $`\stackrel{~}{x}=l\stackrel{~}{\varphi }`$, where the function $`F(x^+)`$ satisfies the following differential equation $$F^{\prime \prime \prime }=\frac{N\mathrm{}}{24\pi }\left(\frac{F^{\prime \prime }}{F}+\frac{1}{2}\left(\frac{F^{}}{F}\right)^2\right),$$ (29) and relates the $`x^+`$ and $`v`$ coordinates $$\frac{dv}{dx^+}=\frac{lq^3}{F}.$$ (30) The evaporating mass is then given by $$m(x^+)=\frac{24\pi }{N\mathrm{}lq^3}F^2F^{\prime \prime \prime },$$ (31) and can be related to the boundary function $`t_{x^+}`$ $$t_{x^+}=\frac{lq^3m(x^+)\mathrm{\Theta }(x^+x_0^+)}{2F^2}.$$ (32) The fact that the functions $`t_\pm `$ can be discontinuous for coordinates $`x^\pm `$ associated to free (or Liouville) fields was pointed out in . The expression (32), and also the function $`F(x^+)`$, admits a series expansion in powers of $`\mathrm{}`$, where the classical term, obtained using the classical relation (25), is given by $$t_{x^+}=\frac{a^2\mathrm{\Theta }(x^+x_0^+)}{(a^2(x^+x_0^+)^2)^2}.$$ (33) As before, the curve $`\stackrel{~}{\varphi }=0`$ $$x^{}=x^+\frac{2F}{F^{}},$$ (34) represents the location of the extremal radius and $`_+\stackrel{~}{\varphi }=0`$ defines the inner and outer apparent horizons in the spacetime of the evaporating black hole $$x^{}=x^+\frac{F^{}}{F^{\prime \prime }}\pm \frac{\sqrt{F^22FF^{\prime \prime }}}{F^{\prime \prime }}.$$ (35) The intersection of these three curves takes place when $$F^22FF^{\prime \prime }=2lq^3m(x^+)=0,$$ (36) i.e. at the end of the evaporation. On the other hand , it is easy to show that $`_+m(x^+)=\frac{N\mathrm{}}{24\pi }\frac{m(x^+)}{F}`$ which can be readily integrated in $`v`$ coordinate giving $`m(v)=\mathrm{\Delta }me^{\frac{N\mathrm{}}{24\pi lq^3}(vv_0)}`$ and so $`m=0`$ is given by $`v=+\mathrm{}`$ and not before (had we started with the classical boundary term (33) we would have obtained a finite evaporation time). From the numerical graph of the function $`F(x^+)`$ it is clear that this happens at a finite FIG.2. Kruskal diagram of semiclassical evolution of RN black hole. The outer apparent horizon shrinks until it meets both inner horizon and extremal radius curve at the AdS boundary. value of $`x^+=x_{int}^+`$ for which $`F(x_{int}^+)=0`$ and moreover from eq. (36) it also follows that $`F^{}(x_{int}^+)=0`$. One can also show that $`F^{\prime \prime }(x_{int}^+)`$ is nonzero while all the other derivatives vanish, so that locally close to the intersection point $`F(x^+)`$ behaves as a parabola with exponentially suppressed corrections. This is enough to prove , see (34) and (35), that the intersection point belongs to the AdS boundary, i.e. $`x_{int}^+=x_{int}^{}`$. The graph of all these curves is shown in Fig. 2, where the saddle point in the outer apparent horizon curve $`r_+`$ signals the transition from the strong to the weak back-reaction regimes as described in . At the end-point the curves $`\varphi =0`$ and $`_+\varphi =0`$ are null and the dilaton function is well represented asymptotically by the extremal solution $$\varphi =\frac{_+^2F(x_{int}^+)}{2}\frac{(x^+x_{\mathrm{i}nt}^+)(x^{}x_{\mathrm{i}nt}^+)}{x^{}x^+}+\mathrm{}.$$ (37) where the dots are nothing but exponentially small corrections. So, as time passes the solution becomes ’more and more extremal’ but without actually coming back to the extremal state due to the infinite evaporation time. Due to the fact that $`T_H=0`$ for the extremal black hole our exact results are in agreement with the third law of thermodynamics applied to black holes (see for instance ), despite the fact that the weak energy condition is violated close to the horizon, as well as with those obtained using the adiabatic approximation . This research has been partially supported by the CICYT and DGICYT, Spain. D. J. Navarro acknowledges the Ministerio de Educación y Cultura for a FPI fellowship. A.F. thanks R. Balbinot for useful discussions and the Department of Theoretical Physics of Valencia University for hospitality during the late stages of this work. D.J.N. and J.N-S. also wish to thank J. Cruz and P. Navarro for comments. Email address: fabbria@bo.infn.it Email address: dnavarro@ific.uv.es Email address: jnavarro@lie.uv.es
warning/0004/quant-ph0004018.html	ar5iv	text	# Calculating the relative entropy of entanglement ## Abstract We extend Vedral and Plenio’s theorem (theorem 3 in Phys. Rev. A 57, 1619) to a more general case, and obtain the relative entropy of entanglement for a class of mixed states, this result can also follow from Rains’ theorem (theorem 9 in Phys. Rev. A 60, 179). The relative entropy of entanglement plays an important role in understanding quantum entanglement , it reduces to the von Neumann reduced entropy of either side for the bipartite pure states ; for general mixed states, the relative entropy of entanglement is hard to calculate. In this note, we extend Vedral and Plenio’s result for pure states to a more general case, and calculate the relative entropy of entanglement for a class of mixed states. The relative entropy of entanglement for 2-party (say, A and B) quantum state $`\sigma `$ is defined by : $$Er(\sigma )\underset{\rho D}{\mathrm{min}}S(\sigma \|\|\rho )$$ (1) where $`D`$ is the set of all disentangled states, and $`S(\sigma \|\|\rho )tr\{\sigma (\mathrm{ln}\sigma \mathrm{ln}\rho )\}`$ is the relative entropy of $`\sigma `$ with respect to $`\rho `$. Let $`\rho ^{}`$ denotes the disentangled state that minimizes the relative entropy $`S(\sigma \|\|\rho )`$. In fact, to calculate $`Er(\sigma )`$ is to find the state $`\rho ^{}`$. Main result: For bipartite quantum state $`\sigma `$ $`=`$ $`{\displaystyle \underset{n_1n_2}{}}a_{n_1n_2}\|\varphi _{n_1}\psi _{n_1}\varphi _{n_2}\psi _{n_2}\|`$ (2) $`=`$ $`{\displaystyle \underset{n_1n_2}{}}a_{n_1n_2}\|\varphi _{n_1}\varphi _{n_2}\|\|\psi _{n_1}\psi _{n_2}\|`$ (3) the relative entropy of entanglement is given by $$Er(\sigma )=\underset{n}{}a_{nn}\mathrm{ln}a_{nn}S(\sigma )$$ (4) and the disentangled state that minimizes the relative entropy is $$\rho ^{}=\underset{n}{}a_{nn}\|\varphi _n\psi _n\varphi _n\psi _n\|$$ (5) Here, $`\|\varphi _n`$ ($`\|\psi _n`$) is a set of orthogonal normalized states of system A (B); $`S(\sigma )`$ is the von Neumann entropy defined by $$S(\sigma )tr\{\sigma \mathrm{ln}\sigma \}$$ (6) If $`\sigma `$ is a pure state, there is $`a_{n_1n_2}=\sqrt{p_{n_1}p_{n_2}}`$, our result reduces to Vedral and Plenio’s result. Proof. The proof is similar to the proof of theorem 3 in . Since we already have a guess for $`\rho ^{}`$, we show that the gradient $`\frac{d}{dx}S(\sigma \|\|(1x)\rho ^{}+x\rho )\|_{x=0}`$ for any $`\rho D`$ is non-negative. On the other hand, if $`\rho ^{}`$ was not a minimum, the above gradient would be strictly negative, which is a contradiction. Thus the result follows. Let $`f(x,\rho )S(\sigma \|\|(1x)\rho ^{}+x\rho )`$, using the identity $`\mathrm{ln}A=_0^{\mathrm{}}[(At1)/(A+t)]𝑑t/(1+t^2)`$, we can get $$\frac{f}{x}(0,\rho )=1_0^{\mathrm{}}tr[(\rho ^{}+t)^1\sigma (\rho ^{}+t)^1\rho ]𝑑t$$ (7) Since $`\rho ^{}=_na_{nn}\|\varphi _n\psi _n\varphi _n\psi _n\|`$, it is not difficult to get $`(\rho ^{}+t)^1\sigma (\rho ^{}+t)^1`$ (8) $`={\displaystyle \underset{nn^{^{}}}{}}(a_{nn}+t)^1a_{nn^{^{}}}(a_{n^{^{}}n^{^{}}}+t)^1\|\varphi _n\psi _n\varphi _n^{^{}}\psi _n^{^{}}\|`$ (9) Set $`g(n,n^{^{}})a_{nn^{^{}}}_0^{\mathrm{}}(a_{nn}+t)^1(a_{n^{^{}}n^{^{}}}+t)^1𝑑t`$, obviously we have that $`g(n,n)=1`$, and for $`nn^{^{}}`$ $$g(n,n^{^{}})=a_{nn^{^{}}}\frac{\mathrm{ln}a_{nn}\mathrm{ln}a_{n^{^{}}n^{^{}}}}{a_{nn}a_{n^{^{}}n^{^{}}}}$$ (10) We now show that $`\|g(n,n^{^{}})\|1`$. Ref. has proved that $$0\sqrt{a_{nn}a_{n^{^{}}n^{^{}}}}\frac{\mathrm{ln}a_{nn}\mathrm{ln}a_{n^{^{}}n^{^{}}}}{a_{nn}a_{n^{^{}}n^{^{}}}}1$$ (11) so we only need to show that $`\|a_{nn^{^{}}}\|\sqrt{a_{nn}a_{n^{^{}}n^{^{}}}}`$. Let $`\|\mathrm{\Psi }=a\|\varphi _n\psi _n+b\|\varphi _n^{^{}}\psi _n^{^{}}`$, $`a`$ and $`b`$ are arbitrary complex numbers. Since $`\sigma `$ is a quantum state, it follows that $$\mathrm{\Psi }\left\|\sigma \right\|\mathrm{\Psi }0$$ (12) for arbitrary pair of $`a`$ and $`b`$, this requires $$a_{nn}a_{n^{^{}}n^{^{}}}a_{nn^{^{}}}a_{n^{^{}}n}=a_{nn}a_{n^{^{}}n^{^{}}}\|a_{nn^{^{}}}\|^20$$ (13) Therefore $$\|g(n,n^{^{}})\|1$$ (14) The following steps are just the same as those in ref. , which are written here for the completeness of this proof. Let $`\rho \|\alpha \alpha \|\|\beta \beta \|`$ where $`\|\alpha =_na_n\|\varphi _n`$ and $`\|\beta =_nb_n\|\psi _n`$ are normalized vectors, it is not difficult to show that $$\frac{f}{x}(0,\rho )1=\underset{n_1n_2}{}g(n_1,n_2)a_{n_2}b_{n_2}a_{n_1}^{}b_{n_1}^{}$$ (15) therefore $`\left\|{\displaystyle \frac{f}{x}}(0,\rho )1\right\|`$ (16) $``$ $`{\displaystyle \underset{n_1n_2}{}}\|g(n_1,n_2)\|\|a_{n_2}\|\|b_{n_2}\|\|a_{n_1}^{}\|\|b_{n_1}^{}\|`$ (17) $``$ $`{\displaystyle \underset{n_1n_2}{}}\|a_{n_2}\|\|b_{n_2}\|\|a_{n_1}^{}\|\|b_{n_1}^{}\|`$ (18) $`=`$ $`({\displaystyle \underset{n}{}}\|a_n\|\|b_n\|)^2`$ (19) $``$ $`{\displaystyle \underset{n}{}}\|a_n\|^2{\displaystyle \underset{n}{}}\|b_n\|^2`$ (20) $`=`$ $`1`$ (21) Then it follows that $`\frac{f}{x}(0,\|\alpha \beta \alpha \beta \|)0`$. Since any disentangled state $`\rho D`$ can be written in the form $`\rho =_ir_i\|\alpha ^i\beta ^i\alpha ^i\beta ^i\|`$, we have that $$\frac{f}{x}(0,\rho )=\underset{i}{}r_i\frac{f}{x}(0,\|\alpha ^i\beta ^i\alpha ^i\beta ^i\|)0$$ (22) Now we show that $`S(\sigma \|\|\rho )S(\sigma \|\|\rho ^{})`$ for all $`\rho D`$. Suppose that $`S(\sigma \|\|\rho )<S(\sigma \|\|\rho ^{})`$ for some $`\rho D`$, then, for $`0<x1`$, $`f(x,\rho )`$ $`=`$ $`S(\sigma \|\|(1x)\rho ^{}+x\rho )`$ (23) $``$ $`(1x)S(\sigma \|\|\rho ^{})+xS(\sigma \|\|\rho )`$ (24) $`=`$ $`(1x)f(0,\rho )+xf(1,\rho )`$ (25) which implies $$\frac{f(x,\rho )f(0,\rho )}{x}f(1,\rho )f(0,\rho )<0$$ (26) This contradicts the fact that $`\frac{f}{x}(0,\rho )0`$ in the limit of small $`x`$. So, for any $`\rho D`$, this is $`S(\sigma \|\|\rho )S(\sigma \|\|\rho ^{})`$, i.e., the state $`\rho ^{}=_na_{nn}\|\varphi _n\psi _n\varphi _n\psi _n\|`$ minimizes the relative entropy $`S(\sigma \|\|\rho )`$ over $`\rho D`$. Then it follows that $$Er(\sigma )=tr\{\sigma (\mathrm{ln}\sigma \mathrm{ln}\rho ^{})\}=\underset{n}{}a_{nn}\mathrm{ln}a_{nn}S(\sigma )$$ (27) This completes the proof. It should be pointed out that the above result can also follow directly from Rains’ theorem 9 in ref. . The Vedral and Plenio’s theorem has been extended to a more general case, the result is, nevertheless, very useful for calculating the relative entropy of entanglement for the mixed states $`\sigma =_{n_1n_2}a_{n_1n_2}\|\varphi _{n_1}\psi _{n_1}\varphi _{n_2}\psi _{n_2}\|`$. For the simplest case of two qubits, the relative entropy of entanglement $`Er`$ of the state $`\sigma `$ $`=`$ $`x\|0000\left\|+(1x)\right\|1111\|`$ (29) $`+\alpha \|0011\left\|+\alpha ^{}\right\|1100\|`$ is given by $`Er(\sigma )`$ $`=`$ $`x\mathrm{ln}x(1x)\mathrm{ln}(1x)`$ (31) $`+\lambda \mathrm{ln}\lambda +(1\lambda )\mathrm{ln}(1\lambda )`$ where $`\lambda =\frac{1}{2}\{1+\sqrt{14[x(1x)\|\alpha \|^2]}\}`$. Other applications of this result can be found in ref. . We thank Prof. E. M. Rains for pointing out an important reference.
warning/0004/hep-th0004024.html	ar5iv	text	# 1 Introduction ## 1 Introduction Noncommutative field theories (NCFT) have been studied for some time in the search for a new class of field thories . Four-dimensional noncommutative gauge theories arizes naturally in open string theories incorporating D-branes , and they stimulated renewed interrest in NCFT. Perturbation analyses of a few kinds of NCFT have recentry been made , but full understanding of their quantum theory is still lacking. In this letter we aim at constructing noncommutative extension of two-dimensional integrable field theories. Integrable field theories, e.g, sine-Gordon model, non-linear $`\sigma `$ model, Wess-Zumino-Witten model, have many important physical properties but allows computation of exact Green functions quantum mechanically. If we succeed in their noncommutative extension, the resulting noncommutative integrable field theory will provide us a concrete example of quantum NCFT. They will also give us a useful guide for understanding quantum NCFT in general. Integrable field theories posess infinitely many conservation laws (infinite-dimensional symmetry) which generate an infinite-dimensional Lie algebra or its deformation. It is vitally important to understand how Hamiltonian and other conservation laws are constructed in NCFT. In this letter we consider the Wess-Zumino-Witten (WZW) model and construct its noncommutative extension. The noncommutative WZW model is found to have the same infinite-dimensional symmetry as the ordinary (commutative) WZW model at their respective critical points. We are interested in the perturbative property of this theory and evaluate its $`\beta `$-function. We assume that there exist some regularization to deal with infrared divergences. In the course of writing this letter there appeared three papers , in which the WZW model in noncommutative space is discussed. Dabrowski et al. study the noncommutative nonlinear $`\sigma `$ models. The infinite-dimensional symmetry discussed in sect.3 of our letter has also been derived by them. They have pointed out some larger symmetry. In and the authers have derived the WZW action by integrating fermion fields in the noncommutative extension of a gauge theory. ## 2 Noncommutative Extension of the WZW Model We take the $`SU(N)`$ WZW model, and define its noncommutative extension (noncommutative WZW model) by the action obtained from the WZW model by replacing product of fields by $``$-product, $$I(g)=\frac{1}{2\lambda ^2}_\mathrm{\Sigma }d^2x\mathrm{Tr}\left(_\mu gg^1^\mu gg^1\right)_{}+\frac{k}{12\pi }_B\mathrm{Tr}\left(d\stackrel{~}{g}\stackrel{~}{g}^1\right)_{}^3.$$ (1) Here $`\mathrm{\Sigma }`$ is the boundary of a three-dimensional manifold $`B`$, and $`\stackrel{~}{g}`$ is an extension of $`g`$ such that $$\mathrm{\Sigma }=B,\stackrel{~}{g}(y)=g(y)\mathrm{for}yB.$$ (2) We consider the flat manifold $`\mathrm{\Sigma }`$ parametrized by ($`x^0`$, $`x^1`$), and take the metric $`g^{\mu \nu }=\mathrm{diag}(1,1)`$. The $``$-product $`(fg)_{}`$ is defined by $$\left(fg\right)_{}(x)\stackrel{def}{=}fg(x)=\mathrm{exp}i\xi \theta ^{\mu \nu }\frac{}{u^\mu }\frac{}{v^\nu }f(u)g(v)\|_{u=v=x},$$ (3) where $`\theta ^{\mu \nu }`$ is a 2nd rank antisymmetric tensor and is of the form $$\left\{\theta ^{\mu \nu }\right\}=\left(\begin{array}{ccc}0& 1& 0\\ 1& 0& 0\\ 0& 0& 0\end{array}\right).$$ (4) The $``$-product satisfies the associativity $`(fgh)_{}\stackrel{def}{=}(f(gh))=((fg)h)`$. The group-valued field $`\stackrel{~}{g}(y)`$ is defined as $$\stackrel{~}{g}=\mathrm{exp}_{}(i\phi ^at_a)=\left(1+i\phi ^at_a+\frac{1}{2!}(i\phi ^at_a)^2+\mathrm{}\right)_{},$$ (5) where $`t^a`$ are Hermitian matrices with the normalization $`\mathrm{Tr}(t_at_b)=\delta _{ab}`$. In order that the product $`\stackrel{~}{g}_1(y)\stackrel{~}{g}_2(y)=\mathrm{exp}_{}(i\phi _1^at^a)\mathrm{exp}_{}(i\phi _2^at^a)`$ makes sense, $`t^a`$ have to be generators of $`U(N)`$ in place of $`SU(N)`$. Hence the action (1) defines the noncommutative $`U(N)`$ WZW model. ## 3 Equation of motion We will investigate the ultraviolet property of the model (1). This can be made in perturbation and using the background field method. To this end we begin by parametrizing the field $`g`$ around the classical background $`g_c`$. We set $`g=g_qg_c`$ and express $`g_q`$ in terms of quantum fluctuation $`\pi ^a`$ as $$g_q(x)=\mathrm{exp}_{}\left(i\lambda \pi (x)\right),\pi =\pi ^at^a.$$ (6) The action is then expanded in powers of $`\pi ^a`$ as $$I(g)=I(g_c)+I_1+I_2+\mathrm{}$$ (7) where $`I_n`$ represents the $`n`$-th order term in $`\pi ^a`$. The first two are $$I_1=\frac{i}{\lambda }_\mathrm{\Sigma }d^2x\left(g^{\mu \nu }\frac{\lambda ^2k}{4\pi }ϵ^{\mu \nu }\right)\mathrm{Tr}\left\{_\mu g_cg_c^1_\nu \pi \right\}_{},$$ (8) $$I_2=_\mathrm{\Sigma }d^2x\mathrm{Tr}\left\{\frac{1}{2}(_\mu \pi )^2+\frac{1}{2}\left(g^{\mu \nu }\frac{\lambda ^2k}{4\pi }ϵ^{\mu \nu }\right)_\mu g_cg_c^1[_\nu \pi ,\pi ]\right\}_{}.$$ (9) We get the equation of motion from (8) $$\left(g^{\mu \nu }\frac{\lambda ^2k}{4\pi }ϵ^{\mu \nu }\right)\mathrm{Tr}_\nu \left(_\mu g_cg_c^1\right)_{}=0.$$ (10) For the specific value of the coupling constant $`\lambda `$, $$\lambda ^2=\frac{4\pi }{k},$$ (11) eq.(10) takes a simpler form. Namely it is reduced to the left-handed (or equivalentry right-handed) current conservation law, $$_+J_{}=0,J_{}=_{}gg^1$$ (12) (or equivalentry, $$_{}J_+=0,J_+=g^1_+g).$$ (13) Here $`_\pm `$ refers to the light-cone coodinates $`x^\pm =2^{\frac{1}{2}}(x^0\pm x^1`$). At the point (11), the action (1) can be shown to be invariant under the transformations of $`g`$ by left and right multiplication, $`g(x)`$ $``$ $`g^{}(x)=g_{}(x^{})g(x),`$ $`g(x)`$ $``$ $`g^{}(x)=g(x)g_+(x^+),`$ (14) where $`g_\pm `$ represents any element of $`U(N)`$ and depends only on $`x^\pm `$, respectively. The fact that $`g_\pm `$ is a function of $`x^\pm `$ means that the symmetry group is infinite-dimensional. It is a loop group if we choose $`x^\pm `$ to be a peoriodic coodinate. We further note the products of (an arbitrary number of) the chiral current $`J_\pm `$ are also conserved. In particular, the singlet part of quadratic terms $$T_\pm (x^\pm )=:J_\pm ^a(x^\pm )J_\pm ^a(x^\pm ):$$ (15) is reminiscent of the Sugawara construction of the Virasoro algebra. This possibility needs a further study. ## 4 One-loop $`\beta `$-function The standard background field method can be applied to the loop computation in NCFT. In this method $`I_n(n2)`$ can be regarded as the interaction term for the $`\pi `$ field. To evaluate one-loop terms, only $`I_2`$ is needed. We introduce the momentum representation of $`I_2`$. The interaction part (second term) of (9) is then, $`I_{\mathrm{int}}(g_c,\pi )`$ $`=`$ $`{\displaystyle \frac{1}{4}}{\displaystyle \frac{f^{\mu \nu }}{(2\pi )^4}}{\displaystyle d^2p_1d^2p_2d^2p_3\delta ^2(p_1+p_2+p_3)G_\mu ^c(p_1)\pi ^a(p_2)\pi ^b(p_3)}`$ (16) $`\times i(p_2p_3)_\nu \left(f_{cab}\mathrm{cos}(p_2p_3)id_{cab}\mathrm{sin}(p_2p_3)\right),`$ where we have written $`f^{\mu \nu }`$ $`=`$ $`g^{\mu \nu }{\displaystyle \frac{\lambda ^2k}{4\pi }}ϵ^{\mu \nu }`$ $`pq`$ $`=`$ $`\xi \theta ^{\mu \nu }p_\mu q_\nu .`$ (17) $`G_\mu ^a=\left(_\mu g_cg_c^1\right)^a`$ is the background. The group factors appearing in eq. (16) are $$f_{abc}=\mathrm{Tr}\left(t_a[t_b,t_c]\right),d_{abc}=\mathrm{Tr}\left(t_a\{t_b,t_c\}\right).$$ (18) There is only one one-loop diagram that contributes to the ultraviolet (UV) divergence of the effective action $`\mathrm{\Gamma }(g_c)`$ (Fig. 1). The contribution of this diagram to $`\mathrm{\Gamma }(g_c)`$ is $`\mathrm{\Gamma }_{\mathrm{Fig}.1}(g_c)`$ $`=`$ $`{\displaystyle \frac{1}{i}}{\displaystyle \frac{d^2p}{(2\pi )^2}\frac{1}{64}f^{\rho \lambda }f^{\mu \nu }G_\rho ^a(p)G_\mu ^b(p)\frac{d^2k}{(2\pi )^2}\frac{(p+2k)_\lambda (p+2k)_\nu }{(k^2\mu ^2)\{(p+k)^2\mu ^2\}}}`$ (19) $`\times \left\{(F+D)_{ab}+\mathrm{cos}(2pk)(FD)_{ab}i\mathrm{sin}(2pk)C_{ab}\right\},`$ where $`F_{ab}=f_{acd}f_{bdc},D_{ab}=d_{acd}d_{bdc},`$ $`C_{ab}=f_{acd}d_{bdc}f_{bcd}d_{adc}.`$ (20) Only the first term in the second line of (19) contributes to the UV divergence. After using the identity for $`U(N)`$ group $$F_{ab}+D_{ab}=4N\delta _{ab},$$ (21) we obtain the UV divergent part of (19) as $$\mathrm{\Gamma }_{\mathrm{div}}^{(1\mathrm{loop})}(g_c)=\frac{N\lambda ^2}{32\pi }\left\{1\left(\frac{\lambda ^2k}{4\pi }\right)^2\right\}\mathrm{ln}\frac{\mu ^2}{\mathrm{\Lambda }^2}S(g_c),$$ (22) where $$S(g_c)=\frac{1}{2\lambda ^2}_\mathrm{\Sigma }d^2x\mathrm{Tr}\left(_\mu g_cg_c^1^\mu g_cg_c^1\right)_{}$$ (23) is the first term in the action (1) for the background field. $`\mu `$ is the renormalization point and $`\mathrm{\Lambda }`$ is the ultraviolet cutoff. We take (22)(with opposite sign) as the one-loop counterterm, then we get the $`\beta `$-function of $`U(N)`$ noncommutative WZW model as $$\beta _{N.C.}=\mu \frac{d}{d\mu }\lambda =\frac{N\lambda ^3}{32\pi }\left\{1\left(\frac{\lambda ^2k}{4\pi }\right)^2\right\}.$$ (24) As pointed out in , the second term in R.H.S of (24) has meaning only for large $`k`$. The above result is the same as that of the ordinary $`SU(N)`$ WZW model . To see this, set $`\xi 0`$ in (19). We then obtain the same result as (24), and hence $$\beta =\frac{N\lambda ^3}{32\pi }\left\{1\left(\frac{\lambda ^2k}{4\pi }\right)^2\right\}.$$ (25) We have used the identity for $`SU(N)`$ group $$F_{ab}=2N\delta _{ab}$$ (26) instead of (21). ## 5 Acnowledgements We would like to thank N. Ishibashi, M. Hayakawa and Y. Yoshida for discussions and helpful suggestions.
warning/0004/hep-th0004076.html	ar5iv	text	# AEI-2000-022hep-th/0004076 Yang-Mills Integrals for Orthogonal, Symplectic and Exceptional Groups ## I Introduction Recent attempts to describe $`D`$-branes through effective actions have revealed the existence of a new class of gauge-invariant matrix models. These models are related to ($`1`$)-branes. They differ from the classic systems of random matrices, which have been extensively studied ever since Wigner’s and Dyson’s work in the 1950’s. The models consist in matrix integrals of $`D`$ non-linearly coupled matrices. They are obtained by complete dimensional reduction of $`D`$-dimensional Euclidean continuum (susy) Yang-Mills theory to zero dimensions, and we term them quite generally Yang-Mills integrals. The supersymmetric $`D=4,6,10`$ integrals with SU$`(N)`$ symmetry have already found several important applications. They are relevant to the calculation of the Witten index of quantum mechanical gauge theory ,, and to multi-instanton calculations , of four-dimensional SU$`(\mathrm{})`$ susy conformal gauge theory. The $`D=10`$ integrals are furthermore the crucial ingredient in the so-called IKKT model , which possibly provides a non-perturbative definition of IIB superstring theory. Finally, it remains to be seen whether Yang-Mills integrals contain information on the full, unreduced field theory through the Eguchi-Kawai mechanism as the size of the matrices gets large. Some very interesting recent considerations along these lines can be found in . Yang-Mills integrals are ordinary, not functional integrals. Despite this tremendous simplification, no systematic analytic tools for their investigation are known to date. We have developed ,,, accurate and reliable Monte Carlo methods which allow to study the new matrix models as long as the dimension of the gauge group is not too large. We have found, e.g., that supersymmetry is generically not necessary for the existence of the integrals . We also computed their asymptotic eigenvalue distributions, which we found to qualitatively differ between the susy and bosonic case as the size of the gauge group gets large . For related, complementary studies see . To date all existing studies have focused on the case where the gauge group is SU$`(N)`$. In the present paper we generalize to the cases of all other semi-simple compact Lie groups of rank $`r3`$. These are, (besides the already known cases SU$`(2)`$, SU$`(3)`$, SU$`(4)`$) the groups SO$`(3)`$, SO$`(4)`$, SO$`(5)`$, SO$`(6)`$, SO$`(7)`$, Sp$`(2)`$, Sp$`(4)`$, Sp$`(6)`$, and G<sub>2</sub>, for which we compute the susy $`D=4`$ partition functions. We were motivated in part by a recent paper of Kac and Smilga which presented conjectures about the values of the bulk part of the Witten index (and therefore for the corresponding integrals). Intriguingly, our results are at variance with their predictions in most cases, indicating that the index calculations for these groups are even more subtle than the corresponding considerations for SU$`(N)`$, where the approach of agrees with the known values. Moore, Nekrasov and Shatashvili recently employed sophisticated deformation techniques to evaluate the SU$`(N)`$ susy bulk index for all $`N`$ and $`D=4,6,10`$. The method apparently leads to the correct result for all SU$`(N)`$ , . Below, we adapt the technique to the more general groups, and we find again excellent agreement with the Monte Carlo calculation. This further indicates that the deformation method is indeed reliable. ## II Yang-Mills integrals for semi-simple compact gauge groups For a general semi-simple compact Lie group $`G`$ we define supersymmetric or bosonic Yang-Mills integrals as $$𝒵_{D,G}^𝒩:=\underset{A=1}{\overset{\mathrm{dim}(G)}{}}\left(\underset{\mu =1}{\overset{D}{}}\frac{dX_\mu ^A}{\sqrt{2\pi }}\right)\left(\underset{\alpha =1}{\overset{𝒩}{}}d\mathrm{\Psi }_\alpha ^A\right)\mathrm{exp}\left[\frac{1}{4g^2}\mathrm{Tr}[X_\mu ,X_\nu ][X_\mu ,X_\nu ]+\frac{1}{2g^2}\mathrm{Tr}\mathrm{\Psi }_\alpha [\mathrm{\Gamma }_{\alpha \beta }^\mu X_\mu ,\mathrm{\Psi }_\beta ]\right],$$ (1) where dim$`(G)`$ is the dimension of the Lie group and the $`D`$ bosonic matrices $`X_\mu =X_\mu ^AT^A`$ and the $`𝒩`$ fermionic matrices $`\mathrm{\Psi }_\alpha =\mathrm{\Psi }_\alpha ^AT_A`$ are anti-hermitean and take values in the fundamental representation of the Lie algebra Lie$`(G)`$, whose generators we denote by $`T^A`$. The integral eq.(1) depends on the gauge coupling constant in a trivial fashion, as we can immediately scale out $`g`$. Nevertheless, there is a natural convention for fixing $`g`$: For an orthogonal set of generators we should pick $`g`$ according to their normalization: $`\mathrm{Tr}T^AT^B=g^2\delta ^{AB}`$. This convention is imposed by the index calculations of the next section. In the bosonic case we simply drop the fermionic variables. The supersymmetric integrals can formally be defined in $`D=3,4,6,10`$, which corresponds to $`𝒩=2,4,8,16`$ real supersymmetries. We are not aware of a mathematically rigorous investigation of their convergence properties. However, our numerical studies indicate that they are absolutely convergent in $`D=4,6,10`$ (but not in $`D=3`$) for all semi-simple compact gauge groups. The convergence properties of the bosonic ($`𝒩=0`$) integrals are discussed in section 6. The variables $`\mathrm{\Psi }_\alpha ^A`$ are real Grassmann-valued and can be integrated out, leading to a bosonic integral with very special measure: $$𝒵_{D,G}^𝒩=\underset{A=1}{\overset{\mathrm{dim}(G)}{}}\underset{\mu =1}{\overset{D}{}}\frac{dX_\mu ^A}{\sqrt{2\pi }}\mathrm{exp}\left[\frac{1}{4g^2}\mathrm{Tr}[X_\mu ,X_\nu ][X_\mu ,X_\nu ]\right]𝒫_{D,G}(X).$$ (2) $`𝒫_{D,G}(X)`$ is a homogeneous Pfaffian polynomial of degree $`\frac{1}{2}𝒩\mathrm{dim}(G)`$ given by $$𝒫_{D,G}=\mathrm{Pf}_{D,G}\mathrm{with}\left(_{D,G}\right)_{\alpha \beta }^{AB}=if^{ABC}\mathrm{\Gamma }_{\alpha \beta }^\mu X_\mu ^C,$$ (3) where the structure constants <sup>§</sup><sup>§</sup>§Note that in we used hermitean generators but defined the structure constants through $`[T^A,T^B]=if^{ABC}T^C`$, so eq.(3) remains valid. are defined through the real Lie algebra $`[T^A,T^B]=f^{ABC}T^C`$. Explicit expressions for the Gamma matrices $`\mathrm{\Gamma }_{\alpha \beta }^\mu `$ and further details on $`𝒫_{D,G}(X)`$ may be found in . ## III Group volumes and bulk indices It is well known that the susy Yang-Mills integrals eq.(2) naturally appear when one computes the Witten index of quantum mechanical gauge theory (i.e. the reduction of the field theory to one, as opposed to zero, dimension) by the heat kernel method. For the details of the method we refer to ,. Specifically, the integrals are related to the bulk part ind$`{}_{0}{}^{D}(G)`$ of the index as $$\mathrm{ind}_0^D(G)=\mathrm{lim}_{\beta 0}\mathrm{Tr}(1)^Fe^{\beta H}=\frac{1}{_G}𝒵_{D,G}^𝒩.$$ (4) The total Witten index “ind$`{}_{}{}^{D}(G)`$” is then the sum of this bulk part and a boundary contribution “ind$`{}_{1}{}^{D}(G)`$”: $`\mathrm{ind}^D(G)=\mathrm{ind}_0^D(G)+\mathrm{ind}_1^D(G)`$. The constant $`_G`$ relating the bulk index ind$`{}_{0}{}^{}{}_{}{}^{D}`$ and the Yang-Mills integral is independent of $`D`$ and can be interpreted as $$_G=\frac{1}{(2\pi )^{\frac{1}{2}\mathrm{dim}(G)}}\mathrm{Volume}\left[\frac{G}{Z_G}\right],$$ (5) i.e. essentially the volume of the true gauge group, which turns out to be the quotient group $`G/Z_G`$, with $`Z_G`$ the center group of $`G`$. In practice, great care has to be taken in using the relation eq.(5), as the volume depends on the choice of the local metric on the group manifold. For the present purposes we simply adapted our method for computing $`_{\mathrm{SU}(N)}`$ (see ) to the relevant gauge groups. An invariant average over the group allows to project onto gauge invariant states and to derive eq.(4) from the quantum mechanical path integral. In the ultralocal limit, the quantum mechanics of $`D1`$ matrices turns into an integral over $`D`$ matrices. Then, this integration is over the anti-hermitean generators of the group $$𝒟U\frac{1}{_G}\underset{A=1}{\overset{\mathrm{dim}(G)}{}}\frac{dX_D^A}{\sqrt{2\pi }}.$$ (6) The normalized Haar measure $`𝒟U`$ on the group elements $`UG`$ simplifies significantly if we restrict attention to the Cartan subgroup of $`G`$. A beautiful result of Weyl allows to explicitly write down the restricted measure. If we parametrize the Cartan torus $`T`$ by angles $`\pi \theta _1\pi `$, $`\mathrm{}`$, $`\pi \theta _r\pi `$, where $`r=`$rank$`(G)`$, the measure reads $$𝒟U𝒟T=\frac{1}{\|W_G\|}\left(\underset{i=1}{\overset{r}{}}\frac{d\theta _i}{2\pi }\right)\|\mathrm{\Delta }_G\|^2,$$ (7) where $`\|W_G\|`$ is the order of the Weyl group $`W_G`$ of $`G`$, and $`\|\mathrm{\Delta }_G\|^2`$ the squared modulus of the Weyl denominator: $$\mathrm{\Delta }_G=\underset{\alpha >0}{}\left[e^{\frac{i}{2}(\theta ,\alpha )}e^{\frac{i}{2}(\theta ,\alpha )}\right].$$ (8) Here the product is over the set of positive roots of the Lie algebra Lie$`(G)`$. In the vicinity of the identity in $`G`$ the angles $`\theta _i`$ are small and we can approximate the measure eq.(7) by $$\frac{1}{\|W_G\|}\left(\underset{i=1}{\overset{r}{}}\frac{d\theta _i}{2\pi }\right)\underset{\alpha >0}{}\left[\frac{1}{2}(\theta ,\alpha )\frac{1}{2}(\theta ,\alpha )\right]^2.$$ (9) Now restricting the flat measure on Lie$`(G)`$ on the right hand side of eq.(6) to the Cartan modes $`\theta _i`$ we get $$\underset{A=1}{\overset{\mathrm{dim}(G)}{}}\frac{dX_D^A}{\sqrt{2\pi }}\frac{_G}{(2\pi )^r}\frac{\|Z_G\|}{\|W_G\|}\left(\underset{i=1}{\overset{r}{}}d\theta _i\right)\underset{\alpha >0}{}\left[\frac{1}{2}(\theta ,\alpha )\frac{1}{2}(\theta ,\alpha )\right]^2.$$ (10) An important subtlety is that we needed to multiply the measure eq.(10) by an additional factor $`\|Z_G\|`$ of the order of the center group $`Z_G`$ of $`G`$, as the averaging over the group manifold localizes on $`\|Z_G\|`$ points. Finally, the constant $`_G`$ in eq.(10) is fixed by noting that the flat measure on Lie$`(G)`$ is normalized with respect to Gaussian integration $$\left(\underset{A=1}{\overset{\mathrm{dim}(G)}{}}\frac{dX_D^A}{\sqrt{2\pi }}\right)\mathrm{exp}\left[\frac{1}{2}\underset{A}{}(X_D^A)^2\right]=1.$$ (11) The Gaussian integration of the right hand side of eq.(10) leads to Selberg-type integrals, see e.g. . Explicit details on how to implement the above procedure for the groups under study can be found in the appendices. With our conventions for the normalization of the generators (SO$`(N)`$: $`\mathrm{Tr}T^AT^B=2\delta ^{AB}`$ and Sp$`(2N)`$, G<sub>2</sub>: $`\mathrm{Tr}T^AT^B=\frac{1}{2}\delta ^{AB}`$) one finds $$_{\mathrm{SO}(N)}=\frac{1}{2C_N}\frac{\pi ^{\frac{N}{2}}}{2^{\frac{N(N5)}{4}}_{j=1}^N\mathrm{\Gamma }(j/2)},$$ (12) where $`C_{2N}=2`$ and $`C_{2N+1}=1`$, as well as $$_{\mathrm{Sp}(2N)}=\frac{1}{2}\frac{2^{2N^2+\frac{N}{2}}\pi ^{\frac{N}{2}}}{_{j=1}^N\mathrm{\Gamma }(2j)}$$ (13) and finally $$_{\mathrm{G}_2}=\frac{36864\sqrt{3}\pi }{5}.$$ (14) ## IV Deformation method In it was suggested that the original susy Yang-Mills integrals eq.(1) may be vastly simplified by a deformation technique. It consists in adding a number of terms to the action which break the number of supersymmetries from $`𝒩=2,4,8,16`$ to $`𝒩=1`$. Keeping one of the supersymmetries means that the partition function is “protected” and should not change under the deformation. This gives the correct resultFor unclear reasons it fails for $`D=3`$. for SU$`(N)`$ and $`D=4,6,10`$. The final outcome is a much simpler integral involving only a single Lie-algebra valued matrix. The remaining integral is still invariant under the gauge group, and one can therefore pass from the full algebra to the Cartan subalgebra degrees of freedom. This was derived in in detail for SU$`(N)`$ but should carry over immediately to other gauge groups. For $`D=4`$ ($`𝒩=4`$) one finds, in the notation of the previous section (here the product is over all roots $`\alpha `$ of Lie$`(G)`$) $$\mathrm{ind}_0^{D=4}(G)=\frac{\|Z_G\|}{\|W_G\|}\frac{1}{E^r}\left(\underset{i=1}{\overset{r}{}}\frac{dx_i}{2\pi i}\right)\underset{\alpha }{}\frac{\left[\frac{1}{2}(x,\alpha )\frac{1}{2}(x,\alpha )\right]}{\left[\frac{1}{2}(x,\alpha )\frac{1}{2}(x,\alpha )E\right]}.$$ (15) This $`r`$-dimensional integral ($`r=`$rank$`(G)`$) is divergent. There are divergences due to the poles of the denominator of eq.(15), as well as at infinity, where the integrand tends to one. These divergences are present since the method starts from the partition sums eq.(1) with Minkowski signature, i.e. with the Wick-rotated versions of our integrals. The Minkowski integrals are divergent without a prescription. The poles are regulated by giving an imaginary part to the parameter $`E`$. The singularity at infinity is regulated by interpreting the integrals as contour integralsThis interpretation furthermore necessitates the inclusion of the factors of $`i`$ in the measure of eq.(15) which would not be present in an ordinary integral over a real Lie algebra.. It would be interesting to complete the arguments by demonstrating that the Wick-rotation leads to precisely these prescriptions. Very encouraging signs for the consistency of this method are that the final result neither depends on the location of the parameter $`E`$ nor on whether the contours are closed in the upper or the lower half plane (it is important though that all $`r`$ contours are closed in the same way). For $`D=6,10`$ expressions very similar to eq.(15) can be found in . We now present the explicit form of the contour integrals eq.(15) for the groups studied in the present work (see appendices for details) $`\mathrm{ind}_0^{D=4}(\mathrm{SO}(2N+1))={\displaystyle \frac{1}{2^NN!}}{\displaystyle \frac{1}{E^N}}{\displaystyle }{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \frac{dx_i}{2\pi i}}{\displaystyle \underset{i<j}{\overset{N}{}}}{\displaystyle \frac{(x_i^2x_j^2)^2}{\left[(x_ix_j)^2E^2\right]\left[(x_i+x_j)^2E^2\right]}}\times `$ $$\times \underset{i=1}{\overset{N}{}}\frac{x_i^2}{x_i^2E^2}$$ (16) $$\mathrm{ind}_0^{D=4}(\mathrm{SO}(2N))=\frac{2}{2^{N1}N!}\frac{1}{E^N}\underset{i=1}{\overset{N}{}}\frac{dx_i}{2\pi i}\underset{i<j}{\overset{N}{}}\frac{(x_i^2x_j^2)^2}{\left[(x_ix_j)^2E^2\right]\left[(x_i+x_j)^2E^2\right]}$$ (17) $`\mathrm{ind}_0^{D=4}(\mathrm{Sp}(2N))={\displaystyle \frac{2}{2^NN!}}{\displaystyle \frac{1}{E^N}}{\displaystyle }{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \frac{dx_i}{2\pi i}}{\displaystyle \underset{i<j}{\overset{N}{}}}{\displaystyle \frac{(x_i^2x_j^2)^2}{\left[(x_ix_j)^2E^2\right]\left[(x_i+x_j)^2E^2\right]}}\times `$ $$\times \underset{i=1}{\overset{N}{}}\frac{x_i^2}{x_i^2(\frac{E}{2})^2}$$ (18) $`\mathrm{ind}_0^{D=4}(\mathrm{G}_2)={\displaystyle \frac{1}{12}}{\displaystyle \frac{1}{E^2}}{\displaystyle \frac{dx_1}{2\pi i}\frac{dx_2}{2\pi i}\frac{(x_1x_2)^2(x_1+x_2)^2x_1^2x_2^2}{\left[(x_1x_2)^2E^2\right]\left[(x_1+x_2)^2E^2\right]\left[x_1^2E^2\right]\left[x_2^2E^2\right]}\times }`$ $$\times \frac{(2x_1+x_2)^2(x_1+2x_2)^2}{\left[(2x_1+x_2)^2E^2\right]\left[(x_1+2x_2)^2E^2\right]}$$ (19) They are easily evaluated for low rank, and we present the results in table 1. We highlighted the cases which were not already indirectly known due to the standard low-rank isomorphisms $`\mathrm{so}(3)=\mathrm{sp}(2)=\mathrm{su}(2)`$, $`\mathrm{so}(4)=\mathrm{su}(2)\mathrm{su}(2)`$, $`\mathrm{so}(6)=\mathrm{su}(4)`$. Note, however, that these identities, as well as the final semi-simple Lie algebra isomorphism $`\mathrm{so}(5)=\mathrm{sp}(4)`$, constitute non-trivial consistency checks on the expression eq.(15), as the precise form of the corresponding contour integrals is different in all these cases. It would be interesting to compute eqs. (16),(17),(18),(19) for arbitrary rank, as has been done in the case of SU$`(N)`$ in . We also checked that the analogous, more complicated $`D=6`$ contour integrals lead to the same bulk indices, as one expects. For the groups not related by an isomorphism to SU$`(N)`$ the rational numbers in table 1 differ from the ones proposed in . I would be important to understand why. We also do not see how the arguments of section 8 of , which seemed to furnish a shortcut explanation of the SU$`(N)`$ results, could be adapted to reproduce the numbers highlighted in table 1. We next turn to numerical verification of these proposed bulk indices. Table 1: $`D=4`$ and $`D=6`$ bulk indices for the orthogonal, symplectic and exceptional groups of rank $`3`$ \| Group \| rank \| ind$`{}_{}{}^{D=4,6}{}_{0}{}^{}`$ \| \| --- \| --- \| --- \| \| SO(3) \| 1 \| $`1/4`$ \| \| SO(4) \| 2 \| $`1/16`$ \| \| SO(5) \| 2 \| $`\mathrm{𝟗}/\mathrm{𝟔𝟒}`$ \| \| SO(6) \| 3 \| $`1/16`$ \| \| SO(7) \| 3 \| $`\mathrm{𝟐𝟓}/\mathrm{𝟐𝟓𝟔}`$ \| \| Sp(2) \| 1 \| $`1/4`$ \| \| Sp(4) \| 2 \| $`\mathrm{𝟗}/\mathrm{𝟔𝟒}`$ \| \| Sp(6) \| 3 \| $`\mathrm{𝟓𝟏}/\mathrm{𝟓𝟏𝟐}`$ \| \| G<sub>2</sub> \| 2 \| $`\mathrm{𝟏𝟓𝟏}/\mathrm{𝟖𝟔𝟒}`$ \| ## V Monte Carlo evaluation of Yang-Mills Integrals As in previous works , we evaluate the Yang-Mills integrals using Monte Carlo methods. Both the Pfaffian polynomial $`𝒫_{D,G}`$ and the action $`𝒮=\frac{1}{4g^2}\mathrm{Tr}[X_\mu ,X_\nu ][X_\mu ,X_\nu ]`$ in eq.(2) are homogeneous functions of the $`X_\mu ^A`$ $$X_\mu ^A\alpha X_\mu ^A(\mu ;A)\{\begin{array}{ccc}𝒫(\{X_\mu ^A\})& & \alpha ^{\frac{1}{2}𝒩\mathrm{dim}(G)}𝒫(\{X_\mu ^A\})\hfill \\ & & \\ 𝒮(\{X_\mu ^A\})& & \alpha ^4𝒮(\{X_\mu ^A\})\hfill \end{array}.$$ (20) We introduce polar coordinates $`(X_1^1,\mathrm{},X_{\mathrm{dim}(G)}^D)=(\mathrm{\Omega }_d,R)`$, with $`d=D\mathrm{dim}(G)`$ the total dimension of the integral. As an example, $`𝒫(\mathrm{\Omega },1)`$ and $`𝒮(\mathrm{\Omega }_d,1)`$ denote the value of the Pfaffian polynomial and the action, respectively, for a configuration $`(X_1^1,\mathrm{},X_{\mathrm{dim}(G)}^D)`$ on the surface of the $`d`$dimensional unit hyper-sphere, with polar coordinates $`(\mathrm{\Omega }_d,R=1)`$). Eq.(20) allows us to to perform the $`R`$-integration analytically for each value of $`\mathrm{\Omega }_d`$, and to express the Yang-Mills integral as an expectation value over these angular variables: $$𝒵_{D,G}^𝒩=\frac{𝒟\mathrm{\Omega }_dz_G(\mathrm{\Omega }_d)}{𝒟\mathrm{\Omega }_d},$$ (21) with $$z_G(\mathrm{\Omega }_d)=2^{d/21}\frac{\mathrm{\Gamma }\left(\frac{d}{4}+\frac{𝒩}{8}\mathrm{dim}(G)\right)}{\mathrm{\Gamma }\left(\frac{d}{2}\right)}\times \frac{𝒫(\mathrm{\Omega },1)}{\left[𝒮(\mathrm{\Omega }_d,1)\right]^{\frac{d}{4}+\frac{𝒩}{8}\mathrm{dim}(G)}}.$$ (22) As discussed previously, the integrand $`z_G(\mathrm{\Omega }_d)`$ is too singular to be obtained by direct sampling of random points on the surface of the unit hyper sphere. Slightly modifying our procedure from , we therefore write $$𝒵_{D,G}^𝒩=\left[\frac{𝒟\mathrm{\Omega }_dz\times z^{\alpha _11}}{𝒟\mathrm{\Omega }_dz}\right]^1\left[\frac{𝒟\mathrm{\Omega }_dz^{\alpha _1}\times z^{\alpha _2\alpha _1}}{𝒟\mathrm{\Omega }_dz^{\alpha _1}}\right]^1\left[\frac{𝒟\mathrm{\Omega }_dz^{\alpha _2}\times z^{\alpha _2}}{𝒟\mathrm{\Omega }_dz^{\alpha _2}}\right]^1.$$ (23) Each of the terms $`[]`$ in eq.(23) is computed in a separate run. For example, the second quotient in eq.(23): $$\frac{𝒟\mathrm{\Omega }_dz^{\alpha _1}\times z^{\alpha _2\alpha _1}}{𝒟\mathrm{\Omega }_dz^{\alpha _1}}=<z^{\alpha _2\alpha _1}>_{\alpha _1},$$ (24) is simply the average of $`z^{\alpha _2\alpha _1}`$ for points $`\mathrm{\Omega }_d`$ on the unit hyper-sphere distributed according to the probability distribution $`\pi (\mathrm{\Omega }_d)z^{\alpha _1}(\mathrm{\Omega }_d)`$ (As it stands, eq.(24) is immediately applicable to $`D=4`$, where the integrand $`z`$ is positive semi-definite . In the general case, we have to sample with $`\|z^{\alpha _1}(\mathrm{\Omega }_d)\|`$ (cf. )). We sample angular variables $`\mathrm{\Omega }_d`$ according to $`z^{\alpha _1}`$ with a Metropolis Markov-chain method, which we now explain: At each iteration of the procedure, two distinct indices $`(A_1,\mu _1)`$ and $`(A_2,\mu _2)`$, and an angle $`0<\varphi <2\pi `$ are chosen randomly. An unbiased trial move $`\mathrm{\Omega }_d\mathrm{\Omega }_d^{}`$ is then constructed by modifying solely the coordinates $`X_{\mu _1}^{A_1}`$ and $`X_{\mu _2}^{A_2}`$: $$\left[\begin{array}{c}X_{\mu _1}^{A_1}\\ \\ X_{\mu _2}^{A_2}\end{array}\right]\left[\begin{array}{c}X_{\mu _1}^{A_1}\\ \\ X_{\mu _2}^{A_2}\end{array}\right]^{}=\sqrt{(X_{\mu _1}^{A_1})^2+(X_{\mu _2}^{A_2})^2}\left[\begin{array}{c}\mathrm{sin}(\varphi )\\ \\ \mathrm{cos}(\varphi )\end{array}\right],$$ (25) all other elements of $`\{X_\mu ^A\}`$ remaining unchanged. The trial move eq.(25) preserves the norm $`R`$ of the vector $`\{X_\mu ^A\}`$, i. e. keeps the configuration on the surface of the unit sphere. Furthermore, it is unbiased (the probability to propose $`\mathrm{\Omega }_d\mathrm{\Omega }_d^{}`$ is the same as for the reverse move). Finally, the move (for the example in eq. (24)) is accepted according to the Metropolis acceptance probability $$P(\mathrm{\Omega }\mathrm{\Omega }^{})=\mathrm{min}(1,\frac{z^{\alpha _1}(\mathrm{\Omega }^{})}{z^{\alpha _1}(\mathrm{\Omega })}).$$ (26) Empirically, we found the values $`\alpha _1=0.95,\alpha _2=0.6`$ to be appropriate. Each of the averages in eq.(23) was computed within between a few hours and more than a thousand hours of computer time (on a work station array), corresponding to a maximum of $`5\times 10^9`$ samples. Results are presented in the table below. Table 2: Direct evaluation of Yang-Mills integrals \| Group \| Monte Carlo result \| \| Exact \| \| --- \| --- \| --- \| --- \| \| $`G`$ \| $`𝒵_{D=4,G}^{𝒩=4}`$ \| \| \| \| SO(3) \| 1.255 $`\pm `$ \| 0.003 \| 1.2533… \| \| SO(4) \| 0.197 $`\pm `$ \| 0.004 \| 0.1963… \| \| SO(5) \| 0.589 $`\pm `$ \| 0.004 \| 0.589… \| \| SO(6) \| 0.0407 $`\pm `$ \| 0.0007 \| 0.04101… \| \| SO(7) \| 0.0169 $`\pm `$ \| 0.0003 \| 0.01708… \| \| Sp(2) \| 1.253 $`\pm `$ \| 0.001 \| 1.2533… \| \| Sp(4) \| 18.65 $`\pm `$ \| 0.2 \| 18.849… \| \| Sp(6) \| 279.2 $`\pm `$ \| 9.7 \| 285.59… \| \| G<sub>2</sub> \| 6943 $`\pm `$ \| 120 \| 7011.4… \| Dividing the Monte Carlo results for $`Z_{D=4,G}^{𝒩=4}`$ by the corresponding group volume factors (cf. eqs (12), (13), and (14)) we arrive at our numerical predictions for the bulk indices ind$`{}_{0}{}^{D=4}(G)`$, which we compare below to the proposed analytical values. Table 3: Monte Carlo results for the $`D=4`$ bulk index \| Group \| Monte Carlo \| \| Exact \| \| \| --- \| --- \| --- \| --- \| --- \| \| $`G`$ \| ind$`{}_{0}{}^{D=4}(G)`$ \| \| (Table 1) \| \| \| SO(3) \| 0.2503 $`\pm `$ \| 0.0006 \| 0.25 ( \| 1/4) \| \| SO(4) \| 0.0627 $`\pm `$ \| 0.0013 \| 0.0625 ( \| 1/16) \| \| SO(5) \| 0.1406 $`\pm `$ \| 0.001 \| 0.1406 ( \| 9/64) \| \| SO(6) \| 0.0620 $`\pm `$ \| 0.001 \| 0.0625 ( \| 1/16) \| \| SO(7) \| 0.0966 $`\pm `$ \| 0.0017 \| 0.0976 ( \| 25/256) \| \| Sp(2) \| 0.2500 $`\pm `$ \| 0.0002 \| 0.25 ( \| 1/4) \| \| Sp(4) \| 0.139 $`\pm `$ \| 0.0015 \| 0.1406 ( \| 9/64) \| \| Sp(6) \| 0.0973 $`\pm `$ \| 0.003 \| 0.0996 ( \| 51/512) \| \| G<sub>2</sub> \| 0.173 $`\pm `$ \| 0.003 \| 0.1747 ( \| 151/864) \| Agreement between the Monte Carlo results and theory is excellent, both in cases where rigorous results are known (SO(3), SO(4), Sp(2)) and where the deformation technique was applied. Among the latter cases, we again indicate in bold type new values, which had not been obtained before. ## VI Bosonic convergence for orthogonal, symplectic and exceptional integrals Our qualitative Monte Carlo method (cf. ,,) allows us to determine the convergence properties of the bosonic Yang-Mills integrals $`Z_{D,G}^{𝒩=0}`$. We have found the following: $$\begin{array}{ccc}\begin{array}{c}𝒵_{D,\mathrm{SO}(N=3,4)}^{𝒩=0}\hfill \\ \\ 𝒵_{D,\mathrm{Sp}(2)}^{𝒩=0}\hfill \end{array}\}\hfill & <\mathrm{}\mathrm{for}\hfill & D5\hfill \\ & & \\ \begin{array}{c}𝒵_{D,\mathrm{SO}(N5)}^{𝒩=0}\hfill \\ \\ 𝒵_{D,\mathrm{Sp}(4,6,\mathrm{})}^{𝒩=0}\hfill \\ \\ 𝒵_{D,\mathrm{G}_2}^{𝒩=0}\hfill \end{array}\}\hfill & <\mathrm{}\mathrm{for}\hfill & D3\hfill \end{array}$$ (27) All other bosonic integrals diverge. We thus obtain conditions which are fully consistent with the group isomorphisms discussed in the appendix. Let us note that we also performed the same qualitative computations for the susy integrals, as an important check of the convergence of the underlying Markov chains during the simulation. ## VII Conclusions and outlook In this paper we provided further evidence that Yang-Mills integrals encode surprisingly rich and subtle structures, which may prove to have important bearings on gauge and string theory. The chief result of the present paper was to demonstrate that Yang-Mills integrals, as well as the methods to study them, can be naturally generalized from the previously studied special unitary symmetries to other gauge symmetries. We numerically evaluated the partition functions for all semi-simple gauge groups of rank $`r3`$ and compared the results to conjectured exact values, which were obtained by a generalization of certain contour integrals derived from a supersymmetric deformation procedure. The connection between the Yang-Mills integrals and the bulk indices is provided by the group volumes, that we computed explicitly. We provided details on these very subtle calculations. Agreement between the approaches is perfect within the tight error margins left by our Monte Carlo technique. It would be very interesting to gain a simpler understanding of the rational numbers collected in table 1, although it is already evident that the bulk indices of the groups in question are more complicated than those of the special unitary case. In particular it would be nice to find general formulas for arbitrary rank. In the present paper we have focused on $`D=4`$ since this is the case where our numerical approach is most accurate. One should clearly study the dimensions $`D=6`$ and, especially, $`D=10`$ as well. It is straightforward, if more involved, to work out the predictions of the deformation method for these cases, at least for low rank gauge groups. In exact values for the total (bulk plus boundary) Witten index ind$`{}_{}{}^{D}(G)`$ were proposed. In $`D=4,6`$ one should have ind$`{}_{}{}^{D=4,6}(G)=0`$ while for $`D=10`$ it is argued to be a positive integer which, for groups other than SU$`(N)`$, can be larger than one. It is important to check the arguments by computing the index from the path integral. The results in the present paper indicate that the bulk contributions ind$`{}_{0}{}^{D}(G)`$ are correctly reproduced by the defomation method; however, we still lack a reliable method for computing the boundary terms ind$`{}_{1}{}^{D}(G)`$. Since the deformation method of successfully reproduces the partition functions, it is natural to ask whether it can be extended to calculate correlation functions of the ensembles eq.(1), such as the quantities studied numerically (so far only for SU$`(N)`$) in ,,. This would likely lead to new insights both in string theory and gauge theory ,. Numerically, it might be interesting to compare SU$`(N)`$, SO$`(N)`$ and Sp$`(2N)`$ for large values of $`N`$, as in the standard large $`N\mathrm{}`$ limit of ‘t Hooft these groups are expected to lead to identical results. ###### Acknowledgements. We thank H. Nicolai, H. Samtleben, G. Schröder and especially A. Smilga for useful discussions. This work was supported in part by the EU under Contract FMRX-CT96-0012. ## A Details and conventions for $`\mathrm{SO}(N)`$ The Lie algebra so$`(N)`$ has $`\frac{1}{2}N(N1)`$ generators which we choose to be the following standard anti-symmetric matrices $`T^{pq}`$ (i.e. the index $`A`$ becomes a double index $`pq`$) $`(T^{pq})_{jk}=\delta _j^p\delta _k^q\delta _k^p\delta _j^q`$ where $`p<q`$ $`(p,q=1,2,\mathrm{},N)`$. The Lie algebra reads then $`[T^{pq},T^{rs}]=\delta ^{qr}T^{ps}\delta ^{qs}T^{pr}\delta ^{pr}T^{qs}+\delta ^{ps}T^{qr}={\displaystyle \underset{t<u}{}}f^{pq,rs,tu}T^{tu}.`$ In this basis the generators are normalized as $`\mathrm{Tr}T^{pq}T^{rs}=2\delta ^{pq,rs}`$ and the structure constants are given through $`f^{pq,rs,tu}=\frac{1}{2}\mathrm{Tr}\left(T^{pq}[T^{rs},T^{tu}]\right)`$. The gauge potentials are then $`X_\mu =_{p<q}X_\mu ^{pq}T^{pq}`$ and the SO$`(N)`$ Yang-Mills integrals in these conventions read $$𝒵_{D,\mathrm{SO}(N)}^𝒩=\underset{p<q}{\overset{\frac{1}{2}N(N1)}{}}\underset{\mu =1}{\overset{D}{}}\frac{dX_\mu ^{pq}}{\sqrt{2\pi }}\mathrm{exp}\left[\frac{1}{8}\mathrm{Tr}[X_\mu ,X_\nu ][X_\mu ,X_\nu ]\right]𝒫_{D,N}(X),$$ (A1) where $`𝒫_{D,N}`$ is the Pfaffian as defined in eq.(3). These conventions are such that, in view of the isomorphism $`\mathrm{so}(3)=\mathrm{su}(2)`$ we have for all $`𝒩`$ (i.e. $`𝒩=0,4,8,16`$) $$𝒵_{D,\mathrm{SO}(3)}^𝒩=𝒵_{D,\mathrm{SU}(2)}^𝒩=𝒵_{D,\mathrm{Sp}(2)}^𝒩,$$ (A2) where the sympletic case is discussed in appendix B. One checks that in these normalizations the isomorphism $`\mathrm{so}(4)=\mathrm{so}(3)\mathrm{so}(3)`$ results in $$𝒵_{D,\mathrm{SO}(4)}^𝒩=2^{\frac{3}{2}(D\frac{1}{2}𝒩)}\left(𝒵_{D,\mathrm{SO}(3)}^𝒩\right)^2,$$ (A3) while the isomorphism $`\mathrm{so}(6)=\mathrm{su}(4)`$ leads to $$𝒵_{D,\mathrm{SO}(6)}^𝒩=2^{\frac{15}{4}(D\frac{1}{2}𝒩)}𝒵_{D,\mathrm{SU}(4)}^𝒩,$$ (A4) where the SU$`(4)`$ integral is defined as in . Finally, the isomorphism $`\mathrm{so}(5)=\mathrm{sp}(4)`$ translates into $$𝒵_{D,\mathrm{SO}(5)}^𝒩=2^{\frac{5}{2}(D\frac{1}{2}𝒩)}𝒵_{D,\mathrm{Sp}(4)}^𝒩,$$ (A5) where the Sp$`(4)`$ integral is defined in appendix B. One verifies that these isomorphisms are in perfect agreement with the results of the present paper as well as with . We next provide the details necessary for verifying the group volume factor $`_{\mathrm{SO}(N)}`$ of eq.(12). The natural Cartan subalgebra is spanned by the generators $`T^{12},T^{34},\mathrm{}`$. The corresponding maximal compact tori are given for SO$`(2N+1)`$ by the $`(2N+1)\times (2N+1)`$ matrix $$T=\left(\begin{array}{ccccc}\mathrm{rot}\theta _1& & & & \\ & \mathrm{rot}\theta _2& & & \\ & & \mathrm{}& & \\ & & & \mathrm{rot}\theta _N& \\ & & & & 1\end{array}\right)$$ (A6) while for SO$`(2N)`$ one has the $`2N\times 2N`$ matrix $$T=\left(\begin{array}{cccc}\mathrm{rot}\theta _1& & & \\ & \mathrm{rot}\theta _2& & \\ & & \mathrm{}& \\ & & & \mathrm{rot}\theta _N\end{array}\right)$$ (A7) Here $`\mathrm{rot}\theta _i`$ are the $`2\times 2`$ rotation matrices $$\mathrm{rot}\theta _i=\left(\begin{array}{cc}\mathrm{cos}\theta _i& \mathrm{sin}\theta _i\\ \mathrm{sin}\theta _i& \mathrm{cos}\theta _i\end{array}\right)$$ (A8) and matrix elements with no entries are zero. The corresponding reduced, normalized Haar measure on SO$`(2N+1)`$ (i.e. eq.(7)) reads $$𝒟T=\frac{2^{2N^2}}{2^NN!}\underset{i=1}{\overset{N}{}}\frac{d\theta _i}{2\pi }\underset{i<j}{\overset{N}{}}\mathrm{sin}^2\left(\frac{\theta _i\theta _j}{2}\right)\mathrm{sin}^2\left(\frac{\theta _i+\theta _j}{2}\right)\underset{i=1}{\overset{N}{}}\mathrm{sin}^2\left(\frac{\theta _i}{2}\right)$$ (A9) while for SO$`(2N)`$ one has $$𝒟T=\frac{2^{2N(N1)}}{2^{N1}N!}\underset{i=1}{\overset{N}{}}\frac{d\theta _i}{2\pi }\underset{i<j}{\overset{N}{}}\mathrm{sin}^2\left(\frac{\theta _i\theta _j}{2}\right)\mathrm{sin}^2\left(\frac{\theta _i+\theta _j}{2}\right).$$ (A10) Eqs.(A9),(A10) may also be used to work out the detailed form of the contour integrals eqs.(16),(17) of section 4: One simply expands the Haar measure around $`\theta _i0`$. Finally we recall the center groups of SO$`(N)`$: One has $`Z_{\mathrm{SO}(2N+1)}=\{\text{1}\}`$ and $`Z_{\mathrm{SO}(2N)}=\{\text{1},\text{1}\}`$ and therefore $`\|Z_{\mathrm{SO}(2N+1)}\|=1`$ and $`\|Z_{\mathrm{SO}(2N)}\|=2`$. ## B Details and conventions for $`\mathrm{Sp}(2N)`$ The Lie algebra sp$`(2N)`$ has $`2N^2+N`$ generators which we choose as follows. Define the $`N\times N`$ matrices $`E^{pq}`$ ($`p,q,j,k=1,2,\mathrm{},N`$) $`(E^{pq})_{jk}=\delta _j^p\delta _k^q`$ Then we define the $`N\times N`$ matrix generators $`T^{a,pq}`$ for $`p<q`$ by $`T^{0,pq}={\displaystyle \frac{1}{2\sqrt{2}}}\left(\begin{array}{cc}E^{pq}E^{qp}& 0\\ 0& E^{pq}E^{qp}\end{array}\right)T^{1,pq}={\displaystyle \frac{1}{2\sqrt{2}}}\left(\begin{array}{cc}0& i(E^{pq}+E^{qp})\\ i(E^{pq}+E^{qp})& 0\end{array}\right)`$ $`T^{2,pq}={\displaystyle \frac{1}{2\sqrt{2}}}\left(\begin{array}{cc}0& (E^{pq}+E^{qp})\\ (E^{pq}+E^{qp})& 0\end{array}\right)T^{3,pq}={\displaystyle \frac{1}{2\sqrt{2}}}\left(\begin{array}{cc}i(E^{pq}+E^{qp})& 0\\ 0& i(E^{pq}+E^{qp})\end{array}\right)`$ and the remaining generators are ($`p=1,\mathrm{},N`$) $`T^{1,pp}={\displaystyle \frac{1}{2}}\left(\begin{array}{cc}0& iE^{pp}\\ iE^{pp}& 0\end{array}\right)`$ $`T^{2,pp}={\displaystyle \frac{1}{2}}\left(\begin{array}{cc}0& E^{pp}\\ E^{pp}& 0\end{array}\right)T^{3,pp}={\displaystyle \frac{1}{2}}\left(\begin{array}{cc}iE^{pp}& 0\\ 0& iE^{pp}\end{array}\right)`$ In this basis the generators are normalized as $`\mathrm{Tr}T^AT^B=\frac{1}{2}\delta ^{AB}`$ and the structure constants are given through $`f^{ABC}=2\mathrm{T}\mathrm{r}T^A[T^B,T^C]`$ where $`A`$ is the multi-index $`(a,pq)`$. The Cartan subalgebra is spanned by the generators $`T^{3,pp}`$ with $`p=1,\mathrm{},N`$. The corresponding maximal compact torus is given by the matrix $$T=\left(\begin{array}{cccccc}e^{i\theta _1}& & & & & \\ & \mathrm{}& & & & \\ & & e^{i\theta _N}& & & \\ & & & e^{i\theta _1}& & \\ & & & & \mathrm{}& \\ & & & & & e^{i\theta _N}\end{array}\right)$$ (B1) The corresponding normalized Haar measure reads $$𝒟T=\frac{2^{2N^2}}{2^NN!}\underset{i=1}{\overset{N}{}}\frac{d\theta _i}{2\pi }\underset{i<j}{\overset{N}{}}\mathrm{sin}^2\left(\frac{\theta _i\theta _j}{2}\right)\mathrm{sin}^2\left(\frac{\theta _i+\theta _j}{2}\right)\underset{i=1}{\overset{N}{}}\mathrm{sin}^2\theta _i$$ (B2) The center of Sp$`(2N)`$ is the group $`Z_{\mathrm{Sp}(2N)}=\{\text{1},\text{1}\}`$ and thus $`\|Z_{\mathrm{Sp}(2N)}\|=2`$. ## C Details and conventions for $`\mathrm{G}_2`$ The Lie algebra $`G_2`$ has 14 generators. For the fundamental representation we choose them as the following explicit $`7\times 7`$ matrices (see e.g. ). Define $`(X^{p,q})_{jk}=\delta _j^p\delta _k^q\delta _k^p\delta _j^q`$ $`(Y^{p,q})_{jk}=i(\delta _j^p\delta _k^q+\delta _k^p\delta _j^q)`$ Then $`T^1`$ $`=`$ $`(X^{1,2}+X^{3,4}){\displaystyle \frac{1}{2\sqrt{6}}}+(X^{5,6}+X^{6,7}){\displaystyle \frac{1}{2\sqrt{3}}}`$ (C1) $`T^2`$ $`=`$ $`(Y^{1,2}+Y^{3,4}){\displaystyle \frac{1}{2\sqrt{6}}}+(Y^{5,6}+Y^{6,7}){\displaystyle \frac{1}{2\sqrt{3}}}`$ (C2) $`T^3`$ $`=`$ $`(X^{1,7}X^{4,5}){\displaystyle \frac{1}{2\sqrt{2}}}`$ (C3) $`T^4`$ $`=`$ $`(Y^{1,7}+Y^{4,5}){\displaystyle \frac{1}{2\sqrt{2}}}`$ (C4) $`T^5`$ $`=`$ $`(X^{1,6}+X^{4,6}){\displaystyle \frac{1}{2\sqrt{3}}}+(X^{2,7}X^{3,5}){\displaystyle \frac{1}{2\sqrt{6}}}`$ (C5) $`T^6`$ $`=`$ $`(Y^{1,6}Y^{4,6}){\displaystyle \frac{1}{2\sqrt{3}}}+(Y^{2,7}+Y^{3,5}){\displaystyle \frac{1}{2\sqrt{6}}}`$ (C6) $`T^7`$ $`=`$ $`(X^{1,3}+X^{2,4}){\displaystyle \frac{1}{2\sqrt{2}}}`$ (C7) $`T^8`$ $`=`$ $`(Y^{1,3}+Y^{2,4}){\displaystyle \frac{1}{2\sqrt{2}}}`$ (C8) $`T^9`$ $`=`$ $`(X^{1,5}+X^{4,7}){\displaystyle \frac{1}{2\sqrt{6}}}+(X^{2,6}X^{3,6}){\displaystyle \frac{1}{2\sqrt{3}}}`$ (C9) $`T^{10}`$ $`=`$ $`(Y^{1,5}Y^{4,7}){\displaystyle \frac{1}{2\sqrt{6}}}+(Y^{2,6}+Y^{3,6}){\displaystyle \frac{1}{2\sqrt{3}}}`$ (C10) $`T^{11}`$ $`=`$ $`(X^{2,5}X^{3,7}){\displaystyle \frac{1}{2\sqrt{2}}}`$ (C11) $`T^{12}`$ $`=`$ $`(Y^{2,5}+Y^{3,7}){\displaystyle \frac{1}{2\sqrt{2}}}`$ (C12) Finally, the matrices $`T^{13}`$ and $`T^{14}`$ are diagonal matrices with elements: $$T^{13}=\mathrm{diag}(\frac{i}{2\sqrt{6}};\frac{i}{2\sqrt{6}};\frac{i}{2\sqrt{6}};\frac{i}{2\sqrt{6}};\frac{i}{\sqrt{6}};0;\frac{i}{\sqrt{6}})$$ (C13) $$T^{14}=\mathrm{diag}(\frac{i}{2\sqrt{2}};\frac{i}{2\sqrt{2}};\frac{i}{2\sqrt{2}};\frac{i}{2\sqrt{2}};0;0;0)$$ (C14) In this basis the generators are normalized as $`\mathrm{Tr}T^AT^B=\frac{1}{2}\delta ^{AB}`$ and the structure constants are given through $`f^{ABC}=2\mathrm{T}\mathrm{r}T^A[T^B,T^C]`$. The Cartan subalgebra is spanned by the generators $`T^{13}`$ and $`T^{14}`$. The corresponding maximal compact torus is given by the matrix $$T=\left(\begin{array}{ccccccc}e^{i\theta _1}& & & & & & \\ & e^{i\theta _2}& & & & & \\ & & e^{i\theta _2}& & & & \\ & & & e^{i\theta _1}& & & \\ & & & & e^{i(\theta _1+\theta _2)}& & \\ & & & & & 1& \\ & & & & & & e^{i(\theta _1+\theta _2)}\end{array}\right)$$ (C15) The normalized Haar measure on G<sub>2</sub> with respect to this torus reads $$𝒟T=\frac{2^{12}}{12}\frac{d\theta _1}{2\pi }\frac{d\theta _2}{2\pi }\mathrm{sin}^2\left(\frac{\theta _1}{2}\right)\mathrm{sin}^2\left(\frac{\theta _2}{2}\right)\mathrm{sin}^2\left(\frac{\theta _1\theta _2}{2}\right)\mathrm{sin}^2\left(\frac{\theta _1+\theta _2}{2}\right)\mathrm{sin}^2\left(\frac{2\theta _1+\theta _2}{2}\right)\mathrm{sin}^2\left(\frac{\theta _1+2\theta _2}{2}\right)$$ (C16) The center of G<sub>2</sub> is trivial: $`Z_{\mathrm{G}_2}=\{\text{1}\}`$ and thus $`\|Z_{\mathrm{G}_2}\|=1`$.
warning/0004/quant-ph0004051.html	ar5iv	text	# Persistent entanglement in arrays of interacting particles ## Abstract We study the entanglement properties of a class of $`N`$ qubit quantum states that are generated in arrays of qubits with an Ising-type interaction. These states contain a large amount of entanglement as given by their Schmidt measure. They have also a high persistency of entanglement which means that $`N/2`$ qubits have to be measured to disentangle the state. These states can be regarded as an entanglement resource since one can generate a family of other multi-particle entangled states such as the generalized GHZ states of $`<N/2`$ qubits by simple measurements and classical communication (LOCC). The notion of entanglement has many facets. A modern perspective is to regard it as a resource for certain communicational and computational tasks . Related to this viewpoint is the problem of identifying equivalence classes of entangled states, and to find relations between these classes. While for pure states of bi-partite systems there is a single “unit” of entanglement – the entanglement contained in a Bell state – it has recently become clear that for systems shared by three and more parties there are several inequivalent classes of entangled states . Progress in the understanding of multi-particle entanglement has been triggered by giving explicit examples of states that did not fit into existing classification schemes. Generally speaking, a sufficiently rich phenomenology of entangled states is needed. It helps us to refine entanglement classification schemes and, arguably, to motivate them in the first place. In this paper, we introduce a class of $`N`$-qubit entangled states which is different from both the GHZ class and the recently introduced W class of $`N`$-qubit states . We also give an operational characterization of these classes in terms of local measurements. In one respect, the states we are going to describe resemble the so-called maximally entangled Greenberger-Horne-Zeilinger (GHZ) states of $`N`$ qubits, while in some other respect they are much more entangled than the GHZ states. To characterize these states, we introduce the notions of maximal connectedness and persistency of entanglement of an entangled state. The first notion emphasizes the possibility that, in an $`N`$-particle state, even when the reduced density matrix of a subset of particles is fully separable , it may still be possible to project that subset of particles into a highly entangled state by performing local measurements on the other particles \[supplemented with classical communication and local operations (LOCC), if the parties are remotely separated\]. The second notion relates the amount of entanglement in a multi-particle system to the operational effort it takes (in terms of local operations) to destroy all entanglement in the system. The states we describe occur, for example, in the quantum Ising model of spin chains and, more generally, spin lattices. We will introduce the states in the context of this specific model; their entanglement properties, however, are discussed in general terms by assuming, as usual, that the qubits are distributed between remote parties which can only act through LOCC. Consider an ensemble of qubits that are located on a $`d`$-dimensional lattice ($`d=1,2,3`$) at sites $`a^d`$ and interact via some short-range interaction described by the Hamiltonian $$H_{\text{int}}=\mathrm{}g(t)\underset{a,a^{}}{}f(aa^{})\frac{1+\sigma _z^{(a)}}{2}\frac{1\sigma _z^{(a^{})}}{2}$$ (1) Concerning the entanglement properties of the states we are going to investigate, this interaction Hamiltonian is equivalent to the quantum Ising model with $`H_{\text{int}}^{}=_{a,a^{}}\frac{1}{4}\mathrm{}g(t)f(aa^{})\sigma _z^{(a)}\sigma _z^{(a^{})}`$, where the indices $`a,a^{}`$ run over all occupied lattices sites. The coupling strength is written as a product $`gf`$, where $`f(aa^{})`$ specifies the interaction range and the $`g(t)`$ allows for a possible overall time dependence. In this letter, we confine ourselves to next-neighbor interactions. A more general situation will be reported in . In the language of quantum information, the interaction (1) realizes simultaneous conditional phase gates between qubits at neighboring sites $`a`$ and $`a^{}`$. For an experimental realization see the discussion at the end of the paper. Consider first the one-dimensional example of a chain of $`N`$ qubits (“spin chain”) with next-neigbor interaction $`f(aa^{})=\delta _{a+1,a^{}}`$. Initially, all qubits are prepared in the state $`(\|0_a+\|1_a)/\sqrt{2}`$, where $`\|0_a\|0_{z,a}`$ and $`\|1_a=\|1_{z,a}`$ are eigenstates of $`(1\sigma _z^{(a)})/2`$ with eigenvalues $`0`$ and $`1`$, respectively. (This is the most interesting situation; if they are prepared in states $`\|0_a`$ or $`\|1_a`$, no entanglement will build up.) The unitary transformation generated by (1) is $`U(\phi )=\mathrm{exp}(\mathrm{i}\phi _a\frac{1+\sigma _z^{(a)}}{2}\frac{1\sigma _z^{(a+1)}}{2})`$ with $`\phi =\text{d}tg(t)`$. For $`g(t)=g=`$const, $`U(\phi )=U(gt)`$ is periodic in time and generates “entanglement oscillations” of the chain. For the specific values $`\phi =0,2\pi ,4\pi ,\mathrm{}`$ the chain is disentangled, while for all other values of $`\phi `$, it is entangled. For the values $`\phi =\pi ,3\pi ,5\pi ,\mathrm{}`$ the chain is in some sense maximally entangled and we will concentrate on this situation in the following. The state can then be written in the form $$\|\varphi _N=\frac{1}{2^{N/2}}\underset{a=1}{\overset{N}{}}\left(\|0_a\sigma _z^{(a+1)}+\|1_a\right)$$ (2) with the convention $`\sigma _z^{(N+1)}1`$. The compact notation employed in (2) is easily understood by multiplying out the right hand side. For $`N=2`$, one obtains $`\|\varphi _2=\frac{1}{2}(\|0_1\sigma _z^{(2)}+\|1_1)(\|0_2+\|1_2)=\frac{1}{2}\left(\|0_1(\|0_2\|1_2)+\|1_1(\|0_2+\|1_2)\right)`$ which is a maximally entangled state. We may write it, up to a local unitary transformation on qubit 2 in the standard form $$\|\varphi _2=_{l.u.}\frac{1}{\sqrt{2}}(\|0_1\|0_2+\|1_1\|1_2),$$ (3) where “l.u.” indicates that the equality holds up to a local unitary transformation on one or more of the qubits . Similarly, one obtains for $`N=3,4`$ $`\|\varphi _3`$ $`=_{l.u.}`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left(\|0_1\|0_2\|0_3+\|1_1\|1_2\|1_3\right),`$ (4) $`\|\varphi _4`$ $`=_{l.u.}`$ $`{\displaystyle \frac{1}{2}}(\|0_1\|0_2\|0_3\|0_4+\|0_1\|0_2\|1_3\|1_4`$ (6) $`+\|1_1\|1_2\|0_3\|0_4\|1_1\|1_2\|1_3\|1_4).`$ While $`\|\varphi _3`$ corresponds to a GHZ state of three qubits , $`\|\varphi _4`$ is not equivalent to a 4-qubit GHZ state . More generally, the states $`\|\varphi _N`$ and the $`N`$-qubit GHZ state $`\|\mathrm{GHZ}_N2^{1/2}(\|0_1\mathrm{}\|0_N)+\|1_1\mathrm{}\|1_N)`$ are not equivalent for $`N>3`$ i.e. cannot be transformed into each other by LOCC (local transformations and classical communication) as we shall see below. How can we compare the entanglement properties of $`\|\varphi _4`$ and $`\|\mathrm{GHZ}_4`$ in operational terms? Imagine that the qubits are distributed between four remote parties, which may perform, as usual, local operations and classical communication. We observe: a) The states share the property that any two of the four qubits can be projected into a Bell state by measuring the other two qubits in an appropriate basis. In other words, the parties may use either of the states $`\|\varphi _4`$ or $`\|\mathrm{GHZ}_4`$ to teleport a qubit between any of the four parties. b) The states are different in that it is harder to destroy the entanglement of state $`\|\varphi _4`$ than that of $`\|\mathrm{GHZ}_4`$ by local operations. In fact, it is impossible to destroy all entanglement of $`\|\varphi _4`$ by a single local operation, such as a von Neumann measurement or complete depolarization of a qubit. For the state $`\|\mathrm{GHZ}_4`$, in contrast, a single local measurement suffices to bring it into a product state . These observations motivate us to introduce the following definitions. A local measurement in the following means a von Neumann measurement on a single qubit. Definition 1: (Max. connectedness) The quantum mechanical state of a set $`𝒞=\{1,2,\mathrm{},n\}`$ of $`n`$ qubits is maximally connected if any two qubits $`jk𝒞`$ can be projected, with certainty, into a pure Bell state by local measurements on a subset of the other qubits. Note that the state obtained may depend on the outcome of the measurements. Definition 2: (Persistency) The persistency of entanglement $`P_e`$ of an entangled state of $`n`$ qubits is the minimum number of local measurements such that, for all measurement outcomes, the state is completely disentangled. Since we are only concerned with pure states, a disentangled state means a product state of all $`n`$ qubits . Obviously, for all $`n`$-qubit states $`0P_en1`$. Definitions 1 and 2 can be straightforwardly generalized to arbitrary $`n`$-partite pure states. Note that the definitions 1 and 2 are invariant under the group of local unitary transformations on any of the qubits . In the sense of these definitions, both states $`\|\varphi _4`$ and $`\|\mathrm{GHZ}_4`$ are maximally connected, while their persistency is $`P_e=2`$ and $`P_e=1`$, respectively. More generally, for the state $`\|\varphi _N`$ we show that (i) it is maximally connected and (ii) its persistency is $`P_e(\|\varphi _N)=N/2`$. Property (ii) quantifies the operational effort that is needed to destroy all entanglement in the qubit chain. We also note that, (iii), the persistency of the states $`\|\varphi _N`$ is equal to their Schmidt measure : If one expands $`\|\varphi _N`$ into a product basis of the $`N`$ qubits, the minimum number of terms in such a generalized Schmidt representation grows exponentially and requires $`2^{N/2}`$ product terms. In that sense, the state $`\|\varphi _N`$ of the qubit chain is indeed much more entangled than most of the known $`N`$ qubit states. We now prove property (i). The cases $`N=2,3`$ are trivial as the state is a Bell or a GHZ state, respectively. For $`N>3`$, the proof goes as follows. Let us denote by $`\|0_{xj}(\|0_j+\|1_j)/\sqrt{2}`$, $`\|1_{xj}(\|0_j\|1_j)/\sqrt{2}`$ the eigenstates of $`\sigma _x^{(j)}`$. We first show that the qubits at the ends of the string, i.e., qubits $`1`$ and $`N`$ can be brought into a Bell state by measuring the qubits $`2,\mathrm{},N1`$. For easier book keeping, we use the notation $`\|\varphi _N\|\{1,2,3,\mathrm{},N\}_{\text{chain}}`$. Then the state can be expanded in the form $`\|\{1,2,3,\mathrm{},N\}_{\text{chain}}=(\|0_1\sigma _z^{(2)}+\|1_1)(\|0_2\sigma _z^{(3)}+\|1_2)\|\{3,4,\mathrm{},N\}_{\text{chain}}`$ where we suppress normalization factors. Measuring the operator $`\sigma _x^{(2)}`$ of qubit 2, we obtain for the remaining (unmeasured) qubits $`1,3,4,\mathrm{},N`$ the state $`{}_{x2}{}^{}ϵ_2\|\{1,2,3,\mathrm{},N\}_{\text{chain}}^{}=\{(1i\sigma _y^{(1)})/\sqrt{2},(\sigma _x^{(1)}+\sigma _z^{(1)})/\sqrt{2}\}`$ $`\|\{1,3,4,\mathrm{},N\}_{\text{chain}}`$ for the outcome $`ϵ_2=\{0,1\}`$, correspondingly. This state is, up to the local unitary transformations specified in the parenthesis, identical to an entangled chain of length $`N1`$, and gives us a recursion formula. We can repeat this procedure and measure qubit 3, and so on. We obtain $`_j`$$`(_{xj}ϵ_j\|)\|\{1,2,3,\mathrm{},N\}_{\text{chain}}=U_1\|\{1,N\}_{\text{chain}}`$ with $`U_1\{1,\sigma _x^{(1)},\sigma _y^{(1)},\sigma _z^{(1)}\}`$ for $`N`$ even and $`U_1\{(\sigma _x^{(1)}\pm \sigma _z^{(1)})/\sqrt{2},(1\pm i\sigma _y^{(1)})/\sqrt{2}\}`$ for $`N`$ odd, up to a phase factor. This is a Bell state. To bring any other qubits $`j,k`$ (w.l.o.g. $`j<k`$ ) from the chain $`\{1,2,\mathrm{},N\}`$ into a Bell state, we first measure the “outer” qubits $`1,2,\mathrm{}j1`$ and $`k+1,k+2,\mathrm{},N`$ in the $`\sigma _z`$ basis, which projects the qubits of the remaining chain $`j,j+1,\mathrm{},k`$ into the state $`U_jU_k\|\{j,j+1,\mathrm{},k1,k\}_𝒞`$ with $`U_j\{1,\sigma _z^{(j)}\}`$, $`U_k\{1,\sigma _z^{(k)}\}`$. A subsequent measurement of the “inner” qubits $`j+1,\mathrm{}k1`$ will then project qubits $`j,k`$ into a Bell state, as shown previously. To prove property (iii), we use the expansion $`\|\{1,2,3,\mathrm{},N,N+1,N+2\}_{\text{chain}}=`$ $`\|\{1,2,3,\mathrm{},N\}_{\text{chain}}(\|0_{N+1}\|1_{N+1}\sigma _z^{(N)})`$ $`(\|0_{N+2}\|1_{N+2}\sigma _z^{(N+1)})`$ which can be written in the form $`\|\varphi _{N+2}=\|\varphi _N\|0_{x,N+1}\|1_{z,N+2}(\sigma _z^{(N)}\|\varphi _N)\|1_{x,N+1}\|0_{z,N+2}`$. Denote the minimum number of product terms in an expansion of $`\|\varphi _N`$ by $`r`$. As this number is invariant under local unitary transformations , it is the same for the state $`\sigma _z^{(N)}\|\varphi _N`$. No term in an expansion of $`\|\varphi _N\|0_{x,N+1}\|1_{z,N+2}`$ can be combined with any term in an expansion of $`(\sigma _z^{(N)}\|\varphi _N)\|1_{x,N+1}\|0_{z,N+2}`$ into a single product term, since any nontrivial linear combination of $`\|0_{x,N+1}\|1_{z,N+2}`$ with $`\|1_{x,N+1}\|0_{z,N+2}`$ gives a non-product state w.r.t. qubit $`N+1`$ and $`N+2`$. The minimum number of product terms for an expansion of $`\|\varphi _{N+2}`$ is thus equal to $`2r`$. Since for $`N=2,3`$ we have $`r=2`$ \[see (3) and (6)\], it follows by induction that $`r=2^{N/2}`$. In other words, the Schmidt measure $`P_\text{S}(\|\varphi _N)`$ of $`\|\varphi _N`$ is equal to $`\mathrm{log}_2(r)=N/2`$. We now prove property (ii). An explicit strategy to disentangle state (2) is to measure $`\sigma _z^{(j)}`$ of all even numbered qubits, $`j=2,4,6,\mathrm{}`$, which can easily be verified. The total number of these measurements is $`N/2`$, which gives an upper bound to the persistency, i.e. $`P_e(\|\varphi _N)N/2`$. On the other hand, the Schmidt measure gives a lower bound to the persistency. This can be seen as follows. Since $`\|\varphi _N`$ can be disentangled by $`P_e`$ measurements, there exists an expansion of the form $`\|\varphi _N=_{j_1,\mathrm{},j_{P_e}=0}^1\|\mu _1^{(j_1)}_{a_1}\|\mu _2^{(j_2)}_{a_2}\mathrm{}\|\mu _{P_e}^{(j_{P_e})}_{a_{P_e}}\|\mathrm{prod}^{(j_1,\mathrm{},j_{P_e})}`$ where $`a_1,\mathrm{},a_{P_e}`$ are the measured atoms, $`\|\mu _1^{(j_1)}_{a_1},\mathrm{},\|\mu _{P_e}^{(j_{P_e})}_{a_{P_e}}`$ the resulting 1-qubit states for the measurement outcomes $`j_1,j_2,\mathrm{},j_{P_e}`$, and $`\|\mathrm{prod}^{(j_1,\mathrm{},j_{P_e})}`$ some (unnormalized) product states of the remaining qubits. This expansion contains at most $`2^{P_e}`$ product terms, and therefore $`P_\text{S}\mathrm{log}_2(2^{P_e})=P_e`$. Together with (iii) we obtain $`P_\text{S}(\|\varphi _N)=N/2P_e(\|\varphi _N)N/2`$ which proves property (ii). Results (ii) and (iii) show that the persistency of entanglement of the state $`\|\varphi _N`$ (2) coincides with its Schmidt measure. This result also holds for the state $`\|\mathrm{GHZ}_N`$. The meaning of these two concepts is, however, not the same. To illustrate this point, consider the so-called W state discussed in Ref. , $`\|W_N=N^{1/2}(\|1_1\|0_2\mathrm{}\|0_N+\|0_1\|1_2\mathrm{}\|0_N+\mathrm{}+\|0_1\|0_2\mathrm{}\|1_N)`$. The Schmidt measure of this state is equal to $`\mathrm{log}_2(N)`$ which means that the amount of entanglement contained in $`\|W_N`$ is smaller, in fact exponentially smaller, than in the state $`\|\varphi _N`$. The persistency of $`\|W_N`$, on the other hand, is given by $`N1`$ which means that the entanglement of $`\|W_N`$ is harder to destroy by local measurements than that of $`\|\varphi _N`$. This observation agrees with the findings of Ref. , who showed that any state obtained from $`\|W_N`$ by tracing over $`N2`$ qubits is inseparable. Note however that, different from $`\|\varphi _N`$ and $`\|\mathrm{GHZ}_N`$, the state $`\|W_N`$ is not maximally connected. In the remainder of the paper we will generalize some of the results to dimensions $`d=2`$ and $`d=3`$, i.e. to qubits arranged on a lattice. These cases are different from the case $`d=1`$ since there is no natural ordering of the qubits. Therefore, the concept of a “chain” of qubits does not apply anymore. The natural generalization to higher dimensions is a “cluster” $`𝒞`$ of qubits as in Fig. 1a. The precise definition of a cluster is the following: Let each lattice site be specified by a $`d`$-tuple of (positive or negative) integers $`a^d`$. Each site $`a`$ has $`2d`$ neighbouring sites. If occupied, these are the sites whose qubit interacts with the qubit at $`a`$. The set $`𝒜^d`$ specifies all sites that are occupied by a qubit. Two sites $`a,a^{}𝒜`$ are connected (in a topological sense) if there exists a sequence of neighboring sites that are all occupied, that is $`\{a^{(n)}\}_{n=1}^N𝒜`$ with $`a^{(1)}=a`$ and $`a^{(n)}=a^{}`$. A cluster $`𝒞𝒜`$ is a subset of $`𝒜`$ with the properties that first, any two sites $`c,c^{}𝒞`$ are connected and second, any sites $`c𝒞`$ and $`a𝒜\backslash 𝒞`$ are not connected. The quantum mechanical state of a cluster that is generated under the Hamiltonian (1) for $`\phi =\pi `$ is $$\|\mathrm{\Phi }_𝒞=\underset{c𝒞}{}\left(\|0_c\underset{\gamma \mathrm{\Gamma }}{}\sigma _z^{(c+\gamma )}+\|1_c\right)$$ (7) with the choice $`\mathrm{\Gamma }=\{(1,0),(0,1)\}`$ for $`d=2`$ and $`\mathrm{\Gamma }=\{(1,0,0),(0,1,0),(0,0,1)\}`$ for $`d=3`$, using the convention that $`\sigma _z^{(c+\gamma )}1`$ when $`c+\gamma 𝒞`$ (the qubit cannot be entangled with an empty site). The special case of the 1D chain (2) is obtained from (7) for the choice $`\mathrm{\Gamma }=\{1\}`$. The cluster states (7) satisfy the following set of eigenvalue equations: $$K_a\|\mathrm{\Phi }_𝒞=\kappa \|\mathrm{\Phi }_𝒞$$ (8) for the family of operators $`K_a=\sigma _x^{(a)}`$ $`_{\gamma \mathrm{\Gamma }\mathrm{\Gamma }}`$ $`\sigma _z^{(a+\gamma )}`$, $`a𝒞`$, where $`\mathrm{\Gamma }\mathrm{\Gamma }`$ specifies the sites of all qubits that interact with $`a`$, and $`\sigma _z^{(a+\gamma )}1`$ when $`a+\gamma 𝒞`$. The eigenvalue $`\kappa =\pm 1`$ is determined by the specific occupation pattern of the neighboring sites. For $`a+\{\mathrm{\Gamma }\mathrm{\Gamma }\}𝒞`$, for example, $`\kappa =(1)^d`$. The operators $`\{K_a\|a𝒞\}`$ form a complete set of commuting observables of which the cluster states $`\|\mathrm{\Phi }_𝒞`$are are eigenstates. Equations (8) can be used to generalize some of the entanglement properties from the 1D case to higher dimensions. Here we just report the results. A detailed proof will be given in a longer paper . We find that all cluster states are maximally connected. It is noteworthy that the property of maximal connectedness of $`\|\mathrm{\Phi }_𝒞`$does not depend on the precise shape of the cluster, and not even on its topological characterization except for being a cluster. Consider a cluster $`𝒞`$ and any two qubits on sites $`c^{},c^{\prime \prime }𝒞`$ as in Fig 1a. To bring these qubits into a Bell state, we first select a one-dimensional path $`𝒫𝒞`$ that connects sites $`c^{}`$ and $`c^{\prime \prime }`$ as in Fig. 1a. Then we measure all neighboring qubits surrounding this path in the $`\sigma _z`$ basis. By this procedure, we project the qubits on path $`𝒫`$ into a state that is, up to local unitary transformations, identical to the state $`\|\mathrm{\Phi }_N`$ of the linear chain. We have thereby reduced the two- and three-dimensional problem to the one-dimensional problem. Equations (8) can also be used to calculate the persistency, as they imply strict correlations among 1-particle measurements. These correlations can be used to minimize the number of measurements required to project $`\|\mathrm{\Phi }_𝒞`$into a product state. In general, the exact value of the persistency depends on the shape of a cluster. For large convex clusters, we can give the asymptotic result $`P_e/N=1/2`$ where $`N\mathrm{}`$ is the number of qubits. Entanglement is often regarded as a resource and thus the question arises which states can be obtained from cluster states by local operations and classical communication (LOCC). A particularly simple class of LOCC is obtained by restricting oneself to projective von Neumann measurements on selected qubits. We note without proof that from a block $`𝒞`$ of $`L^d`$ qubits, one can obtain any state of the form $`\alpha \|00\mathrm{}0_𝒞^{}+\beta \|11\mathrm{}1_𝒞^{}`$ of any subset of qubits $`𝒞^{}𝒞\{2\}^d`$. For $`\alpha =\beta =1/\sqrt{2}`$, this includes, in particular, the family of generalized (multi-particle) GHZ states on this subset. An illustration is given in Fig. 1b. Even though the thereby obtained states are highly entangled, their Schmidt entanglement measure is always smaller than of the original cluster state, and so the total amount of entanglement decreases. With the experimental progress in cooling and trapping of neutral atoms, one has identified systems such as “optical lattices” in which the interaction (1) can be implemented by cold atomic collisions or other techniques. These systems allow one, in particular, to switch on and off the coupling $`g(t)`$ between all qubits simultaneously by a manipulation of the parameters of the trapping lasers. The unitary transformation $`U(\phi )`$ \[before eq.(2)\] with $`\phi =\pi `$ can thereby be realized by a single global operation. This enables one, in principle, to create a variety of multi-particle entangled states such as $`\|\mathrm{GHZ}_\mathrm{M}`$ with $`M3`$ by the entanglement operation $`U(\phi )`$, followed by 1-qubit measurements and subsequent 1-qubit rotations (compare Fig. 1b). In conclusion we have introduced a class of highly entangled multi-qubit states. The cluster states have a large persistency of entanglement which quantifies the operational effort needed to disentangle these states. For the chain of qubits in the state $`\|\varphi _N`$, we have shown that the value of the persistency agrees with the Schmidt measure of $`\|\varphi _N`$. In that sense, the state $`\|\varphi _N`$ is indeed much more entangled than most known $`N`$ qubit states. The cluster states can be regarded as a (scalable) resource for other multi-qubit entangled states, such as multiparticle GHZ states. Experimentally, these states could be generated and studied in optical lattices or similar systems. We thank H. Aschauer, B.-G. Englert, J. Hersch, L. Hardy, A. Schenzle, and C. Simon for discussions. One of us (HJB) enjoyed delightful discussions with Jens Eisert that emerged from a hiking tour during the Benasque workshop 2000. We are also grateful to J. Eisert for comments on the manuscript. This work has been supported in part by the Schwerpunktsprogramm QIV of the DFG.
warning/0004/hep-ph0004012.html	ar5iv	text	# DO-TH 00/08April 2000 Off-forward parton distributions and Shuvaev’s transformations ## 1 Introduction The treatment of nonforward high-energy processes, such as deeply virtual Compton scattering (DVCS) and hard exclusive electroproduction of vector mesons, in perturbative QCD gives rise to a new class of parton distributions, the so-called skewed parton distributions (SPDs) and double distributions (DDs) , generalizing the well-known conventional parton distributions and, at the same time, the nucleon form factors. In the following we restrict ourselves to the off-forward parton distributions (OFPDs) introduced by Ji , which are equivalent to nonforward or off-diagonal parton distributions. Therefore, our results can be easily generalized to nonforward and off-diagonal parton distributions. The off-forward parton distributions $`H_p(x,\xi ,t,\mu )`$, which parametrize nonforward matrix elements of light cone bilocal operators $`P^{}\|𝒪_p(n/2,n/2)\|P\|_{n^2=0}`$, depend on the momentum fraction $`x`$ of the average nucleon momentum $`\overline{P}:=(P+P^{})/2`$, which the initial state parton $`p`$ carries, on the “skewedness” $`\xi =n\mathrm{\Delta }/2n\overline{P}`$ with $`\mathrm{\Delta }:=P^{}P`$, on the momentum transfer invariant $`t:=\mathrm{\Delta }^2`$, and on the renormalization scale $`\mu `$. For vanishing $`\xi `$ and $`t`$, they are identical to the usual forward parton distributions. Detailed reviews on off-forward parton distributions can be found, e.g., in Refs. . Recently, Shuvaev demonstrated that the off-forward parton distributions can be related (at least in leading order) by simple transformations to so-called effective forward parton distributions (EFPDs), the renormalization scale dependence of which is governed by the conventional forward evolution equations. These relations have led to some progress in determining the shape of the off-forward parton distributions for small values of $`\xi `$ , since the EFPDs can be identified with the usual partons for small values of $`\xi `$ and arbitrary scale $`\mu `$. In the present paper, we express nonforward amplitudes directly in terms of effective forward parton distributions. Furthermore, we define a family of self-consistent models for EFPDs, in which the effective forward parton distributions are obtained from the conventional forward parton distributions and nucleon form factors at arbitrary scale $`\mu `$. In the next section we briefly review the basic properties of the off-forward parton distributions, and we define the effective forward parton distributions. In Sec. 3, we recalculate Shuvaev’s inverse transformations, which relate the EFPDs to the OFPDs, and we derive their support in $`x`$, which is not identical to $`1x1`$ as for the conventional forward parton distributions.<sup>1</sup><sup>1</sup>1We use throughout parton distributions with both signs of $`x`$, i.e., $`𝑞(x)=\overline{q}(x)`$ and $`𝑔(x)=𝑔(x)`$. In Sec. 4, Shuvaev’s transformation is brought into a form that is convenient for a further analytical and numerical treatment. In Sec. 5, we connect the effective forward parton distributions directly to nonforward amplitudes, and we briefly discuss the reliability of simple approximative formulas. In Sec. 6, we introduce our model. Finally, in Sec. 7, we summarize our results, and we draw the conclusions. ## 2 Off-forward and effective parton distributions The long-distance behavior of hard scattering processes, which is not calculable in (QCD) perturbation theory, is factorized in matrix elements of light-cone bilocal operators. A Fourier transformation of diagonal matrix elements results in the conventional quark and gluon densities $`q(x)`$ and $`g(x)`$. Analogously, the off-forward parton distributions are defined by nonforward matrix elements: $`P^{},S^{}\left\|\overline{\psi }_q(\frac{n}{2})n/\text{ }𝒢\psi _q(\frac{n}{2})\right\|P,S\|_{n^2=0}`$ $`=\overline{U}(P^{},S^{})n/\text{ }U(P,S){\displaystyle _1^{+1}}e^{ix(n\overline{P})}H_q(x,\xi ,t)\mathrm{d}x+𝒪(\mathrm{\Delta }),`$ (1a) $`P^{},S^{}\left\|F_{\mu \lambda }^a(\frac{n}{2})n^\mu n^\nu 𝒢_{ab}F^{b\lambda }{}_{\nu }{}^{}(\frac{n}{2})\right\|P,S\|_{n^2=0}`$ $`={\displaystyle \frac{1}{2}}\overline{U}(P^{},S^{})n/\text{ }U(P,S)\left(n\overline{P}\right){\displaystyle _1^{+1}}e^{ix(n\overline{P})}H_g(x,\xi ,t)\mathrm{d}x+𝒪(\mathrm{\Delta }),`$ (1b) where $`𝒢_{(ab)}`$ is the Wilson gauge link. Note that we have an additional factor of $`x`$ in the definition of the gluon distribution compared to the original definition of Ji , $`H_g=xH_g^{\text{Ji}}`$, which removes an “artificial” singularity for finite $`\xi `$ . Due to time reversal invariance and hermiticity, Ji’s OFPDs are even functions of $`\xi `$, so it is sufficient to treat only positive values of $`\xi `$. The different $`𝒪(\mathrm{\Delta })`$ contributions can be found, for example, in Refs. . For vanishing $`\mathrm{\Delta }`$ the off-forward parton distributions reduce to the diagonal partons: $`H_q(x,0,0)`$ $`=𝑞(x),`$ (2a) $`H_g(x,0,0)`$ $`=x𝑔(x).`$ (2b) The renormalization of the defining operators leads to a scale dependence of the off-forward parton distributions. The evolution of the OFPDs takes a simple form at the one-loop level for the Gegenbauer moments, $`G_n^q(\xi ,t,\mu )`$ $`:={\displaystyle \frac{2^n[n!]^2}{(2n+1)!}}{\displaystyle _1^{+1}}\xi ^nC_n^{(3/2)}(\frac{x}{\xi })H_q(x,\xi ,t,\mu )\mathrm{d}x,`$ (3a) $`G_n^g(\xi ,t,\mu )`$ $`:={\displaystyle \frac{32^n(n1)!n!}{(2n+1)!}}{\displaystyle _1^{+1}}\xi ^{n1}C_{n1}^{(5/2)}(\frac{x}{\xi })H_g(x,t,\mu ,\xi )\mathrm{d}x,`$ (3b) since they evolve exactly as the Mellin moments in the diagonal case . For example, for the Gegenbauer moments of the nonsinglet off-forward quark distributions one has $$G_n^{q,ns}(\xi ,t,\mu )=\left(\frac{\alpha _s(\mu )}{\alpha _s(\mu _0)}\right)^{\gamma _{0n}/2\beta _0}G_n^{q,ns}(\xi ,t,\mu _0),$$ (4) where $`\gamma _{0n}`$ and $`\beta _0`$ are the leading coefficients of the nonsinglet anomalous dimension and the beta function. This fact allows the definition of effective forward parton distributions, whose Mellin moments equal the corresponding Gegenbauer moments of the OFPDs : $`{\displaystyle _1^{+1}}x^nq_{\xi ,t}(x,\mu )\mathrm{d}x`$ $`=G_n^q(\xi ,t,\mu ),`$ (5a) $`{\displaystyle _1^{+1}}x^ng_{\xi ,t}(x,\mu )\mathrm{d}x`$ $`=G_n^g(\xi ,t,\mu ).`$ (5b) Their scale dependence is governed by the conventional evolution equations, and they reduce to the diagonal quark and gluon densities for $`\xi ,t0`$. The crucial point is that the effective and off-forward parton distributions can be related to each other by Shuvaev’s transformations . As these transformations do not depend on the momentum transfer invariant $`t`$ and the scale $`\mu `$, we can safely skip them in the following. ## 3 Shuvaev’s inverse transformation We start with connecting the effective forward parton distribution to the off-forward ones. This can be done by Shuvaev’s inverse integral transformation : $`q_\xi (x)`$ $`={\displaystyle _1^{+1}}𝒦_q^1(x,\xi ;y)H_q(y,\xi )\mathrm{d}y,`$ (6a) $`g_\xi (x)`$ $`={\displaystyle _1^{+1}}𝒦_g^1(x,\xi ;y)H_g(y,\xi )\mathrm{d}y.`$ (6b) We briefly sketch the main steps of the derivation of the integral kernels $`𝒦_{q,g}^1(x,\xi ;y)`$ in Ref. , in order to determine the support properties of the EFPDs. The calculation is based on the formal inversion of the Mellin moments in Eq. (5) $`q_\xi (x)`$ $`={\displaystyle \frac{1}{\pi }}\mathrm{disc}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{G_n^q(\xi )}{x^{n+1}}},`$ (7a) $`g_\xi (x)`$ $`={\displaystyle \frac{1}{\pi }}\mathrm{disc}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{G_n^g(\xi )}{x^{n+1}}},`$ (7b) with $`\mathrm{disc}F(x)`$ $`={\displaystyle \frac{1}{2i}}\underset{\epsilon 0}{lim}\left[F(x+i\epsilon )F(xi\epsilon )\right].`$ The Gegenbauer moments $`G_n^{q,g}(x)`$ are defined in Eq. (3). The factorial functions in Eq. (3) are replaced using the integral representation of the beta function {Eqs. (6.1.18) and (6.2.1) in Ref. }. We obtain $`q_\xi (x)`$ $`={\displaystyle _1^{+1}}\left[{\displaystyle \frac{1}{\pi }}\mathrm{disc}{\displaystyle _1^{+1}}{\displaystyle \frac{1}{2x\sqrt{1s}}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}C_n^{(3/2)}(\frac{y}{\xi })\left({\displaystyle \frac{s\xi }{2x}}\right)^n\mathrm{d}s\right]H_q(y,\xi )\mathrm{d}y,`$ (8a) $`g_\xi (x)`$ $`={\displaystyle _1^{+1}}\left[{\displaystyle \frac{1}{\pi }}\mathrm{disc}{\displaystyle _1^{+1}}{\displaystyle \frac{3\sqrt{1s}}{2x^2}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}C_n^{(5/2)}(\frac{y}{\xi })\left({\displaystyle \frac{s\xi }{2x}}\right)^n\mathrm{d}s\right]H_g(y,\xi )\mathrm{d}y.`$ (8b) The expressions in square brackets are the integral kernels $`𝒦_{q,g}^1(x,\xi ;y)`$. Before we state their final form, we have to look at the generating functions of the Gegenbauer polynomials {Eq. (22.9.3) in Ref. }: $$\mathrm{exp}\left(\nu \mathrm{Log}(12wz+z^2)\right)=\underset{n=0}{\overset{\mathrm{}}{}}C_n^{(\nu )}(w)z^n.$$ (9) The generating functions on the left-hand side analytically continue the power series on the right-hand side to the complete complex plane. In Fig. 1 we show the circles of convergence of the power series and the discontinuities of the generating functions, which arise from negative arguments of the complex logarithm. We see that we have to distinguish two cases: $`\|y\|>\xi `$, which corresponds to the parton-distribution-like region of the OFPDs, and $`\|y\|<\xi `$, which corresponds to the meson-wave-function-like region. Let us begin with the latter case. One might think that for $`\|y\|<\xi `$ one has no contribution to the integral kernels $`𝒦_{q,g}^1(x,\xi ;y)`$ because the generating functions are analytical for any finite and real $`z`$. However, the discontinuity at $`z=\mathrm{}`$ produces delta functions and their derivatives $`\delta ^{(n)}(x)`$. From the circle of convergence we find that their contribution is in any case restricted to $`\|x\|<\xi /2`$ in the effective forward parton distributions. For $`\|y\|>\xi `$ the discontinuity can be easily calculated. We face strong singularities in the generating functions at the end points of the cut, therefore, we have to take derivatives of less singular functions, so that the $`s`$-integral in Eq. (8) is convergent. An examination of the circle of convergence shows that $$q_\xi (x)=0\text{ and }g_\xi (x)=0,\text{if }\|x\|>x_b:=\frac{1}{2}\left(1+\sqrt{1\xi ^2}\right).$$ (10) This defines the support area of the EFPDs, which is shown in Fig. 2. Finally, the complete result for the integral kernels of Shuvaev’s inverse transformation is $`𝒦_q^1(x,\xi ;y)`$ $`={\displaystyle \frac{2x}{\pi \|\xi \|}}{\displaystyle \frac{}{y}}{\displaystyle _0^1}{\displaystyle \frac{𝜃(s^2+s\frac{4xy}{\xi ^2}\frac{4x^2}{\xi ^2})}{s\sqrt{(1s)(s^2+s\frac{4xy}{\xi ^2}\frac{4x^2}{\xi ^2})}}}\mathrm{d}s`$ $`+𝜃(1\|\frac{y}{\xi }\|){\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}(1)^n{\displaystyle \frac{2^n[n!]^2}{(2n+1)!}}\xi ^nC_n^{(3/2)}(\frac{y}{\xi })\delta ^{(n)}(x),`$ (11a) $`𝒦_g^1(x,\xi ;y)`$ $`={\displaystyle \frac{4x}{\pi \|\xi \|}}{\displaystyle \frac{^2}{y^2}}{\displaystyle _0^1}{\displaystyle \frac{𝜃(s^2+s\frac{4xy}{\xi ^2}\frac{4x^2}{\xi ^2})\sqrt{1s}}{s^2\sqrt{s^2+s\frac{4xy}{\xi ^2}\frac{4x^2}{\xi ^2}}}}\mathrm{d}s`$ $`+𝜃(1\|\frac{y}{\xi }\|){\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}(1)^n{\displaystyle \frac{32^n(n1)!n!}{(2n+1)!}}\xi ^{n1}C_{n1}^{(5/2)}(\frac{y}{\xi })\delta ^{(n)}(x).`$ (11b) We cannot carry out the derivatives, which would lead to infinite contributions from the derivatives of the theta functions and divergent integrals. The infinite sums represent the contribution of the meson-wave-function-like part of the off-forward parton distributions and do not appear in Ref. , since the discontinuity at $`z=\mathrm{}`$, resp. $`x=0`$, was overlooked. The occurence of delta functions and their derivatives in the effective forward parton distributions are comparable to meson-exchange-type contributions in double distributions . The expressions for the integral kernels $`𝒦_{q,g}^1(x,\xi ;y)`$ show that the Gegenbauer moment inversion is not practicable, in general. Therefore, it can generally not be used for a simple solution of the evolution equations of off-forward parton distributions. ## 4 Shuvaev’s transformation As we will later see, the predictive power of the formalism lies in relating the effective forward parton distributions to the off-forward ones by Shuvaev’s integral transformation $`H_q(x,\xi )`$ $`={\displaystyle _1^{+1}}𝒦_q(x,\xi ;y)q_\xi (y)\mathrm{d}y,`$ (12a) $`H_g(x,\xi )`$ $`={\displaystyle _1^{+1}}𝒦_g(x,\xi ;y)g_\xi (y)\mathrm{d}y.`$ (12b) The full derivation of the integral kernels $`𝒦_{q,g}(x,\xi ;y)`$ can be found in Ref. . We merely state the finite result: $`𝒦_q(x,\xi ;y)`$ $`={\displaystyle \frac{1}{\pi \sqrt{\|y\|}}}{\displaystyle \frac{}{y}}{\displaystyle \frac{y}{\sqrt{\|y\|}}}{\displaystyle _1^{+1}}𝜃(\frac{y(1s^2)}{x\xi s}1)\sqrt{{\displaystyle \frac{x\xi s}{y(1s^2)x+\xi s}}}\mathrm{d}s,`$ (13a) $`𝒦_g(x,\xi ;y)`$ $`={\displaystyle \frac{1}{\pi \sqrt{\|y\|}}}{\displaystyle \frac{}{y}}\sqrt{\|y\|}{\displaystyle _1^{+1}}𝜃(\frac{y(1s^2)}{x\xi s}1)\sqrt{{\displaystyle \frac{(x\xi s)^3}{y(1s^2)x+\xi s}}}\mathrm{d}s.`$ (13b) Again, performing the derivatives would give divergent integrals and infinite contributions from the end points. The derivatives of the theta function give rise to the “suspicious overall sign” that is mentioned in Ref. . Equation (13) is equivalent to previous results presented in Refs. . It is useful to express the integral kernels in terms of standard elliptic integrals, because it is then possible to perform the derivatives analytically. First, we give the symmetry properties of the integral kernels: $`𝒦_q(x,\xi ;y)`$ $`=𝒦_q(x,\xi ;y),`$ $`𝒦_q(x,\xi ;y)`$ $`=+𝒦_q(x,\xi ;y),`$ (14a) $`𝒦_g(x,\xi ;y)`$ $`=𝒦_g(x,\xi ;y),`$ $`𝒦_g(x,\xi ;y)`$ $`=𝒦_g(x,\xi ;y).`$ (14b) Of course, $`𝒦_{q,g}(x,\xi ;y)`$ obey the fundamental $`(\xi \xi )`$-symmetry of Ji’s off-forward parton distributions. The relations on the right-hand side show that it is sufficient to restrict the calculation to positive values of $`y`$. We define $$a:=\frac{x}{\xi }\frac{\xi }{2y}\frac{1}{2y}\sqrt{4y^24yx+\xi ^2},b:=\frac{x}{\xi }\frac{\xi }{2y}+\frac{1}{2y}\sqrt{4y^24yx+\xi ^2},$$ (15) which correspond up to the $`x/\xi `$, which we added for convenience, to the zeroes of the denominator in the square root in Eq. (13). As the radicand has to be positive, we must further restrict the possible range of $`y`$ to $$y>x_a:=\frac{1}{2}(x+\sqrt{x^2\xi ^2}),\text{if }x\xi .$$ (16) With help of the integral tables in Ref. we obtain our final result: $`𝒦_q(x,\xi ;y)`$ $`=\{\begin{array}{cc}𝛿(yx_a)\sqrt{\frac{1}{x_a}\sqrt{x^2\xi ^2}}\hfill & \\ +𝜃(yx_a)\frac{\xi }{\pi y^2}\sqrt{\frac{\xi }{yb}}\frac{b}{ba}\left(R_F(0,\frac{a}{b},1)\frac{1}{3}\frac{b+a}{b}R_D(0,\frac{a}{b},1)\right)\hfill & \text{for }x\xi ,\hfill \\ 𝜃(x+\xi )\frac{\xi }{\pi y^2}\sqrt{\frac{\xi }{y(ba)}}\frac{b}{ba}\left(R_F(0,\frac{a}{ba},1)\frac{1}{3}\frac{b+a}{ba}R_D(0,\frac{a}{ba},1)\right)\hfill & \text{for }x<\xi ,\hfill \end{array}`$ (17a) $`𝒦_g(x,\xi ;y)`$ $`=\{\begin{array}{cc}𝛿(yx_a)\sqrt{\frac{1}{x_a}(x^2\xi ^2)^{3/2}}+𝜃(yx_a)\frac{\xi ^2}{\pi y^2}\sqrt{\frac{\xi b}{y}}\frac{b}{ba}\hfill & \\ \times \left(\frac{2ba}{b}R_F(0,\frac{a}{b},1)\frac{2}{3}\frac{b^2ab+a^2}{b^2}R_D(0,\frac{a}{b},1)\right)\hfill & \text{for }x\xi ,\hfill \\ 𝜃(x+\xi )\frac{\xi ^2}{\pi y^2}\sqrt{\frac{\xi (ba)}{y}}\frac{b}{ba}\hfill & \\ \times \left(\frac{2ba}{ba}R_F(0,\frac{a}{ba},1)\frac{2}{3}\frac{b^2ab+a^2}{(ba)^2}R_D(0,\frac{a}{ba},1)\right)\hfill & \text{for }x<\xi ,\hfill \end{array}`$ (17b) where $`R_F`$ and $`R_D`$ are Carlson’s elliptic integrals of the first and second kind (see, e.g., ) with $`R_F(x,y,z)`$ $`:={\displaystyle \frac{1}{2}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{1}{\sqrt{(t+x)(t+y)(t+z)}}}\mathrm{d}t,`$ (18) $`R_D(x,y,z)`$ $`:={\displaystyle \frac{3}{2}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{1}{(t+z)\sqrt{(t+x)(t+y)(t+z)}}}\mathrm{d}t.`$ (19) These integral kernels accumulate the main properties of the $`x`$ and $`\xi `$ dependence of off-forward parton distributions. For $`x\xi `$, the OFPDs essentially look like forward quarks and gluons: $`𝒦_q(x,\xi ;y)`$ $`=𝛿(yx)+{\displaystyle \frac{1}{x}}𝒪(\frac{\xi ^2}{x^2}),`$ (20a) $`𝒦_g(x,\xi ;y)`$ $`=x𝛿(yx)+𝒪(\frac{\xi ^2}{x^2}).`$ (20b) The forward evolution concentrates the effective forward parton distributions at $`y0`$. Therefore, the small-$`y`$ behavior of the integral kernels $`𝒦_{q,g}(x,\xi ;y)`$ reproduces the well-known asymptotic forms of the off-forward valence and singlet quark, and gluon distributions : $`\begin{array}{cc}\hfill H_q^v(x,\xi )& =𝜃(1\|\frac{x}{\xi }\|){\displaystyle \frac{\xi ^2x^2}{\xi ^3}}{\displaystyle _0^\epsilon }\left({\displaystyle \frac{3}{2}}+𝒪(\frac{y}{\xi })\right)q_\xi ^v(y)\mathrm{d}y,\hfill \\ \hfill H_q^s(x,\xi )& =𝜃(1\|\frac{x}{\xi }\|){\displaystyle \frac{x\left(\xi ^2x^2\right)}{\xi ^5}}{\displaystyle _0^\epsilon }\left({\displaystyle \frac{15}{2}}+𝒪(\frac{y}{\xi })\right)yq_\xi ^s(y)\mathrm{d}y,\hfill \end{array}`$ (21a) $`H_g(x,\xi )`$ $`=𝜃(1\|\frac{x}{\xi }\|){\displaystyle \frac{(\xi ^2x^2)^2}{\xi ^5}}{\displaystyle _0^\epsilon }\left({\displaystyle \frac{15}{8}}+𝒪(\frac{y}{\xi })\right)yg_\xi (y)\mathrm{d}y.`$ (21b) Additionally, these equations prove that the integrals in Eq. (12) are well defined and convergent. Even the physical interpretation of the different regions in $`x`$ holds in the formalism of EFPDs. Inserting Eq. (14) in Eq. (12) yields $`H_q(x,\xi )`$ $`={\displaystyle _0^1}\left(𝒦_q(x,\xi ;y)q_\xi (y)+𝒦_q(x,\xi ;y)q_\xi (y)\right)\mathrm{d}y,`$ (22a) $`H_g(x,\xi )`$ $`={\displaystyle _0^1}\left(𝒦_g(x,\xi ;y)+𝒦_g(x,\xi ;y)\right)g_\xi (y)\mathrm{d}y.`$ (22b) For $`\|x\|>\xi `$ the off-forward parton distributions are related to corresponding effective forward (anti)partons with a minimum momentum $`\|x_a\|`$. For $`\|x\|<\xi `$ the picture of a meson wave function is supported by a simultaneous contribution of effective forward partons and antipartons with any momentum $`y`$. The different expressions for the integral kernels $`𝒦_{q,g}(x,\xi ;y)`$ for $`\|x\|\xi `$ show, analogous to , that the off-forward parton distributions are not analytic at $`\|x\|=\xi `$. The analyticity of the OFPDs for $`x\xi `$ requires that the effective forward parton distributions need to be analytic for $`\|x\|\xi /2`$ only. ## 5 Nonforward amplitudes and EFPDs The use of off-forward parton distributions is required in deeply virtual Compton scattering (DVCS) and hard exclusive electroproduction processes. Detailed information can be found in Ref. , and references therein. Here, we are only interested in the part of the amplitudes that refers to the OFPDs: $`𝒜_q(\xi )`$ $`:={\displaystyle _1^{+1}}\left({\displaystyle \frac{1}{x\xi +i\epsilon }}+{\displaystyle \frac{1}{x+\xi i\epsilon }}\right)H_q(x,\xi )\mathrm{d}x,`$ (23a) $`𝒜_g(\xi )`$ $`:={\displaystyle _1^{+1}}\left({\displaystyle \frac{1}{x\xi +i\epsilon }}+{\displaystyle \frac{1}{x+\xi i\epsilon }}\right){\displaystyle \frac{1}{x}}H_g(x,\xi )\mathrm{d}x,`$ (23b) where we have neglected the $`𝒪(\mathrm{\Delta })`$ contributions, analogous to Ref. . The imaginary part of the amplitudes is related to the diagonal elements $`H_{q,g}(\xi ,\xi )`$, which can be expressed by the effective forward parton distributions: $`\pi H_q^s(\xi ,\xi )`$ $`=\mathrm{Im}𝒜_q(\xi )`$ $`={\displaystyle _0^{\sqrt{1\xi /2x_b}}}4\sqrt{1z^2}q_\xi ^s(\frac{\xi }{2(1z^2)})\mathrm{d}z,`$ (24a) $`{\displaystyle \frac{2\pi }{\xi }}H_g(\xi ,\xi )`$ $`=\mathrm{Im}𝒜_g(\xi )`$ $`={\displaystyle _0^{\sqrt{1\xi /2x_b}}}32z^2\sqrt{1z^2}g_\xi (\frac{\xi }{2(1z^2)})\mathrm{d}z,`$ (24b) with the quark singlet $`q_\xi ^s(x)=q_\xi (x)q_\xi (x)`$. The imaginary part is essentially dominated by the behavior of the EFPDs around $`x\xi /2`$. For small values of $`\xi `$ this region can be accurately described by $`xq_\xi ^s(x)x^{\lambda _q},`$ (25a) $`xg_\xi (x)x^{\lambda _g}.`$ (25b) If we insert Eq. (25) into Eq. (24) and set the upper integration limits to one, we achieve approximation formulas for the imaginary part of the amplitudes in Eq. (23): $`R_q^{\mathrm{Im}}:={\displaystyle \frac{\mathrm{Im}𝒜_q(\xi )}{q_\xi ^s(\frac{\xi }{2})}}`$ $`{\displaystyle \frac{2\sqrt{\pi }\Gamma (\lambda _q+\frac{5}{2})}{\Gamma (\lambda _q+3)}},`$ (26a) $`R_g^{\mathrm{Im}}:={\displaystyle \frac{\mathrm{Im}𝒜_g(\xi )}{g_\xi (\frac{\xi }{2})}}`$ $`{\displaystyle \frac{8\sqrt{\pi }\Gamma (\lambda _g+\frac{5}{2})}{\Gamma (\lambda _g+4)}}.`$ (26b) A similar ratio was already presented in Ref. , where, however, the imaginary part was compared to diagonal partons at $`x=2\xi `$, which leads to an extra factor $`2^{2+2\lambda _{q,g}}`$. In Fig. 3 we show a comparison of the exact ratio, derived from Eq. (24), to the approximation in Eq. (26) for the effective distributions in Eq. (25). The change of the integration limit has no remarkable effect up to $`\xi 0.1`$. The accuracy of the gluon ratio is slightly worse compared to the quark ratio, because of an additional factor of $`z^2`$ in the integrand in Eq. (24). Since the quotients of the gamma functions in Eq. (26) have a weak $`\lambda _{q,g}`$ dependency and Fig. 3 shows a good stability under a change of $`\lambda _{q,g}`$, we can conclude that Eq. (26) is an excellent approximation of the imaginary part of the amplitudes for the values of $`\xi `$, where the effective forward parton distributions can be reliably identified with the conventional forward quark and gluon densities. The calculation of the real part is straightforward but tedious \[we have used Eq. (13) rather than Eq. (17)\]. The principal value integration can be performed exactly and the final result consists of integrals without any strong singularities: $`\mathrm{Re}𝒜_q(\xi )`$ $`={\displaystyle _0^1}{\displaystyle \frac{2}{z^2}}\left(z+{\displaystyle \frac{1}{\sqrt{1+z}}}{\displaystyle \frac{1}{\sqrt{1z}}}\right)q_\xi ^s(\frac{\xi z}{2})\mathrm{d}z`$ $`+{\displaystyle _{\xi /2x_b}^1}2\left({\displaystyle \frac{1}{z}}+\sqrt{{\displaystyle \frac{z}{1+z}}}\right)q_\xi ^s(\frac{\xi }{2z})\mathrm{d}z,`$ (27a) $`\mathrm{Re}𝒜_g(\xi )`$ $`={\displaystyle _0^1}{\displaystyle \frac{4}{z^3}}\left(z^28+4\sqrt{1+z}+4\sqrt{1z}\right)g_\xi (\frac{\xi z}{2})\mathrm{d}z`$ $`+{\displaystyle _{\xi /2x_b}^1}4\left({\displaystyle \frac{1}{z}}8z+4\sqrt{z(1+z)}\right)g_\xi (\frac{\xi }{2z})\mathrm{d}z.`$ (27b) We note that the expressions in brackets in the first integrals in Eq. (27) are always negative, therefore, the two integrals partly cancel each other. Again, we insert the small-$`x`$ behavior of Eq. (25) into these integrals and set the lower limits of the second integrals to zero, so that everything can be evaluated and yields the following ratios between the real and imaginary parts: $$R_{q,g}^{\mathrm{Re}}:=\frac{\mathrm{Re}𝒜_{q,g}(\xi )}{\mathrm{Im}𝒜_{q,g}(\xi )}\mathrm{tan}\frac{\pi \lambda _{q,g}}{2}.$$ (28) This is identical to the result, achieved by dispersion relations, in Refs. . From Fig. 4 we see that the quality of the approximations of the real parts is significantly more sensitive — note the logarithmic scale — to the change of the integration limit and to a variation of $`\lambda _{q,g}`$. Additionally, the right-hand side of Eq. (28) depends strongly on $`\lambda _{q,g}`$, and the first integrals in Eq. (27) have dominant contributions from two regions: around $`x0`$ and $`x\xi /2`$, i.e., we must require that $`\lambda _{q,g}`$ is essentially constant for small $`x`$. Therefore, only for very small $`\xi `$, when Eq. (25) is a valid approximation for a large range of $`x`$ for the usual forward quark and gluon distribution, the latter can be used to reliably predict the real part of the amplitude. Nevertheless, for small values of $`\lambda _{q,g}`$, where the real part is strongly suppressed, the absolute values of the amplitudes are determined to a good precision for small $`\xi `$ . ## 6 Moment-diagonal models In this section, we try to build a model for the effective forward parton distributions. The situation does not seem to be very promising, because Shuvaev’s inverse transformation cannot generally be used and we face a difficult support area for the EFPDs in Fig. 2. References were dealing with a model for off-forward parton distributions, the Gegenbauer moments of which are independent of $`\xi `$. Because the corresponding EFPDs are also independent of $`\xi `$, these models manifestly violate the support area in Fig. 2, as was also recognized in Ref. , though it should be a good approximation for small values of $`\xi `$. This model would have had the great advantage that it would have been stable against a change of the input scale, since the evolution of the Gegenbauer moments is identical to that of the Mellin moments. Nevertheless, this idea can be used to find valid models for the effective forward parton distributions. Because a common $`n`$-, $`\xi `$-, or $`t`$-dependent factor does not have any influence on the evolution equations, we can generalize the model with $`\xi `$-independent Gegenbauer moments to a class of moment-diagonal models with a common proportionality factor: $$G_n^{q,g}(\xi ,t,\mu ):=\text{const}(n,\xi ,t)\times M_n^{q,g}(\mu ).$$ (29) The $`t`$ dependence is usually factorized, $`\text{const}(n,\xi )\times F_1(t)`$, where $`F_1(t)`$ is the Dirac form factor . These models allow, provided that the Mellin inverse can be performed and leads to valid supports, a direct and simple calculation of off-forward parton distributions and amplitudes, with help of the formulas (17), (24), and (27), for arbitrary $`x`$, $`\xi `$, $`\mu `$, and for the region of $`t`$, where the $`𝒪(\mathrm{\Delta })`$ contributions can be neglected. The simplest model for proton OFPDs of this class, which fulfills all known theoretical constraints for $`t=0`$ (e.g., see Ref. ) and gives a good approximation of the $`t`$-dependence, is $$G_n^{q,g}(\xi ,t,\mu ):=2\left(\frac{\xi }{2}\right)^{n+1}T_{n+1}(\xi ^1)F_1^p(t)M_n^{q,g}(\mu ),$$ (30) where $`M_n^{q,g}(\mu )`$ are the Mellin moments of the usual quark and gluon distributions in the proton, $`F_1^p(t)`$ is the Dirac form factor of the proton with $`F_1^p(0)=1`$, and $`T_n(x)`$ are Chebyshev polynomials of the first kind {Eq. (22.2.4) in Ref. }. It is advantageous to use the Glück–Reya–Vogt 1998 (GRV 98) parton distributions , since they are given for $`x`$ values down to $`10^9`$, which allows us to compute the real part in Eq. (27) for small values of $`\xi `$ with a high accuracy. With use of Eq. (22.3.25) and (4.4.27) in Ref. , a Mellin inversion yields $`q_\xi (x)`$ $`=\left\{𝜃(\frac{1+\sqrt{1\xi ^2}}{2}\|x\|)𝑞(\frac{2x}{1+\sqrt{1\xi ^2}})+𝜃(\frac{1\sqrt{1\xi ^2}}{2}\|x\|)𝑞(\frac{2x}{1\sqrt{1\xi ^2}})\right\}F_1^p(t),`$ (31a) $`g_\xi (x)`$ $`=\left\{𝜃(\frac{1+\sqrt{1\xi ^2}}{2}\|x\|)𝑔(\frac{2x}{1+\sqrt{1\xi ^2}})+𝜃(\frac{1\sqrt{1\xi ^2}}{2}\|x\|)𝑔(\frac{2x}{1\sqrt{1\xi ^2}})\right\}F_1^p(t).`$ (31b) We see that the effective forward parton distributions are a simple combination of two rescaled forward parton densities with the correct support area and an appropriate common $`t`$-dependent factor. The first one gives the conventional forward quark and gluons for vanishing $`\xi `$ and $`t`$. The contribution of the second summand is restricted to $`\|x\|<\xi /2`$, i.e., it influences only the meson-wave-function-like region $`\|x\|<\xi `$ of OFPDs, and is negligible for small values of $`\xi `$. Therefore, most of the numerical results of Refs. can be accurately transferred to the model in Eq. (31). It is an interesting fact that the argument of the first forward parton density in Eq. (31) is very similar to the Georgi–Politzer $`\xi `$-scaling variable , $`\xi =2x_B/(1+\sqrt{1+4x_B^2M_N^2/Q^2})`$, that originally described target mass effects in deep inelastic scattering. Hence, one can argue that the arguments of the parton densities in Eq. (31) reflect skewedness effects. But such arguments can as well be relics of the Gegenbauer polynomials that appear in the derivation of the $`\xi `$-scaling variable. ## 7 Summary and conclusions In this paper, we presented with Eqs. (17), (24), and (27) simple expressions that relate effective forward parton distributions, which evolve like conventional forward partons, to off-forward parton distributions and nonforward amplitudes. We emphasized that the off-forward parton distributions and nonforward amplitudes can be directly determined from the conventional forward parton distributions and nucleon form factors at arbitrary scale $`\mu `$ for moment-diagonal models. Exemplary, we stated a simple self-consistent model for the EFPDs of the proton in terms of the GRV 98 parton distributions and the Dirac form factor of the proton, which allows us to predict off-forward parton distributions and nonforward amplitudes for arbitrary $`x`$, $`\xi `$, $`\mu `$, and (not to large) $`t`$. These predictions should not differ too much from results of other models at least at small $`\xi `$, as the results in Refs. show. Nevertheless, it would be illuminating if further self-consistent moment-diagonal models exist, especially models that have a qualitatively different behavior in the meson-wave-function-like region, such as the off-forward parton distributions of chiral soliton model calculations , since the real part of nonforward amplitudes is dominated by this region and gets important for large $`\xi `$. Because of the complicated support area of EFPDs, an investigation of the double distributions of moment-diagonal models might be helpful. ## Acknowledgments We thank M. Glück and E. Reya for proposing this investigation and for instructive remarks during the initial stage of this work.
warning/0004/cond-mat0004013.html	ar5iv	text	# The effect of rare regions on a disordered itinerant quantum antiferromagnet with cubic anisotropy ## I Introduction Quenched disorder can have very drastic influences on the critical behavior of a system undergoing a continuous phase transition. According to the Harris criterion the critical behavior of a clean system is unaltered by disorder, if the correlation length critical exponent $`\nu `$ obeys the inequality $`\nu >2/d`$, where $`d`$ is the spatial dimensionality of the system. In the opposite case, $`\nu <2/d`$, the clean critical behavior is unstable, and the disorder either leads to a new, different universality class, or to an unconventional critical point, or even to the destruction of the phase transition. Another, less well understood consequence of quenched disorder is the formation of rare locally ordered regions in the disordered phase. For a transition occuring at a finite temperature, this can be explained in the following way. In general, disorder leads to the suppression of the critical temperature from its clean value $`T_c^0`$ to $`T_c`$. In the temperature region between $`T_c^0`$ and $`T_c`$ the system does not show long-range order. However, there will be arbitrarily large regions which are devoid of impurities and thus order locally. The probability of finding such regions usually decreases exponentially with their size, they represent non-perturbative degrees of freedom. These locally ordered regions are known as rare regions, and the order parameter fluctuations induced by them as local moments or instantons. Griffiths showed that the rare regions lead to a non-analytic free energy everywhere in the temperature region between $`T_c^0`$ and $`T_c`$, now called the Griffiths region or Griffiths phase. In generic classical systems this is a very weak effect, and the non-analyticity in the free energy is only an essential one. However, the Griffiths singularities become stronger if the disorder is spatially correlated. McCoy and Wu studied a two-dimensional Ising model where the disorder is perfectly correlated in one spatial direction and uncorrelated in the other. In this model the rare regions lead to the divergence of the susceptibility at some temperature $`T_\chi `$ within the Griffiths region. A very interesting question is what is the influence of the rare regions on the critical behavior of a system. Dotsenko et al. studied this question for a weakly disordered classical ferromagnet. They found that the conventional theory of critical behavior in this system is unstable with respect to replica symmetry breaking. They also showed that the rare regions actually induce replica symmetry breaking perturbations and thus destabilize the conventional critical fixed point. While so far no final conclusion about the fate of the transition in the weakly disordered ferromagnet could be reached, the occurrence of replica symmetry breaking raises the possibility of an unconventional transition with activated scaling, as is believed to occur in the random field Ising model . For quantum phase transitions which occur at zero temperature as a function of some non-thermal control parameter, one expects an even stronger influence of the rare regions than for classical transitions. The reason is that a quantum model with uncorrelated quenched disorder is effectively equivalent to a classical model with the disorder being perfectly correlated in one dimension (the imaginary time dimension). Fisher investigated the critical behavior of a one-dimensional quantum Ising spin chain in a transverse field which is equivalent to the classical McCoy-Wu model. He found that due to the rare regions the critical behavior is of the activated form. This has been confirmed by numerical simulations which also suggest that this sort of behavior may not be restricted to one-dimensional systems. In two recent papers we have considered the effect of rare regions on quantum phase transitions of itinerant electrons in $`d>1`$. We have developed a systematic approach, representing the local moments by inhomogeneous saddle point solutions of the field theory. The interaction between the local moments and the fluctuations leads to a new term in the effective action which is of the form of annealed static disorder. In the case of the quantum antiferromagnetic transition this new term results in the destruction of the conventional critical fixed point if the number $`p`$ of order parameter components is smaller than 4. No new fixed point could be identified, the system displays runaway flow to large disorder strength. On the other hand, for the quantum ferromagnetic transition the rare regions do not affect the critical behavior since a self-induced long-range interaction suppresses all fluctuations including those produced by the local moments. In this paper we apply the approach developed in Ref. to a model of an itinerant antiferromagnet with an additional interaction term with cubic symmetry. This model is equivalent to a weakly disordered classical ferromagnet with cubic anisotropy in which the disorder is perfectly correlated in some of the spatial dimensions but uncorrelated in the remaining dimensions. The conventional theory for this model (without taking rare regions into account) has been developed by Yamazaki, Holz, Ochiai and Fukuda. The purpose for this work is threefold. We want investigate (i) whether the conventional critical fixed point is stable under the influence of the rare regions. If it is unstable we want to find out (ii) whether a new stable fixed fixed point exists which describes a rare region driven transition. Finally we want to study (iii) the influence of the rare regions on the fluctuation-driven first-order transition occurring in our system. The layout of the paper is as follows. In Sec. II we derive the effective field theory by taking into account the disorder induced rare regions. In Sec. III, we carry out the renormalization group analysis. Finally, Sec. IV is left for a summary of our results. ## II An effective action for disordered antiferromagnets with cubic anisotropy ### A The model In 1976 Hertz derived an order parameter field theory for the description of the antiferromagnetic quantum phase transition of itinerant electrons. Later this model was generalized to the dirty case by making the distance from the critical point a random function of position . Here we consider an extension of this order parameter field theory by incorporating an additional $`\varphi ^4`$ term which possesses a (hyper-)cubic symmetry. In terms of the $`p`$-component order parameter field $`\mathit{\varphi }`$ (with components $`\varphi _i`$) the total action can be written as $$S[\mathit{\varphi }]=S_\mathrm{G}[\mathit{\varphi }]+S_{\mathrm{int}}[\mathit{\varphi }]+S_{\mathrm{cubic}}[\mathit{\varphi }],$$ (2) with the Gaussian part, the interaction part and the cubic anisotropic part given by $$S_\mathrm{G}[\mathit{\varphi }]=\frac{1}{2}𝑑x𝑑y\underset{i}{}\varphi _i(x)\mathrm{\Gamma }(xy)\varphi _i(y),$$ (3) $$S_{\mathrm{int}}[\mathit{\varphi }]=u𝑑x\underset{i,j}{}\varphi _i(x)\varphi _i(x)\varphi _j(x)\varphi _j(x),$$ (4) $$S_{\mathrm{cubic}}[\mathit{\varphi }]=\lambda 𝑑x\underset{i}{}\varphi _{i}^{}{}_{}{}^{4}(x).$$ (5) Here we use a 4-vector notation to combine the real space coordinate $`𝐱`$ and imaginary time $`\tau `$, $`x=(𝐱,\tau `$), $`𝑑x=𝑑𝐱_0^{1/T}𝑑\tau `$. The bare two point function, $`\mathrm{\Gamma }(𝐱𝐲,\tau \tau ^{})`$ $`=`$ $`\mathrm{\Gamma }_0(𝐱𝐲,\tau \tau ^{})`$ (7) $`+\delta (𝐱𝐲)\delta (\tau \tau ^{})\delta t(𝐱),`$ consists of the deterministic part derived by Hertz whose Fourier transform reads $$\mathrm{\Gamma }_0(𝐪,\omega _n)=t_0+𝐪^2+\|\omega _n\|,$$ (8) and a disorder part in the form of a ”random mass” term. Here $`𝐪`$ is the wave vector, $`\omega _n`$ is a bosonic Matsubara frequency and $`\delta t(𝐱)`$ is a random function of position and is endowed with the following statistical properties: $$\delta t(𝐱)=0,$$ (10) $$\delta t(𝐱)\delta t(𝐲)=\mathrm{\Delta }\delta (𝐱𝐲).$$ (11) ### B Inhomogeneous saddle points and annealed disorder In the conventional approach to critical behavior in systems with quenched disorder the disorder average is carried out at the beginning of the calculation by means of the replica trick . A subsequent perturbative analysis of the resulting, spatially homogeneous effective theory misses the rare regions we are interested in since they are non-perturbative degrees of freedom. We therefore follow the approach developed in Ref. , and work with a particular realization of the disorder rather than integrating it out. Let us consider spatially inhomogeneous, but time-independent saddle point solutions of the action (II A) (time-dependent saddle-point solutions – if any – will always have a higher free energy since the disorder is static). Depending on the sign of the cubic interaction term the structure of the saddle points in the $`p`$-dimensional order parameter space will be different. When $`\lambda >0`$ the free energy is minimized by saddle point solutions that lie on the diagonals of a $`p`$-dimensional hypercube, while when $`\lambda <0`$ the free energy is minimized by solutions that lie on the axis of the hypercube. In either case the modulus $`\varphi _{\mathrm{sp}}`$ of these minimizing saddle point solutions fulfills the equation $$\left(t_0+\delta t(𝐱)_𝐱^\mathrm{𝟐}\right)\|\mathit{\varphi }_{\mathrm{sp}}(𝐱)\|+4u_{\mathrm{eff}}\|\mathit{\varphi }_{\mathrm{sp}}(𝐱)\|^3=0,$$ (13) $$u_{\mathrm{eff}}=\{\begin{array}{cc}u+\frac{\lambda }{p}\hfill & \text{for }\lambda >0\hfill \\ u+\lambda \hfill & \text{for }\lambda <0\hfill \end{array}.$$ (14) Although $`\varphi _{\mathrm{sp}}(𝐱)=0`$ is always a solution, there will be spatially inhomogeneous solutions if $`\delta t(𝐱)`$ has sufficiently deep and wide troughs . Let us now consider the Griffiths region, i.e. the region where the average distance $`t_0`$ from the critical point is positive but where there are isolated islands which support a non-zero $`\varphi _{\mathrm{sp}}`$. If we have $`N`$ such islands which are sufficiently apart from each other the global saddle point solutions may be written as $`\mathit{\varphi }_{\mathrm{sp}}^{\{\sigma _I\}}(𝐱)\mathrm{\Phi }^{\{\sigma _I\}}(𝐱)`$ $`=`$ $`{\displaystyle \underset{I=1}{\overset{N}{}}}\psi _I(𝐱)\sigma _I`$ (15) where $`\psi _I(𝐱)`$ is a solution of (II B) on the island $`I`$ and $`\sigma _I`$ is a unit vector in spin space (on one of the axis for $`\lambda <0`$ or on one of the diagonals for $`\lambda >0`$). Since the direction of the order parameter on each of the $`N`$ islands can be chosen independently, (15) describes an exponentially large number of degenerate saddle points, $`(2p)^N`$ for $`\lambda <0`$ and $`(2^p)^N`$ for $`\lambda >0`$. To be precise, the saddle points are not exactly degenerate due to the residual interaction of the (exponentially small) tails of the order parameter between the islands. The complicated structure of the free energy landscape connected with the existence of an exponentially large number of almost degenerate saddle points will finally turn out to be responsible for the failure of the conventional approach. We now consider fluctuations around the saddle points (15). Since the saddle points are separated by large free energy barriers an expansion around one of them will not give a good representation of the partition function of the entire system. Instead we will restrict ourselves to small fluctuations and simply add the contributions coming from all of the saddle points. Thus the partition function for a particular realization $`\delta t(𝐱)`$ of the disorder can be written as $$Z[\delta t(𝐱)]\underset{\{\sigma _I\}}{}_<D[\phi (x)]e^{S[\mathrm{\Phi }^{\{\sigma _I\}}(𝐱)+𝝋(x),\delta t(𝐱)]}.$$ (16) Here $`_<`$ indicates that the integration is restricted to small fluctuations $`𝝋`$ only. We now carry out the sum over the saddle point configurations. The residual interaction between the islands will lead to slight deviations of the saddle point function from the ideal one given in (15). This is taken into account by replacing the sum over the saddle points by an integral over a probability distribution $$P[\mathrm{\Phi }]e^{\frac{1}{T}{\scriptscriptstyle 𝑑x^{\mathrm{sp}}(\mathrm{\Phi })}}.$$ (17) The temperature factor in the exponent reflects the fact that the saddle points are classical (static) degrees of freedom . Expanding in powers of the fluctuations, we obtain the following effective action for the fluctuations $`𝝋`$ (still for a particular disorder realization) $`S_{\mathrm{eff}}S^{\mathrm{SP}}`$ $`=`$ $`S_\mathrm{G}[𝝋]+S_{\mathrm{int}}[𝝋]+S_{\mathrm{cubic}}[𝝋]`$ (18) $`+`$ $`T\overline{w}{\displaystyle 𝑑x𝑑yC(x,y)\underset{i,j}{}\phi _i^2(x)\phi _j^2(y)}`$ (19) $`+`$ $`\mathrm{higher}\mathrm{order}\mathrm{terms}.`$ (20) The correlation function $`C(x,y)`$ measures, up to a constant factor determined by the precise form of $``$, whether $`𝐱`$ and $`𝐲`$ belong to the same island, and $`\overline{w}=[(2+4/p)u+6\lambda /p]`$ is a positive constant. The $`\overline{w}`$ term is produced by the interaction of the fluctuations with the rare regions. It is our approximation of the effect of these non-perturbative degrees of freedom. Terms of higher than fourth order in $`𝝋`$ also arise, but they are renormalization group irrelevant at both the Gaussian and the nontrivial fixed points of the conventional theory (see below). Having identified the effects of the rare regions we now use the replica trick to perform the quenched disorder average over $`\delta t(𝐱)`$ which implies an average over position and size of the rare regions. The resulting effective action reads $`S`$ $`{}_{\mathrm{eff}}{}^{}[𝝋^\alpha (x)]=`$ (21) $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{\alpha }{}}{\displaystyle \underset{i}{}}{\displaystyle 𝑑x𝑑y\mathrm{\Gamma }_0(xy)𝝋_i^\alpha (x)𝝋_i^\alpha (y)}`$ (22) $`+`$ $`u{\displaystyle \underset{\alpha }{}}{\displaystyle \underset{i,j}{}}{\displaystyle 𝑑𝐱𝑑\tau \left(𝝋_i^\alpha (𝐱,\tau )\right)^2\left(𝝋_j^\alpha (𝐱,\tau )\right)^2}`$ (23) $`+`$ $`\lambda {\displaystyle \underset{\alpha }{}}{\displaystyle \underset{i}{}}{\displaystyle 𝑑𝐱𝑑\tau \left(𝝋_i^\alpha (𝐱,\tau )\right)^4}`$ (24) $``$ $`\mathrm{\Delta }{\displaystyle \underset{\alpha ,\beta }{}}{\displaystyle \underset{i,j}{}}{\displaystyle 𝑑𝐱𝑑\tau 𝑑\tau ^{}\left(𝝋_i^\alpha (𝐱,\tau )\right)^2\left(𝝋_j^\beta (𝐱,\tau ^{})\right)^2}`$ (25) $``$ $`T\overline{w}{\displaystyle \underset{\alpha ,\beta }{}}{\displaystyle \underset{i,j}{}}{\displaystyle 𝑑𝐱𝑑\tau 𝑑\tau ^{}\left(𝝋_i^\alpha (𝐱,\tau )\right)^2\left(𝝋_j^\alpha (𝐱,\tau ^{})\right)^2}`$ (26) Here the first four terms are identical to the result of the conventional treatment. The 5th term has the form of static, annealed disorder and represents the interaction of the fluctuations with the rare regions in the Griffiths phase. For more details of this derivation see Ref. . ## III Renormalization Group Analysis ### A Flow equations We first consider the effective action (LABEL:eq:2.10) at tree level. As usual, let us define the scale dimension of a length $`L`$ to be $`[L]=1`$, and that of imaginary time $`\tau `$ to be $`[\tau ]=z`$ with $`z`$ being the dynamical critical exponent. We first analyze the Gaussian fixed point. From the Gaussian part of the action (LABEL:eq:2.10) we see that $`\omega _n`$ scales as $`q^2`$, implying that $`z=2`$. The scale dimension of the field is $`[𝝋]=d/2`$. Power counting for the interaction and disorder terms of the action gives the scale dimensions of $`u,\lambda ,\mathrm{\Delta }`$ and $`\overline{w}`$ as $`[u]=[\lambda ]=[\overline{w}]=2d`$ and $`[\mathrm{\Delta }]=4d`$. Here we have used the fact that in Matsubara formalism the temperature scales like a frequency, $`[T]=z`$. Consequently, $`u,\lambda `$ and $`\overline{w}`$ are irrelevant for $`d>2`$, while $`\mathrm{\Delta }`$ is irrelevant only for $`d>4`$. This implies that in the physical dimension $`d=3`$ the Gaussian fixed point is unstable, and we must carry out a loop expansion of the effective action (LABEL:eq:2.10) close to $`d=4`$. All terms of higher order in $`𝝋`$ that arise in addition to those given in (LABEL:eq:2.10) have negative scale dimensions at and close to $`d=4`$. Thus, they are irrelevant by power counting with respect to both the Gaussian and the conventional non-trivial fixed points. As in the conventional theory we carry out the perturbation theory in $`d=4ϵ`$ spatial dimensions and $`ϵ_\tau `$ time dimensions. In this way the perturbation expansion becomes a double expansion in terms of $`ϵ`$ and $`ϵ_\tau `$. The renormalization group flow equations are obtained by performing a frequency momentum shell RG procedure. To one-loop order, we obtain the following flow equations, $`{\displaystyle \frac{du}{dl}}`$ $`=`$ $`\stackrel{~}{ϵ}u4(p+8)u^2+48u\mathrm{\Delta }24u\lambda ,`$ (29) $`{\displaystyle \frac{d\lambda }{dl}}`$ $`=`$ $`\stackrel{~}{ϵ}\lambda 36\lambda ^2+48\lambda \mathrm{\Delta }48u\lambda ,`$ (30) $`{\displaystyle \frac{d\mathrm{\Delta }}{dl}}`$ $`=`$ $`ϵ\mathrm{\Delta }+32\mathrm{\Delta }^28(p+2)u\mathrm{\Delta }+8p\mathrm{\Delta }\overline{w}24\mathrm{\Delta }\lambda ,`$ (31) $`{\displaystyle \frac{d\overline{w}}{dl}}`$ $`=`$ $`\stackrel{~}{ϵ}\overline{w}+4p\overline{w}^28(p+2)u\overline{w}+48\mathrm{\Delta }\overline{w}24\lambda \overline{w}.`$ (32) Here we have defined $`\stackrel{~}{ϵ}=ϵ2ϵ_\tau `$. Of course, also the distance $`t`$ from the critical point will be renormalized. However, we only consider the flow on the critical surface $`t=0`$ since we are interested in the stability of the critical fixed points. Note that the coefficient of the rare region term $`\overline{w}`$ only couples to $`\mathrm{\Delta }`$. The flow of $`u`$ and $`\lambda `$ is only indirectly influenced by the rare regions (via a modification of the flow of $`\mathrm{\Delta }`$). This will be important later on. ### B Fixed points and their stability The flow equations (III A) possess sixteen fixed points. Their fixed point values are given in Table I, the eigenvalues of the corresponding linearized renormalization group transformations are listed in Table II. For eight of the sixteen fixed points (Nos. 1–8 in Table I) the fixed point value of the rare region term is $`\overline{w}^{}=0`$. These fixed points have already been studied in Ref. using the conventional approach. In the following, we concentrate on the case $`ϵ>0`$ and $`\stackrel{~}{ϵ}=ϵ2ϵ_\tau <0`$ relevant for the itinerant quantum antiferromagnet. We first consider the dirty Heisenberg fixed point (No. 6) and the dirty cubic fixed point (No. 8). These are the stable fixed points of the conventional theory for the cases of $`p<4`$ and $`p>4`$, respectively. Analyzing the stability matrix for the dirty Heisenberg fixed point shows that it is unstable since the eigenvalue $`e_4`$ is positive for $`p<4`$. In contrast, the dirty cubic fixed point remains stable for $`p>4`$ since all eigenvalues of the stability matrix are negative. Thus we conclude that the rare regions destroy the conventional dirty Heisenberg critical behavior for $`p<4`$ while they do not influence the conventional dirty cubic critical behavior for $`p>4`$. We now turn to the new fixed points with $`\overline{w}^{}0`$ (Nos. 9 – 16 in Table I). Fixed points 9, 11, 13 and 15 are unphysical because their fixed point values $`\overline{w}^{}`$ are negative. Since the bare $`\overline{w}`$ is positive and according to eq. (32) the flow cannot cross the $`(\overline{w}=0)`$-plane these fixed points can never be reached. Depending on the number $`p`$ of order parameter components the remaining fixed points (Nos. 10, 12, 14, and 16) are either also unphysical, or they are unstable. Consequently, for $`p<4`$ and to one-loop order there is no stable fixed point. Renormalization group trajectories which in the conventional theory would go to the dirty Heisenberg fixed point show runaway flow to large disorder strength. This runaway flow could either indicate a unconventional phase transition, e.g. an infinite disorder critical point as in the one-dimensional random Ising model or a percolative rather than a homogeneous transition or even a destruction of the phase transition. Within the present approach we cannot be decide between these alternatives. The influence of the rare regions on the stability of the fixed points in our model is similar to that in the isotropic case . For $`p<4`$ the conventional fixed point is destroyed in both models. For $`p>4`$ the conventional fixed point is stable. In our model this is the dirty cubic fixed point while in the isotropic case this stable fixed point is the dirty Heisenberg fixed point. ### C The fluctuation-driven first-order transition In addition to the continuous phase transitions associated with the critical points discussed above there is also the possibility for a first-order transition in the model considered here. Let us first discuss the mechanism for a clean system and discuss the effects of disorder and rare regions later. According to a mean-field stability analysis of the effective action (LABEL:eq:2.10) with $`\mathrm{\Delta }=\overline{w}=0`$ the inequalities $`u+\lambda >0`$ (for $`u>0`$) and $`u+\lambda /p>0`$ (for $`u<0`$) have to be fulfilled for the theory to be stable. Now consider a bare theory with $`u<0,\lambda >0`$ or $`u>0,\lambda <0`$ but still fulfilling the above stability conditions. In these cases the flow equations (III A) can lead the renormalization group trajectories to the mean-field unstable region. This indicates a fluctuation-driven first-order transition . It was later shown that the fluctuation-driven first-order in this model survives the presence of quenched disorder, at least within the conventional theory. Let us now consider the influence of the rare regions. As already mentioned, the rare region coefficient $`\overline{w}`$ does not couple into the flow equations for $`u`$ and $`\lambda `$ but only into the flow equation for $`\mathrm{\Delta }`$. Thus a renormalization group trajectory going to the mean-field unstable region within the conventional theory will generically also do so in the presence of rare regions, the only modification being a different disorder value at the stability boundary. Therefore, we conclude that the fluctuation-driven first-order transition also occurs when taking the rare regions into account. However, since the rare regions modify the flow of the disorder strength $`\mathrm{\Delta }`$, the boundaries of the first-order region may change. ## IV Summary and conclusions We have investigated the influence of rare regions on the quantum phase transition of a disordered itinerant antiferromagnet with cubic anisotropy. The local magnetic moments forming on the rare regions in the Griffiths phase generate a new term in the order parameter field theory which has the form of static annealed disorder. We have found that for order parameter dimension $`p>4`$ this new term does not change the critical behavior, which is characterized by the dirty cubic fixed point. In contrast, for $`p<4`$ the rare region term renders the conventional critical fixed point unstable. The renormalization group trajectories show runaway flow to large disorder. Within our approach which is essentially perturbative, even though it includes some non-perturbative degrees of freedom (the local moments) we cannot determine the ultimate fate of the transition. It could be an unconventional phase transition, e.g. an infinite disorder critical point or a percolative rather than a homogeneous transition or even the destruction of the phase transition. We have also found that the fluctuation-driven first-order transition occurring in this model remains qualitatively unchanged by the rare regions, while the precise position of the first-order region in parameter space will change. The authors acknowledge helpful discussions with D. Belitz, J. Cardy, and T.R. Kirkpatrick. R.N. thanks the hospitality of TU Chemnitz during two vistits where part of the research was performed. This work was supported in part by the DFG under grant nos. SFB393/C2 and Vo659/2, by the NSF under grant no. DMR-98-70597, and by EPSRC under grant no. GR/M 04426.
warning/0004/astro-ph0004009.html	ar5iv	text	# Age and metallicity for six LMC clusters and their surrounding field population ## 1 Introduction In spite of its enormous importance for understanding galaxy evolution in adequate detail, the chemical enrichment process in galaxies is still poorly known, which is especially true for the field star component. The Large Magellanic Cloud (LMC) is a natural target to study the chemical evolution because of its proximity. Also, its structure seems to be less complex than that of the Milky Way which might imply that the chemical enrichment history can be described by a simple global age-metallicity relationship (AMR). First efforts to determine the AMR of LMC clusters have been made with integrated broad band photometry of clusters (Westerlund westerlund97 (1997)). Recent work continuing these studies is, for example, Bica et al. (bica98 (1998)) and Girardi et al. (girardi95 (1995)). Another major step towards an understanding of the LMC cluster AMR has been undertaken by Olszewski et al. (olszewski91 (1991)), who used medium resolution spectroscopy of individual giants to measure the metallicity for around 70 clusters, with a quoted uncertainty of $`\pm 0.2`$ dex. In addition many photometric studies of stars in LMC clusters (e.g. with the Washington system by Bica et al. bica98 (1998)) contributed to the unveiling of the cluster AMR. The current wisdom on the cluster AMR that has been established by these studies is, that the mean metallicity of younger clusters is distinctly higher than that of old clusters by more than 1.2 dex. However, it is difficult to trace the AMR over the entire LMC history with this cluster sample, since for a long time, only one cluster (ESO121-SC03) with an age between 3 Gyr and 11 Gyr had been found (Mateo et al. mateo86 (1986), Bica et al. bica98 (1998)). Recently, Sarajedini (sarajedini98 (1998)) found three more clusters with an age of about 4 Gyr (NGC 2121, NGC 2155 and SL 663). The AMR as derived from LMC clusters shows a very large scatter (Olszewski et al. olszewski91 (1991)), which, if intrinsic and not due to measurement uncertainties, would argue for a more complex chemical enrichment history. In addition there are hints that at least some clusters have smaller mean metallicities than the surrounding field population (e.g. Bica et al. bica98 (1998), Richtler et al. richtler89 (1989)). Thus possibly the chemical evolution of the cluster and field stars is to some degree decoupled. However, this is not without contradiction (e.g. Korn et al. korn00 (2000)). Santos Jr. et al. (santos99 (1999)) claimed that the metallicity dispersion of the field seems to be smaller than that of the cluster system of similar age. For the field population the metallicity distribution is known primarily for the young stars since mainly F & G supergiants have been spectroscopically investigated (e.g. Hill et al. hillV95 (1995), Luck & Lambert luck92 (1992), Russell & Bessell russell89 (1989)). A compilation of young LMC field stars abundances which have been derived with high resolution spectroscopy can be found in the appendix (Table 8). Thévenin & Jasniewicz (thevenin92 (1992)) study 9 field stars in the LMC with medium resolution spectroscopy (5 $`\AA `$) and found an average abundance of $`[Fe/H]=0.25\pm 0.08`$ which is higher than the mean value of field stars that has been derived with high resolution spectroscopy ($`0.38\pm 0.11`$ dex). Dopita et al. (dopita97 (1997)) measured element abundances of planetary nebulae (PNs) in the LMC and derived their age by modelling the hot, central star. They found four PNs that are older than $`4`$ Gyr. Their AMR shows only little enrichment from $`15`$ to $`5`$ Gyr ago, while the metallicity doubled in the last $`23`$ Gyr. The study of the older stellar field component has been limited to studies using broad band photometry (e.g. Holtzman et al. holtzman99 (1999) and Elson et al. elson97 (1997)). We used a different approach and measured the metallicity of individual stars by using the medium wide Strömgren filter system, that gives a good metallicity discrimination for giants and supergiants red-wards of $`by=0.4`$ mag. This method has already been used by Grebel & Richtler (grebel92 (1992)), Hilker et al. (1995b ) and Hilker et al. (1995a ) to determine age and metallicity of NGC 330, NGC 1866 and NGC 2136/37. Ardeberg et al. (ardeberg97 (1997)) used HST observations transformed into the Strömgren system to derive the SFH and the metallicity of LMC bar stars. Their investigation differs from our approach by the calibration they employed which is based on bluer stars and includes the gravity dependent $`c1`$ Strömgren colour index. In the current work we investigate mainly young LMC clusters and their surrounding fields, namely NGC 1651, NGC 1711, NGC 1806, NGC 2031, and NGC 2257, an old cluster. We have also re-analysed NGC 2136/37 because of the availability of Strömgren isochrones and a new calibration for photometric metallicities, which improves the calibration for more metal poor stars. This ensures the homogeneity of the sample and also tests if systematic shifts are present between the older investigations and the new one. An important aspect of this new work is exactly this homogeneity of the metallicities allowing one to assess the real magnitude of the intrinsic dispersion among metallicities of clusters of similar age. Two of the clusters (NGC 2136 and NGC 2031) are particularly interesting because they contain Cepheid variables, whose metallicities are important to know for distance scale problems. NGC 1866 might serve as example. Its metallicity has been determined by Hilker et al. (1995b ) via Strömgren photometry which was used for the distance determination using its Cepheid members by Gieren et al. (gieren94 (1994)). ## 2 Data & Reduction The data have been obtained during two observing runs with the 1.54-m Danish telescope at La Silla, Chile. NGC 2136 and NGC 2031 were observed during 13.11. - 15.11.1992, and NGC 1651, NGC 1711, NGC 1806, NGC 2257 during 4.1. - 7.1.1994. The observing log is shown in Table 5 in the Appendix. We used the UV coated Thomson THX 31560 chip, that has a field of view of $`6.5^{}\times 6.5^{}`$ and a scale of 0.377″/pixel. The Danish imaging Strömgren filters $`v`$, $`b`$ and $`y`$ were used. The measurements in these filters have been transformed into a Johnson $`V`$ magnitude, the colour $`by`$ and the colour index $`m1=(vb)(by)`$. Except for NGC 1711 each cluster has been observed on at least two different nights to have photometrically independent measurements. The reduction included bias subtraction, flat-field correction and the elimination of CCD defects. We used DaoPhot II in the MIDAS and IRAF environments for the photometry. For the calibration we employed six different E-region standard stars from the list of Jønch-Sørensen (jonchsorensen93 (1993)) and four fields with secondary Strömgren standards measured by Richtler (richtler90 (1990)), namely M 67, SK -66 80, NGC 2257 and NGC 330. The heavily crowded field of NGC 330 with the secondary standards measured with photoelectric photometry is problematic since slight differences in centring the aperture on a standard star lead to deviations in the obtained magnitude on the order of $`0.02`$ mag. Thus several stars in this field have been removed from the calibration. The photometric error of a single measurement is given by the standard deviation of the standard stars: $`\sigma _y=0.031`$, $`\sigma _{by}=0.028`$ and $`\sigma _{m1}=0.036`$. We obtained the errors by averaging the residuals of the calibration stars in all nights. The calibration errors for the 1992 run are smaller: $`\sigma _y=0.016`$, $`\sigma _{by}=0.023`$ $`\sigma _{m1}=0.025`$. However, since standard stars are always measured in the central region of the CCD chip, this error does not include flat-field errors which are much harder to quantify (especially also due to the problematic field concentration in telescopes with focal reducers (Andersen et al. andersen95 (1995)). From inspection of the sky background the accuracy of the flat-field is $`12\%`$ and we therefore assigned an additional error of $`0.015`$ to each magnitude. The photometric standards were measured with apertures and thus it was necessary to determine the aperture - PSF shift carefully. The remaining uncertainty is of the order of $`0.03`$ mag. Even if several nights have been averaged this calibration error has been kept, since the calibrations in each night are not truly independent: the colour terms have been determined using the standard star observations from all nights together. Thus we overestimated the calibration error for the clusters by a factor $`<\sqrt{nights}`$ if observations from several nights have been averaged. The calibration error causes the deviation of the measured metallicity from the “true” metallicity of a star to be a function of its colour (shown in Fig. 1). The corresponding metallicity error is larger for blue stars since lines of constant metallicity approach each other on the blue side (see Fig. 8 for illustration). In the following stars bluer than $`by=0.6`$ are excluded from the metallicity determination because of this strong rise of the metallicity uncertainty. ## 3 Metallicity determination via Strömgren colours The major advantage of the Strömgren system compared to broad band photometric systems is the ability to get the metallicity of a star nearly independent of its age The reason for this independence is the minor luminosity effect in the metallicity determination, which amounts to less than $`\pm 0.1`$ dex over a luminosity interval of $`4<M_V<3`$. We have used a new metallicity calibration of the Strömgren $`m1(by)`$ two-colour relation by Hilker (hilker99 (1999)) which is valid in the colour range $`0.5<by<1.1`$. For redder stars the calibration breaks down due to the onset of absorption by TiO and MgH molecules in the $`y`$ band. The used calibration equation is $$\left[\frac{Fe}{H}\right]=\frac{m1_0+a1(by)_0+a2}{a3(by)_0+a4}$$ with $`a1=1.277\pm 0.050,a2=0.331\pm 0.035`$ $`a3=0.324\pm 0.035,a4=0.032\pm 0.025`$ This calibration has been derived using primarily giant stars, however as investigated by Grebel & Richtler (grebel92 (1992)) it should also apply to supergiants. For stars bluer $`by=0.7`$ this has been predicted by Gustafsson & Bell (gustafsson79 (1979)). We will discuss this question in greater detail in Sect. 3.4. The reason for the metal sensitivity is the line blocking in the $`v`$ filter, which is best measurable for G and K stars. The measured flux depends largely on the strength of the Fe I lines, but also CN and CH bands contribute. Systematic deviations of less than $`0.1`$ dex are expected from theoretical isochrones due to a small luminosity dependence. ### 3.1 The CN anomaly A severe problem in the interpretation of Strömgren colours is the contribution of the CN molecule absorption (band head at $`421.5`$ nm) in the $`v`$-filter ($`410`$ nm, width $`20`$ nm) to the line blocking. CN variations have been observed in several galactic globular clusters, however, the exact mechanism is not yet fully understood. An increased CN abundance leads to an increased photometric metallicity and thus to a decreased age if it is derived via isochrones. As a rule of thumb we estimated with the aid of Geneva isochrones that an increased metallicity of $`0.2`$ dex will decrease the age by $`20\%`$. A recent investigation using Strömgren photometry of two globular clusters, of which one has CN anomalous stars, the other not, illustrates the effect of CN anomaly on the Strömgren metallicity (Richter et al. richter99 (1999)). We cannot account for this CN anomaly. In this study we have to live with this uncertainty, but there is evidence that this effect is only modest: in the nearby giant sample (see next section) three stars were assigned to be CN enriched, but they do not deviate within the standard deviation from the CN-normal stars, however most probably due to the uncertain reddening correction. Pilachowski et al. (pilachowski96 (1996)) found that for population II halo stars the CN anomaly does not play an important role, in contrast to M 13 for example. This cannot be explained by simple selection effects. Mc Gregor & Hyland (mcgregor84 (1984)) found a general CN deficiency by weaker CO bands in the LMC than in galactic supergiants of the same temperature. Concerning our LMC field stars we found a good agreement for the young stars with spectroscopic analyses as well as for the older field population with observed clusters (see below). Therefore anomalous CN abundances should not play a devastating role. Ultimately this can only be checked with a spectroscopic investigation of the CN behaviour of cool LMC giants and supergiants. ### 3.2 The influence of reddening uncertainties Photometrically measured metallicities are very sensitive to reddening errors, which is a major error source of the determined metallicities. For example an underestimation of the reddening of 0.02 mag in $`E_{by}`$ leads to an average underestimation of the metallicity by $`\mathrm{\Delta }[Fe/H]=0.1`$ dex for a fully populated RGB with an age around $`10^{9.0}`$ yr. This problem is symptomatic for photometric investigations. For example Bica et al. (bica98 (1998)) who used Washington photometry to derive ages and metallicities of old LMC clusters and the field, stated that “an increase of the assumed reddening by E(B-V)=0.03 decreases the derived metallicity by 0.12 dex”. The degeneracy between reddening and abundance becomes a severe problem for old clusters and field stars, while for young clusters it is possible to determine the reddening quite accurately because the colour of the hot, bright main sequence stars is nearly independent of temperature and thus of metallicity. The dependence of the derived metallicity on the assumed reddening is illustrated in Fig. 2, with NGC 1806 as an example (we note that this figure greatly exaggerates the realistic uncertainty of the reddening for this particular cluster and just serves to illustrate the trend). Differential reddening is another aspect of this problem. We cannot exclude it, however there are also no hints in favour of strong differential reddening. Olsen (olsen99 (1999)) investigated four fields in the LMC (three in the bar, one in the inner disk), where the reddening is expected to be larger than in our further outside lying fields. However, they detected strong differential reddening only around NGC 1916. For the other fields it is not significant. We estimated with the aid of Monte Carlo simulations, that as long as the differential reddening is less than $`E_{BV}=0.03`$ (peak to peak), the uncertainty is small compared to the photometric uncertainty. In any case differential reddening results in a broadening of the metallicity and hence the age distribution. ### 3.3 Unresolved binaries and blending The observation of unresolved binaries, which is likely to be the case for a considerable fraction of stars, leads to a change in the photometric metallicity. Fortunately this effect plays a negligible role for stars on the RGB where the mass-luminosity relation is steep and even small differences in the initial mass result in large differences in luminosity. This has been checked with the aid of synthetic CMDs and two-colour diagrams, which is described in Sect. 12 below. In crowded fields blending of stars is a related problem which is nicely illustrated in Fig. 1 in the work of Ardeberg et al. (ardeberg97 (1997)). The first correction is to exclude the most crowded inner part of the clusters thus we excluded stars within $`19\mathrm{}`$ from the cluster center (details in the cluster sections). The most probable, but unimportant, case is blending with a faint, red main sequence star: they are not luminous enough to change the photometric metallicity of a RGB star. To estimate the probability of blending with other stars we used the following approach: we assumed that blending takes place if the luminosity centres of gravity of two stars are nearer than the PSF radius divided by $`\sqrt{2}`$ (sampling theorem). Since the PSF radius was always around $`3`$ pixels the area in which only one unblended star can be is $`(23/\sqrt{2})^2\pi `$. Next we counted the stars on the blue and red side of the observed CMD (bluer and redder than $`(by)=0.4`$) in luminosity intervals fainter than the main star and calculated the probability that one of these fainter stars lie within the area of the main component. We label the blending with a star “strong blending”, if the luminosity difference of this star and the main component is less than $`2`$ mag. Blending with a star that is between $`2`$ and $`4`$ magnitudes fainter is called ”weak blending” and corresponds to a luminosity ratio of at least $`6`$ that results in shifts of $`<0.2`$ dex. The probability of weak blending is underestimated since incompleteness has not been considered. However, weak blending primarily results in a broadening of the metallicity distribution of $`<0.15`$ dex. Also the strong blending is slightly underestimated since a truly blended star is counted in the observed CMD only once, but since the probability for strong blending is well below $`10\%`$ (see below) we regard this approximation to be justified for our fields. Blending with a main sequence star results in a shift of the combined pair towards bluer colour and and smaller m1. The shift in $`(by)`$ dominates and thus the resulting metallicity is in general larger for this pair than for the individual RGB star. Very strongly blended stars even leave the selected colour range and thus our selection also ensures the exclusion of heavily blended stars. The fields of NGC 1711 and NGC 1806 shall serve as examples for the expected blending probability. These fields are rather crowded compared to the fields around NGC 1651, NGC 2257 and NGC 2136, but comparable with the crowding around NGC 2031. The strong blending probability with a blue main sequence star decreases from $`8\%`$ for a $`18.5`$ mag RGB star to less than $`1\%`$ for a $`15`$ mag star in the field of NGC 1711. The probability for strong blending with a red star decreases from $`3\%`$ to $`<1\%`$. Around NGC 1806 the strong blending with red stars is more probable: it decreases from $`10\%`$ to $`<1\%`$ for a luminosity of the main component between $`18.5`$ mag and $`15`$ mag. The probability of weak blending is for a $`17`$ mag RGB stars $`10\%`$ around NGC 1806 and $`5\%`$ around NGC 1711. We conclude that $`5\%`$ of our selected stars have a Strömgren metallicity that deviates by more than $`0.2`$ dex from its “true” value due to blending. Around $`10\%`$ of the stars are blended with a resulting shift of $`<0.2`$ dex. The deviations in age due to blending are more difficult to determine: on the one hand the increasing luminosity would result in an underestimation of the age, on the other hand for strong blending with a blue star, the metallicity would be underestimated, which generally leads to an overestimation of the age. We conclude that also for the age, blending results primarily in a broadening of the age distribution, however, with a distribution that is more extended towards younger ages, i.e. it is more probable to underestimate than to overestimate the age (this has been found with the aid of the mentioned simulation). In the discussion in Sect. 14 we will give an additional argument that the blending in clusters result in shifts of less than $`0.1`$ dex compared to the field, assumed that cluster and field population of the same age have the same metallicity based on the observations. ### 3.4 AGB versus RGB stars AGB stars are in the age range of $`10^{8.4}`$ yr to $`10^{9.0}`$ yr the dominating giant stars. These stars are potentially problematic since their surface abundances might have changed considerably compared to the initial composition. However, in M 13 Pilachowski et al. (pilachowski96 (1996)) found that the AGB stars are less CN enriched than the RGB stars. Frogel & Blanco (frogel90 (1990)) identified several AGB stars around some of our clusters, for which Olszewski et al. (olszewski91 (1991)) obtained the metallicity. However, these stars are always too red to allow a photometric metallicity determination. ### 3.5 Are the Strömgren metallicities independent of the luminosity class? The re-calibration of the Strömgren metallicity of red stars by Hilker (hilker99 (1999)) is based on giant stars for the lower metallicity range and approach the calibration of Grebel & Richtler (grebel92 (1992)) for higher metallicities (however also below 0 dex). In the earlier work by Grebel & Richtler (grebel92 (1992)) no difference between giant and supergiant stars has been found. ¿From the theoretical point of view, only a very small luminosity effect is expected (Bell & Gustafsson bell78 (1978) and the used isochrones by Grebel & Roberts 1995a ). To reinvestigate observationally the dependence of the Strömgren metallicity on the luminosity class, we selected supergiants (luminosity class I & II) with metallicity measurements and Strömgren photometry from the compilation of Cayrel de Strobel et al. (strobel97 (1997)) and SIMBAD, respectively. The major problem of this approach is the largely unknown reddening towards these galactic field supergiants. To exclude stars that are most probable highly reddened we took only supergiants with a galactic latitude $`\|l\|>20^0`$ and brighter than $`6`$ mag in $`V`$ into account. The remaining supergiants are shown in Fig. 3, where the difference of the Strömgren metallicity to the measured metallicity (taken from the list of Cayrel de Strobel et al. strobel97 (1997)) versus the measured metallicity is displayed. Fig. 3 open circles are used for variable stars, and open star symbols for carbon stars and stars with a CN anomaly. We did not attempt to correct for the individual reddening, which explains partially the considerable scatter. The vertical lines shows the location of the literature metallicity of $`0.1`$ dex and the horizontal line indicates where spectroscopic and photometric metallicities are equal. It can be seen in Fig. 3 that the calibration seems to hold for supergiants with a metallicity of less than $`0.1`$ dex (we have corrected for a systematic shift of $`0.1`$ dex which can easily be explained with an average reddening of $`E_{BV}=0.02`$). For larger metallicities the Strömgren metallicities seem to underestimate the “true” metallicity. However, the calibration has been made only for stars of subsolar metallicity, thus the deviation of more metal rich stars is not surprising. ## 4 The age determination To determine ages of stars and clusters we employed isochrones provided by the Geneva (Schaerer et al. geneva93 (1993)) and Padua (Bertelli et al. padua94 (1994)) groups transformed into the Strömgren system by Grebel & Roberts (1995a ); they are called ”Geneva” and ”Padua” isochrones in the following. The isochrones show a zero point difference to the empirical calibration in the sense that an isochrone of a given metallicity is too red and/or $`m1`$ is too low compared to the calibration, which is of the order of $`0.15`$ dex. To solve this discrepancy we increased the $`m1`$ values of the isochrones by $`0.04`$ mag to bring them into accordance with the empirical calibration. We have chosen the $`m1`$ colour index, since it contains the $`v`$ filter, which is the most critical one in the filter-band integration of the model spectra due to it’s fairly short wavelength. For this central wavelength the applied stellar atmospheres of red giants and supergiants are not very precise (Bressan priv. comm.). Unfortunately, only Geneva isochrones with metallicities of more than $`1.4`$ dex and ages of less than $`10^{9.9}`$ yr were available. Padua isochrones on the other hand covered only the age range between $`10^{7.0}10^{9.0}`$ yr and $`10^{10.0}10^{10.24}`$ yr, thus most of our results are based on the Geneva isochrones. The very red part of the RGB is not red enough to describe the location of the observed RGB stars correctly, the isochrones are slightly too steep for $`(BV)>1.1`$. With the observed RGBs of NGC 1651, NGC 1806 and the field surrounding NGC 1711, we introduced an empirical linear colour term to bring the isochrones into agreement with these stars (for $`(by)>0.7`$ the applied shift is: $`0.6((by)0.7)`$, the maximum correction is $`\mathrm{\Delta }(by)=0.12`$). After we applied this correction the isochrones fit also to the younger clusters, which is a hint that the colour term is valid for all gravities. However, nothing can be said for old stars ($`>10^{9.5}`$ yr) since no clear RGB with such an age was available to test the empirical colour term (NGC 2257 is too metal poor and too old for this purpose). The $`m1`$ had to be changed according to the colour term in $`by`$, which has been performed on the basis of the empirical metallicity calibration. The isochrones show a small age-metallicity degeneracy: a substantial age difference between a $`10^{7.5}`$ yr and a $`10^{9.9}`$ yr isochrone leads only to a difference in the photometric metallicity of $`\mathrm{\Delta }[Fe/H]=0.2`$, in the sense that younger stars would appear more metal rich. Since we adjusted a $`10^{9.0}`$ yr isochrone (via the $`m1`$ shift) to the empirical metallicity calibration, the resulting metallicity uncertainty is less than $`0.1`$ dex. Throughout the paper we assumed a distance modulus of 18.5 for the LMC based on surface brightness analysis of Cepheids (Gieren et al. gieren98 (1998)), results from SN1987A (Panagia et al. panagia91 (1991)) and on the recent revision of the ”classical” Cepheid distance calibration (Madore & Freedman madore98 (1998)). A distance uncertainty has a direct effect on the age determination in the sense that a smaller distance to the LMC would result in lower ages. ## 5 Selecting cluster and field stars ### 5.1 Cluster stars The first step in measuring the metallicity of a cluster is to separate its members from the surrounding field population. We performed this mainly by selecting stars within a certain radial distance from the cluster. The selection radius is defined as the radius where the cluster star density starts to be higher than $`2\sigma `$ over the background star density, which has been derived with a radial density profile of the stars in the frame. The innermost part ($`<20\mathrm{}`$) of the clusters has been excluded because it is impossible to derive reliable photometry for stars in this crowded region, especially due to blending. For young clusters also the luminosity of a star is a good criterion to separate cluster and field stars, since it is easy to distinguish bright, young cluster stars from old field RGB stars. Clearly this criterion does not separate field and cluster stars of similar age. The lower luminosity limit that was used to exclude RGB field stars in this approach has been determined by visual inspection of the CMD. It is straightforward for NGC 1711, NGC 2031 and NGC 2136/37, where the cluster stars are much brighter than the field RGB, however, this criterion could not be applied for NGC 1651, NGC 1806 and NGC 2257. We did not perform a statistical field star subtraction for two reasons: we wanted to have the most reliable cluster stars and including stars from a larger radius might lead to a bias towards field stars that have a similar age and metallicity, since the statistical field star subtraction has to work in chunks of colour and luminosities. If the colour and luminosity range of the bins in which the field stars are subtracted are chosen too small than the Poisson error is large, if they are chosen too large, then the resulting distribution is not “cleaner” than the one in our approach. Moreover, since the numbers of stars are frequently much smaller than for the field population, the error in the number of stars would heavily depend on the field star population, especially if the incompleteness varies strongly with radial distance from the cluster. This varying incompleteness would lead to a considerable uncertainty. The disadvantage of our approach is certainly that there will always be some field stars left. Stars with a photometric error (DaoPhot) of more than $`\mathrm{\Delta }(by)=0.1`$ and $`\mathrm{\Delta }m1=0.1`$ have been discarded. This selection ensures more reliable results. Finally, only stars redder than $`(by)_0=0.6`$ have entered the metallicity measurement to reduce the systematic shifts in the derived metallicity due to a possible error in the applied reddening correction, which is illustrated in Fig. 2. Also stars being redder than $`(by)=1.1`$ have been discarded due to additional lines in the $`y`$ filter. From the remaining sample, individual stars have been excluded if they deviate strikingly from the mean metallicity or from the mean RGB location of the other stars, because still a few field and foreground stars might be present, as mentioned above. ### 5.2 Field stars Field stars have been selected with a radial selection criterion as well: we regard as field stars those with a radial distance of more than $`30\mathrm{}`$ plus the radius used for selecting the corresponding cluster The additional $`30\mathrm{}`$ have been added to ensure that cluster stars are a minor fraction among the field stars. Only field stars with a relative photometric error of less than $`\mathrm{\Delta }(by)=0.1`$ and $`\mathrm{\Delta }m1=0.1`$ have been kept. In case of the young clusters NGC 1711, NGC 2031 and NGC 2136/37 we used all RGB stars having a distance of more than $`70\mathrm{}`$ from the cluster center, since they clearly do not belong to the cluster. To minimise the influence of photometric and calibration errors on the derived metallicity, it is necessary to introduce the colour criterion $`(by)_0>0.6`$, as it has been done in the case of clusters. However, it is dangerous to limit the sample just in colour since this introduces a large bias towards metal poor stars with larger metallicity errors <sup>1</sup><sup>1</sup>1$`\mathrm{\Delta }m1`$ is the dominating error in the Strömgren two-colour diagram. Therefore, the metallicity error will be larger for more metal poor stars than for more metal rich stars of the same colour.. This is not a big problem for younger clusters, where the giants extend far into the red and therefore the metallicity measurement does not depend so severely on the blue stars. To circumvent this problem and to have nevertheless a reasonable homogeneous selection criterion, we included only stars that are redder than an inclined line in the $`m1(by)`$ diagram that is nearly perpendicular to $`[Fe/H]=1`$ dex, a metallicity which is in the middle of the expected metallicity range in the LMC. This line is shown for example in the two-colour plot of the field population around NGC 1711 (Fig. 8). ### 5.3 Galactic foreground stars Foreground stars of our own Galaxy contaminate the field and cluster sample in the observed fields. Most of the foreground stars are red clump stars which show up as a broad vertical strip at $`by0.4`$. Ratnatunga & Bahacall (ratnatunga85 (1985)) presented a galaxy model and give the amount of galactic foreground stars in luminosity and colour bins. Their results are compiled in Table 2. With these numbers in mind it is obvious that we expect only few galactic foreground stars in the colour and luminosity range we used for the metallicity and age determination. ## 6 The reddening towards individual clusters and the surrounding field Because of the large influence of the reddening on the measured metallicity, as described in Sect. 3, it is necessary to get a hand on the reddening correction towards the observed regions. For this purpose we used the theoretical upper main sequence ($`(by)_0<0.1`$ and $`M_V<0`$) for the reddening determination, because of the negligible metallicity metallicity effects. It is essential not to fit the isochrone to the centre of the main sequence, since the isochrones are calculated for non-rotating stars and evolutionary effects on the upper main sequence. Rotation shifts a star redwards and thus one has to fit the isochrone more to the blue border of the observed main sequence. Also unresolved binaries on the main sequence are redder than the observed isochrones. However, since photometric errors are also present the fit of a blue envelope would be exaggerated. For the extinction correction the relations of Crawford & Barnes (crawford70 (1970)) have been employed ($`E_{by}=0.7E_{BV}`$ and $`E_{m1}=0.3E_{by}`$). Since the reddening is derived on the assumption that the colour of the isochrone is correct for the very bright blue main sequence we only give the possible uncertainty in the adjustment of the isochrones to the main sequence as a reddening error. However, our reddening is always smaller than the reddening given by Schlegel et al. (schlegel98 (1998)), which is a hint that there is a zero point shift in $`by`$ between either our calibration or the isochrones on the order of $`\mathrm{\Delta }by=0.03`$. It might of course be a zero point shift in the Schlegel et al. values as well, especially when considering that the given reddening values are frequently larger than the one by other authors (see Sect. 6 - Sect. 11). ## 7 NGC 1711 ### 7.1 The cluster Several attempts have been made to determine the age of NGC 1711. One of them used also isochrone fitting (Sagar & Richtler sagar91 (1991)), however, with an assumed metallicity of $`0.4`$ dex. The previous results on the age of NGC 1711 are compiled in Table 6 in the appendix. No metallicity measurement for this cluster has been published yet. The CMD of the entire CCD field is shown in Fig. 34. Using the upper main sequence ($`V<18.5`$) we have deduced a reddening of $`E_{by}=0.06\pm 0.02`$, relative to the isochrone, which corresponds to $`E_{BV}=0.09\pm 0.05`$ (the calibration error has already been included). The determined reddening agrees, within the errors, with $`E_{BV}=0.14`$ given by Cassatella et al. (cassatella96 (1996)) and with the reddening derived from measurements by Schwering & Israel (schwering91 (1991)) ($`E_{BV}=0.11`$). Burstein & Heiles (burstein82 (1982)) gives $`E_{BV}=0.12`$. Concerning the reddening, it is important to note that the surrounding field, where one can observe young stars with a similar age, shows a $`0.02`$ mag higher reddening, indicating that NGC 1711 is located in front of the LMC disk. Fig. 4 illustrates this difference. This also holds for a small concentration of brighter stars south-east of NGC 1711. Unfortunately, the number of stars in this group is not large enough to allow a reliable age or metallicity determination. However, because of the same colour difference between field and cluster, it might form a binary cluster with NGC 1711, a configuration, which seems to be common in the LMC <sup>2</sup><sup>2</sup>210% of the LMC clusters are thought to be paired (e.g. Dieball & Grebel dieball98 (1998) and references therein). By inspecting the radial number density of stars around NGC 1711, we have found that cluster stars begin to dominate ($`2\sigma `$ level) the field stars at a radial distance of $`90\mathrm{}`$. This radius has been used for the radial selection. To exclude older field RGB stars, we regarded only stars brighter than $`V_0=15.5`$ as potential cluster members for the metallicity analysis. The CMD and two-colour diagram of NGC 1711 is presented in Fig. 8 and Fig. 8, respectively. In the two-colour diagram we plotted the calibration error separately and assigned only the photometric error to the individual stars. The effect of these two errors is completely different: the metallicity error due to the photometric errors decrease with increasing size of the sample, while in contrast the metallicity error due to the calibration can only be decreased if observations from different nights are averaged. We measured a metallicity of $`[Fe/H]=0.57\pm 0.06`$ dex for NGC 1711. The error is the standard deviation of the individual stars divided by the square root of the number of used stars ($`5`$). Reddening and calibration error account for an additional error of $`0.16`$ dex, thus we finally obtained $`[Fe/H]=0.57\pm 0.17`$. To be able to fit the red supergiants one needs an isochrone with a metallicity of at least $`0.4`$ dex, despite the measured metallicity. We showed in Sect. 3.2, that the calibration is valid for supergiants, as long as they have a metallicity below $`0`$ dex. From Fig. 3 one could estimate that a supergiant with a metallicity around $`0.2`$ dex could be mistaken for a $`0.5`$ dex star. However, we prefer a different explanation: as described in Sect. 3, a slight age dependence exist accounting for $`\pm 0.1`$ dex, in the sense that younger stars appear more metal poor. With this in mind we derived a metallicity of $`[Fe/H]=0.45\pm 0.2`$ (we assigned an additional error of $`0.05`$ because of the uncertainty of the shift). A third possibility is that the isochrones does not sufficiently extend towards the red for these bright stars, despite the empirical correction. This is mainly a problem in the treatment of overshooting and a common problem for red supergiants (Bressan priv. comm.). We arrived at an age of $`10^{7.70\pm 0.05}`$ yr using Geneva isochrones. The isochrone is overlayed in Fig. 8. ### 7.2 The surrounding field population We used only bright ($`V<16`$) stars that are more than $`110\mathrm{}`$ away from the cluster centre and all RGB stars with a distance of $`>50\mathrm{}`$. The field population consists of a main sequence slightly more reddened and an older population clearly distinguishable by its giant branch. A remarkable feature is the vertical extension of the red clump (VRC), which has been observed in other areas of the LMC as well. This has been interpreted by Zaritsky & Lin (zaritsky97 (1997)) as a signature of an intervening population towards the LMC, but Beaulieu & Sacket (beaulieu98 (1998)) showed that also normal stellar evolution could lead to such a feature, if a $`10^{8.5}`$ yr to $`10^{9.0}`$ yr old population is present. Even the fainter extension (or fainter second red clump), might be present (Girardi girardi99 (1999), Piatti et al. piatti99 (1999)), however the numbers are definitely too small to allow an unambiguous identification. With a reddening of $`E_{BV}=0.11`$ we can derive the metallicity of the field stars which is illustrated in Fig. 8. The average metallicity of the field population is $`[Fe/H]=0.53`$ dex and the standard deviation $`0.42`$ dex. The unambiguously young field stars, which are marked with star symbols in Fig. 8 and Fig. 8, have a mean metallicity of $`0.56\pm 0.27`$ which is not systematically larger than the one of the older population, even when accounting for the slight age dependence of the metallicity. Stars in the narrow metallicity range $`0.75<[Fe/H]<0.45`$ (filled circles) do not exhibit a uniform age. An upper age limit of the field stars is $`10^{8.9}`$ yr. This means that between $`10^{7.7}`$ yr and approximately $`10^{8.9}`$ yr no clear age-metallicity dependence can be seen in this field. ## 8 NGC 1806 ### 8.1 The cluster NGC 1806 is older than NGC 1711, which can immediately be seen from its pronounced RGB (see Fig. 12). No previous CCD CMD is available in the literature. A faint main sequence of a younger field population is visible in the field around NGC 1806 (Fig. 34). This population has been used to determine a reddening of $`E_{by}=0.12\pm 0.02`$ ($`E_{BV}=0.17\pm 0.03`$), which agrees, within the errors, with the values given by Cassatella et al. (cassatella87 (1987)) ($`E_{BV}=0.12`$). Schwering & Israel (schwering91 (1991)) give $`E_{BV}=0.10`$ and Schlegel et al. (schlegel98 (1998)) $`E_{BV}=0.24`$. Burstein & Heiles (burstein82 (1982)) derived a reddening of $`E_{BV}=0.06`$ towards this direction. We assumed that the cluster is reddened by the same amount as the field main sequence stars. The strong dependence of the derived metallicity on the reddening correction has been demonstrated for this cluster in Fig. 2. Only stars within the radial distance of $`60\mathrm{}`$ have been taken for the metallicity determination. In addition, some stars lying apart from the average RGB location have been excluded. These stars are marked with open star symbols in the cluster CMD (Fig. 12). The inclusion of these stars would not change the derived metallicity. We obtained a metallicity of $`[Fe/H]=0.71\pm 0.06`$ dex for NGC 1806. The calibration and reddening uncertainty result in an additional error of $`\mathrm{\Delta }[Fe/H]=\pm 0.23`$ dex, hence $`[Fe/H]=0.7\pm 0.24`$ dex. The reddest cluster star (that has been excluded because of its red colour) is an identified AGB star (Frogel & Blanco frogel90 (1990), LE 6), for which Olszewski et al. (olszewski91 (1991)) determined a metallicity of $`0.7`$ dex, which is in good agreement with our value for this cluster. The two-colour diagram is presented in Fig. 12. Some remarkable stars are those being redder $`(by)=1.0`$ and apparently more metal poor than the other cluster stars. If these stars were members of a true metal poor population one would have expected to find bluer stars of the same metallicity, which is not the case. Because of this reason we think that these stars belong to NGC 1806, but they possess additional absorption lines in the $`y`$ band compared to the bluer RGB stars. This is theoretically expected for red stars with solar metallicity, however the theoretical Geneva isochrones of this metallicity do not extend far into the red regime and no Padua isochrone with appropriate metallicity and age has been available. ¿From the identified AGB star one might speculate that these deviating stars are AGB stars in NGC 1806, which is supported by the best fitting isochrone (see below). The age determination is illustrated in Fig. 12. The best fitting isochrone yields an age of $`10^{8.7\pm 0.1}`$ yr. This is much younger than the age of $`10^{9.6\pm 0.1}`$ yr derived by Bica et al. (bica96 (1996)) using the SWB classification. The red branch of this isochrone consists mainly of AGB stars especially in the employed colour range for the metallicity determination. ### 8.2 The surrounding field population The metallicity of the field population has been derived with all stars that have a distance of at least $`100\mathrm{}`$ from the cluster centre. With the above stated reddening the mean metallicity of the field population can be obtained as $`0.67`$ dex with a (relatively small) standard deviation of $`0.23`$ dex. A feature that is visible in the two-colour diagram of NGC 1806 (Fig. 12) is that stars that are redder than $`by1`$ seem not to follow a straight line for a given metallicity, but rather get smaller $`m1`$ values with increasing $`by`$. The deviation is of the order $`0.5`$ dex. This behaviour is similar to what is observed among the cluster stars of NGC 1806. Again, we argue that these stars might deviate from the line of constant metallicity given by the calibration. If these stars belong to a true metal-poor population one would expect to find more metal-poor stars with a colour between $`0.9<(by)_0<1.1`$; this is not the case. No star is found in the whole colour range with a metallicity of lower than -1.3 dex, while four stars are found in an even smaller colour range of $`\mathrm{\Delta }(by)=0.1`$. Therefore, one would expect to find at least eight stars in the bluer colour range when assuming a homogeneously populated RGB, which is not unreasonable since the bottom part of a RGB around this age usually is even more populated than the upper part of the RGB. These peculiar red stars could be foreground stars as well, however it is intriguing that they are mixed with the other RGB stars in the CMD. Therefore we think that they belong to the LMC. Thus we reconfirm the statement that even low metallicity stars are only good tracers for metallicity as long as they are bluer than $`by=1.1`$. ## 9 NGC 2136/37 ### 9.1 The cluster NGC 2136/37 is a potential triple cluster system (Hilker et al. 1995a ) and thus another example of the common multiplicity among LMC clusters. The main components have an angular separation of $`1\mathrm{}.34`$. We have re-investigated this cluster due to the availability of Strömgren isochrones and because of the new calibration. NGC 2136 contains at least eight Cepheids making the knowledge of its metallicity particularly interesting for the Cepheid distance scale and the metallicity dependence of the PLC relation. The Cepheids have not been included in the derivation of the metallicity. After the inspection of the radial number distribution of stars around the cluster centre we selected all stars with a distance of less than $`75\mathrm{}`$. Additionally we excluded probable RGB stars of the field population with $`V>16.5`$. The reddening can be determined with the upper main sequence and we obtained $`E_{by}=0.07\pm 0.02`$ ($`E_{BV}=0.10\pm 0.03`$). This agrees with the reddening of $`E_{BV}=0.09`$ given by Schwering & Israel (schwering91 (1991)). Burstein & Heiles (burstein82 (1982)) obtained a reddening of $`E_{BV}=0.075`$. The resulting metallicity of NGC 2136 is $`[Fe/H]=0.55\pm 0.06`$ dex (see Fig. 16). Two stars have been excluded, one metal rich one with solar metallicity and a metal poor one with $`1`$ dex. The more metal rich star is most probably a binary star or the centre of a background galaxy, since its $`\chi ^2`$ value given by the DaoPhot PSF fitting routine is worse than for stars with comparable luminosity. The more metal poor star might be a remaining field star. The stars used for the metallicity determination are shown together with the $`[Fe/H]`$ histogram in Fig. 16. Including the calibration and reddening error we got $`0.55\pm 0.23`$ dex. This is the same value as Hilker et al. (1995a ) one obtained. The age of this cluster is $`10^{8.0\pm 0.1}`$ yr. In Fig. 16 Geneva isochrones with an age / metallicity of $`0.4`$ dex / $`10^{7.9}`$ yr and $`0.7`$ dex / $`10^{8.1}`$ yr are overlayed. For NGC 2137 we have chosen a radial selection radius of $`20\mathrm{}`$. Within this radius two stars remain after applying the usual selection criteria. These stars are plotted with stars symbols in the two-colour plot and CMD of NGC 2136/37 (Fig. 16 and Fig. 16). The metallicities and ages of NGC 2136 and NGC 2137 agree well. Therefore, it is plausible that these clusters are a physical pair and not just a chance superposition, as Hilker at al. (1995a ) already stated. ### 9.2 The surrounding field population The two-colour diagram of the field star population is shown in Fig. 16 and the corresponding CMD in Fig. 16. The stars brighter than $`V=16.2`$ are younger than the majority of the RGB stars and are marked with open stars in the two-colour diagram and the CMD of the field population. Unlike the case of the field population around NGC 1711, the younger stars have a lower metallicity than the dominating older RGB field stars, however, also a large fraction of RGB stars have the same metallicity. We measured $`[Fe/H]=0.75`$ dex for the mean metallicity and $`0.59`$ dex for the standard deviation. For the younger population (star symbols in Fig. 16) we derive an abundance of $`[Fe/H]=0.46`$ and a standard deviation of $`0.11`$ dex. In Fig. 16 an isochrone with a metallicity of $`0.4`$ dex and an age of $`10^{8.0}`$ yr is plotted that fit these stars. ## 10 NGC 2031 ### 10.1 The cluster To select the members of NGC 2031 we chose (from the radial density distribution of stars) $`75\mathrm{}`$ as a good radius to separate cluster and field stars effectively. In addition to our usual selection criteria we have also excluded the very metal poor star with $`2.3\pm 0.3`$ dex, which is most probably a foreground star, judging from its very deviant metallicity. Another excluded star is slightly above the cluster RGB. However, its metallicity fits well to the metallicity of the cluster. The stars used for the metallicity and age determination are shown in Fig. 20 and Fig. 20. For this cluster we have found a reddening of $`E_{by}=0.06\pm 0.03`$ ($`E_{BV}=0.09\pm 0.04`$). Mould et al. (mould93 (1993)) quote $`E_{BV}=0.18\pm 0.05`$ based on HI measurements. Schlegel et al. (schlegel98 (1998)) give the even larger value of $`E_{BV}=0.3`$ as galactic foreground reddening in this direction. On the other hand derived Schwering & Israel (schwering91 (1991)) a reddening of $`E_{BV}=0.1`$. Burstein & Heiles (burstein82 (1982)) derived a reddening of $`E_{B=V}=0.07`$. Using our reddening value a metallicity of $`[Fe/H]=0.52\pm 0.21`$ can be derived (the error includes the reddening error and the calibration error). With this metallicity we found the best fitting isochrone to be $`10^{8.1\pm 0.1}`$ yr. This agrees well with the age determined by Mould et al. (mould93 (1993)) ($`10^{8.14\pm 0.05}`$ yr with $`0.4`$ dex). ### 10.2 The surrounding field population The mean metallicity of the field population for this cluster is $`[Fe/H]=0.75`$ dex and the standard deviation $`0.44`$ dex. Definitely young stars are selected in the CMD (Fig. 20) and are marked with open star symbols, while the older stars which can be used to determine a metallicity are marked with filled circles. The metallicity distribution of the older stars is shown in the small panel with the solid line and the distribution of the younger stars with a dashed line. It can be seen, that we do not find young metal poor stars, although we also find older stars with the same metallicity as the younger ones. However, the younger stars are in average more metal rich than the older stars. ## 11 NGC 1651 ### 11.1 The cluster The RGB of NGC 1651 merges with the RGB of the field population which can be seen in the CMD of all stars found in the field around NGC 1651 (Fig. 34). This makes the visual separation of field and cluster stars in the CMD impossible. It is only possible to distinguish a few young field stars unambiguously from the mixture of field and cluster RGB. We decided, after inspecting the radial number density of stars around NGC 1651, to use a radius of $`55\mathrm{}`$ for the field and cluster separation. In addition we excluded three stars lying well below the cluster RGB. They are most probably foreground stars in the Galactic halo because a LMC star with such a luminosity and metallicity would have an unreasonable age of more than $`20`$ Gyr. Also two stars above the mean RGB of the cluster stars have been excluded. All these excluded stars have been marked in the CMD and two-colour diagram of NGC 1651. The remaining stars used for the metallicity and age determination are shown as filled circles in the two-colour diagram (Fig. 24) and CMD ( Fig. 24). The reddening given by Schwering & Israel (schwering91 (1991)) is $`E_{BV}=0.08`$, Mould et al. (mould97 (1997)) used $`E_{BV}=0.1`$ for this cluster, Schlegel et al. (schlegel98 (1998)) found $`E_{BV}=0.14`$ towards this direction and Burstein & Heiles give $`E_{BV}=0.1`$. With a reddening of $`E_{BV}=0.1`$ the cluster metallicity would be $`0.28\pm 0.02`$ dex. However, no fitting isochrone with such reddening and metallicity can be found. The theoretical RGB stars of this rather large metallicity would be too red or - for younger ages - the main sequence should be visible. The only possibility to fit an isochrone is to use a smaller reddening and hence a lower metallicity. Only with such a metallicity a self-consistent fit in the CMD and the two-colour diagram can be found. Using $`E_{by}=0.03`$ ($`E_{BV}=0.04`$) the cluster metallicity derived is $`[Fe/H]=0.58\pm 0.02`$ dex. The two-colour diagram for this reddening is shown in Fig. 24. The calibration error accounts for an additional error of $`0.19`$ dex. Using this metallicity the age is $`10^{9.3\pm 0.1}`$ yr. The isochrone is overlayed to the cluster CMD in Fig. 24. However, this solution is not unique: also with a smaller reddening acceptable fits are possible, resulting in larger ages and smaller metallicities, hence we only got a range of parameters for this cluster: $`0.01<E_{by}<0.05`$, $`0.65<[Fe/H]<0.45`$, $`10^{9.4}>lg(Age)>10^{9.0}`$. The calibration uncertainty has to be included into these upper and lower limits. Three stars around this cluster have been spectroscopically investigated by Olszewski et al. (olszewski91 (1991)). Two have been identified, both being very red ($`(by)>1.3`$. They are also identified AGB stars (Frogel & Blanco frogel90 (1990)). These stars have a metallicity of $`1.33`$ dex and $`1.6`$ dex, a more metal rich ($`0.37`$ dex) AGB star could not be identified. NGC 1651 has been observed with the HST by Mould et al. (mould97 (1997)). They derived an age of $`10^{9.2\pm 0.1}`$ yr (using $`[Fe/H]=0.4`$ dex). The age is in good agreement despite the metallicity discrepancy. The elongated and tilted red clump of the cluster (see Fig. 24) is a feature that remains worth mentioning. The elongation is approximately along a reddening vector, however strong differential reddening should not cause this shape since the RGB does not show a similar large colour spread. We consider it more probable that this it is an intrinsic feature of an HB of a certain age and metallicity. The red clump of the field population does not show this elongated shape, it is rather a slightly fainter clump. Such an elongated red clump has also been observed by Piatti et al. (piatti99 (1999)) with in the Washington system in three of their $`21`$ investigated fields. In case of NGC 2209 they discuss the possibility that an increased helium content or differential reddening could cause such a red clump morphology. ### 11.2 The surrounding field population The dominating field population around NGC 1651 $`[Fe/H]=0.75`$ using the low reddening of $`E_{by}=0.04`$. The two-colour diagram for the field population is shown in Fig. 24 and the corresponding CMD in Fig. 24. The red clump of the field population is not elongated and slightly fainter. In this field, two groups of stars show up with distinct metallicities: one group with an approximately solar abundance ($`[Fe/H]=0.03\pm 0.03`$ dex) and one group around $`1`$ dex. The metal poor stars can be fitted with a Geneva isochrone of $`1`$ dex and an age on the order of $`10^{9.6}`$ yr. However, it is impossible to find a fitting isochrone for the apparent more metal rich stars (around $`0`$ dex), since these stars would have an age around $`10^{8.5}`$ yr and thus many more main sequence stars should be present. Only with the assumption of no reddening being present, one would have derived a metallicity of $`[Fe/H]=0.21\pm 0.19`$ dex for these stars. With such a metallicity they could have been fitted with a $`10^{8.8\pm 0.1}`$ yr Geneva isochrone. ## 12 NGC 2257 ### 12.1 The cluster The oldest cluster in our sample is NGC 2257, which can be seen from the cluster CMD (Fig. 26) that is very similar to the CMDs of Galactic globular clusters. Especially the pronounced blue horizontal branch (HB) is a sign for an old, metal poor population. NGC 2257 lies $`9^0`$ away from the centre of the LMC to the north east. Because of this large distance the field is very sparsely populated, thus we renounced a radial selection, because no radius can be found at which the field stars dominate. No reddening determination via a main sequence is possible. However, only a reddening of less than $`E_{by}=0.06`$ results in consistency of a reasonable age with a reasonable metallicity. We adopt here the reddening used by Testa et al. (testa95 (1995)), $`E_{BV}=0.04`$ and an error of $`\mathrm{\Delta }E_{BV}=0.04`$. Schwering & Israel’s map (schwering91 (1991)) shows a reddening of $`E_{BV}=0.03`$ at position of NGC 2257, Schlegel et al. schlegel98 (1998)) give $`E_{BV}=0.06`$ and Burstein & Heiles (burstein82 (1982)) $`E_{BV}=0.04`$. With a reddening of $`E_{BV}=0.04`$ the metallicity measurement results in $`[Fe/H]=1.63\pm 0.21`$ dex (including the calibration errors). In the two-colour diagram the metallicity determination is illustrated in Fig. 26. The age has been determined like in the previous cases, however since no Geneva isochrone with such an age and metallicity had been available we used Padua isochrones. The best matching (-1.7 dex) isochrone is overlaid on the CMD of this field in Fig. 26. The RGB of an old population evolves slowly and thus its location does not differ much between different ages. Therefore, we can only define an age range from $`10^{10}`$ yr to $`10^{10.3}`$ yr for this cluster, when limiting the age determination to the RGB. If the isochrones represent well the colour of the HB, then the age uncertainty will drop severely, since the HB of younger clusters is much redder (around $`(by)0.3`$) than that of older ones. In this case the allowed age range is merely $`10^{10.24}10^{10.30}`$ yr. The linear age uncertainty is much larger than for the younger clusters, discussed above. The reason for this is not an increased photometric error, but rather an intrinsic effect of the age-luminosity evolution: the brightness (in magnitudes) of the RGB decreases with age approximately logarithmically, therefore the error in the age of young and old stars should be in first order the same in the exponent. For the younger cluster this error was around 0.1, comparable to the error in case of NGC 2257. Testa et al. (testa95 (1995)) obtained deep B and V HST observations and concluded from the Turn-Off location and an assumed metallicity of -1.7 dex that the age is $`10^{10.1}`$ yr. Geisler et al. (geisler97 (1997)) found an age of $`10^{10.07}`$ yr. Recently, Johnson et al. (johnson99 (1999)) used deep HST observations of three true LMC globular clusters including NGC 2257 to derive their relative ages compared to Galactic globulars. They found that these old clusters have an age that is not distinguishable from that of M 3 and M 92. ### 12.2 The surrounding field population In spite of the low stellar density, a few field stars could have been identified due to their deviating age and metallicity. In the two-colour diagram (Fig. 26) stars having a metallicity around $`1`$ dex and being redder than $`by=0.6`$ have been marked with open star symbols. These stars form in the CMD (Fig. 26) a RGB lying below the cluster RGB. With a metallicity of $`[Fe/H]=1`$ dex the age of the population is $`10^{9.4}`$ yr. No other field component (except probably some Galactic foreground stars around $`(by)=0.5`$) have been found. The age and metallicity obtained for the field stars agrees well with the values found for the field population around the investigated clusters that are closer to the LMC centre. This indicates that, if a radial metallicity gradient exists, it cannot be very pronounced. This is consistent with the result of Olszewski et al. olszewski91 (1991) who did not find a radial metallicity gradient in the cluster population of the same age. However, there is a metallicity gradient in the sense that younger clusters and field stars tend to be more concentrated (Santos Jr. et al. santos99 (1999)) and thus the mean metallicity of all stars should show a metallicity gradient. In contrast Kontizias et al. kontizias93 (1993) found a metallicity gradient in the outer cluster system while none was observable for the inner clusters. The field population in the vicinity of NGC 2257 has been studied by Stryker (stryker84 (1984)). Her analysis revealed that the metallicity of the field population is larger than that of NGC 2257 and that “star formation occurred in the field long after the formation of the cluster”. She estimated the age of the field component to be $`67`$ Gyr old. We recalculated the age of this population based on her photographic CMD: we estimated the luminosity difference between the turn-off of this population and its red HB to be $`\mathrm{\Delta }V=2\pm 0.3`$. Walker et al. (walker99 (1993)) also found a younger field population around NGC 2257 with HST observations. The field population around this cluster is best fitted by a Padua isochrone with an age of $`3.5`$ gyr and a metallicity of $`0.6`$ dex. Using $`[Fe/H]=1`$ and applying the calibration given in Binney & Merrifield (binney98 (1998)) we end up with an age of $`4\pm 1`$ Gyr for this population which is comparable to our derived age. ## 13 The method to derive an AMR and SFH of the field population In order to extract detailed information concerning an AMR and/or the chemical enrichment history of the field, it is necessary to use a more sophisticated analysis than just deriving a mean metallicity and its standard deviation. Thus we estimated an age for each star with a measured metallicity. For these measurements a set of isochrones with different ages for each metallicity has to be available. Therefore we interpolated Geneva isochrones linearly to generate isochrones that are continuously distributed in metallicity. For each star the isochrone of the appropriate abundance with the minimal luminosity difference to the star has been identified and the corresponding age has been assigned to the star. This method obviously resulted in discrete age binning for the field stars (the size of the age bins is $`log(t)=0.1`$). We did not interpolate in age, since the age uncertainty due to the discrete age sampling is smaller than the error due to the photometric error, the calibration error and the errors connected to reddening and blending. Our method is not free from ambiguity, since for most stars one can derive two solutions, one for the RGB and one for the AGB. For stars older than $`10^{9.3}`$ yr one can neglect this effect because a) the luminosity difference between the AGB and RGB and thus the inferred age difference is small compared to the other uncertainties in determining an age for these evolved stars and b) the fraction of AGB to RGB stars in our colour range is small due to the lifetime difference. For younger ages the AGB stars play a considerable role: they dominate in the used colour range for populations with an age between $`10^{8.5}`$ yr and $`10^{9.1}`$ yr (see e.g. the overlayed isochrone in Fig. 12). To account for this problem we used the following approach: we determined an AGB age and a RGB age for each star. If the AGB and the RGB age were older than $`10^{9.3}`$ yr we assigned only a RGB age to the star. In case that we found a RGB age between $`10^{8.4}`$ yr and $`10^{8.7}`$ yr we used the AGB age to account for the dominance of the AGB stars in our used colour range. For the other ages we used either a mean AGB & RGB age, if none of the isochrones had a luminosity difference of less than $`0.1`$ mag or we used the age that correspond to the isochrone with a luminosity difference of less than $`0.1`$ mag to the star. The uncertainties in the extension of the isochrone towards the red is not critical for this investigation as long as one is not concerned with the number of stars with a certain age/metallicity. More severe are possible problems in the shape of the theoretical models, which might result in a systematic shift or distortion in the age scale, thus all our results on the field population is valid for the only currently available Strömgren isochrone set. However, since these isochrones fit reasonably well to the studied clusters, we are confident that the AMR is quite robust concerning the applied isochrones. We cannot circumvent this problem and it is necessary to get Strömgren isochrones for more recent stellar models with which the results can be compared. Since the age resolution on the RGB is not very good the derived age for an individual star is not more precise than a factor of $`2`$ for older stars, therefore all results can only be interpreted statistically. The applied procedure resulted in an AMR and an age number distribution (AND). The latter can be used to derive the Star Formation History (SFH) of the combined population. However, one has to bear in mind that the AND is not just the star formation rate (SFR) counted in logarithmic bins. First and most obvious is the fact, that stars of different ages have different masses on the RGB. Therefore one has to account for the IMF to get the SFH. Secondly, systematic shifts, for example due to differential reddening or binaries, have to be considered. Finally one has to be aware of the fact that the conclusions depend sensitively on the used set of isochrones. Taking all these effects into account is a highly complex problem and there is little hope to disentangle them analytically. To get nevertheless a handle on these effects and an idea about the accuracy of our method, we created synthetic CMDs of field populations with different ages and metallicities using a Monte Carlo algorithm and Geneva isochrones. The program allows one to include the (measured) errors, differential reddening, depth structure, a binary fraction and arbitrary SFHs. Binaries have been chosen randomly (according to a white random distribution) and are not important for the further investigation that regards only red giants. Even apart from the problems in interpreting the AND, the age distribution and thus the SFH depends rather strongly on the assumed reddening and possible differential reddening. Thus the results should be regarded more in comparison with the SFH derived with deep photometry (e.g. Gallagher et al. gallagher96 (1996), than as an independent measurement of the starformation history. Holtzman et al. holtzman99 (1999), Elson et al. elson97 (1997), Romaniello et al. romaniello99 (1999)) and serve as a consistency check. ### 13.1 Tests with simulated data To test our applied method we generated several artificial data sets. In Fig. 27a,b we show the resulting AMR and AND of two of these simulations. The left one (a) consists of three populations with $`10^{8.4}`$ yr, $`0.2`$ dex and $`10^{9.1}`$ yr, $`0.4`$ dex and $`10^{9.8}`$ yr, $`0.7`$ dex, respectively. The number ratio old-to-young-stars is $`1:1:10`$. No error was applied. In the right two panels (Fig. 27b) we have included the photometric errors, a binary fraction of 70% and differential reddening of $`\mathrm{\Delta }E_{BV}=0.04`$. The IMF used in both simulations (a Salpeter IMF down to $`0.8M_{}`$) is not important for the resulting AMR because of the small mass interval on the RGB. The open circles show the input age and metallicity, while filled circles are used for the extracted mean values for the metallicity. However, the IMF is an essential input parameter when deriving a SFH. Fortunately, it is not very critical for relative number ratios as long as the IMF did not change with time. This is an ad hoc assumption in our simulation since we cannot constrain any mass function with our method. The AMR of the input population in the shown simulation follow the AMR proposed by Pagel & Tautvais̆vienė (pagel98 (1998)). Thus this simulation demonstrates our ability to be able to recover the shape of the AMR proposed by these authors (on the basis of the given isochrones). It is impossible to see potential bursts in the SFH and it is clear that it is very difficult to derive the SFH from this procedure, only rough estimations of the SFH can be made. ### 13.2 Tests with the observed fields We tested our method also with the aid of the observed cluster stars. We used radially selected samples around the cluster centre and determined for them automatically the mean age and metallicity. The results are shown in Table 3 where also the error of the mean metallicity and the standard deviation ( $`\sigma `$(\[Fe/H\]) and $`\sigma `$(log(Age \[y\]))) is given for each cluster. The results for all clusters (except for NGC 2257, see below) are in good agreement with the ages that had been found with isochrone ”fitting”. For NGC 2257 we did not expect to find the cluster’s age and metallicity since or method applies only to stars having an age of less than $`10`$ Gyr and a metallicity of more than $`1.5`$ dex. The standard deviations around the mean values are considerable, however, systematically $`25\%`$ smaller than for the field populations. In Fig. 28 we show for two clusters one with a pronounced RGB and the other with a AGB the result of this method. We use the stars up to a distance of $`60\mathrm{}`$ and $`75\mathrm{}`$ around NGC 1806 and NGC 1651, respectively, as a combined input. The clusters can be seen as peaks at the corresponding age ($`10^{9.2}`$ yr, $`10^{8.7}`$ yr). The selection radius for NGC 1651 was two times larger than the one we used to derive the cluster’s age and metallicity via isochrone fitting. In Fig. 28 the age distribution of the field stars scaled to the same area as the cluster stars is plotted with the dotted line. ## 14 The AMR of the combined field population A major problem in deriving a reliable AMR and SFH is the small number statistic of field stars available for such an investigation after applying the selection criteria. To overcome the statistical problems, we have summed up all the CCD fields to work on a “global” LMC field population. With this approach, we are able to present an overall picture of the field AMR of the observed regions. Our sample of field stars (comprising 693 RGB stars, for which an age and a metallicity has been measured) enabled us to derive a ”global” AMR and AND, where ”global” rather means an average over our pointings. The distances of our inner clusters from the LMC centre are between $`1.6^0`$ and $`4^0`$. This corresponds to a projected distance difference of approximately $`2`$ kpc between the investigated fields. However, we did not weight these different fields and thus our results are more influenced by the stellar population around NGC 2031 and NGC 1806 than by for example the field stars around NGC 1651. Since after $`1`$ Gyr the stars should be well mixed within the LMC (Gallagher et al. gallagher96 (1996); they assumed $`1`$ km/sec as velocity dispersion of a typical LMC star). Thus this sample has a global meaning at least for the stars being older than $`1`$ Gyr. The CMD and two-colour diagram of this combined population is shown in Fig. 30 and Fig. 29, respectively. The derived AMR is shown in Fig. 31 and tabulated in Table7. The plotted “error” bars give the standard deviation in metallicity of stars with the same age and is not an error of the mean metallicity, since we do not believe that in each age bin all stars have the same metallicity even if we could measure the age with much higher accuracy. The age resolution on the giant branch drops considerably for stars that are older than $`10^{9.2}`$ yr because the spacing in luminosity between isochrones of different ages shrinks. This could cause the flat appearance of the AMR for these old ages, which is compatible with the shown ”error” bars. The upper age limit of our investigation is $`10^{9.9}`$ yr because no older Geneva isochrones were available, the lowest metallicity of the Geneva isochrone set is $`1.3`$ dex. For stars that are more metal poor we used the $`1.3`$ dex isochrones to derive the age (if the luminosity difference between the star and the isochrone was $`<0.1`$ mag), which means an underestimation. With these stars included, the oldest bin of the AMR contains stars, which in reality are older and more metal poor than the above stated limits. Therefore the number of stars in the age range $`10^{9.5}`$ yr - $`10^{9.9}`$ yr is slightly overestimated. The colour selection can introduce a bias, since the RGB of a metal rich population lies completely in the employed colour range, while for example only half of a $`1.3`$ dex RGB extends so far red which can be seen from the employed isochrones. The fraction of the RGB within the selected colour range is in good approximation independent of the age as long as stars older than $`1`$ Gyr are considered (concluded from visual inspection of the employed Geneva isochrones). For younger ages the problem becomes less severe, since the red supergiants extend further red. Since all our conclusions on the metallicity are based on simple means this introduces a bias towards higher metallicities if equal aged stars with different metallicity exist. To estimate the amount of the most extreme shift, we assume two populations with the same age but one having solar abundance and the other having a metallicity of $`1.3`$ dex. The difference in the mean metallicity, if only half of the metal poor giants are observed compared to the mean metallicity of the whole sample is $`\mathrm{\Delta }[Fe/H]_{mean}=1.3(1/21/3)=0.2`$ dex. Thus we conclude that deviations due to this problem are well less than $`0.2`$ dex. The CN anomaly leads to an overestimation of the metallicity and hence to an underestimation of the age. The deviation of the age according to the deviation of an overestimation of the photometric metallicity is nearly parallel to the observed AMR for ages larger than $`10^{8.5}`$ yr. Even with CN anomalous stars we should be able to distinguish between the proposed AMR and for example the AMR propposed by Pagel & Tautvais̆viennė (pagel98 (1998)): the influence of CN anomalous stars on the second AMR would result in an even more pronounced difference as it is already seen, since metal poorer stars would be shifted to even more metal rich and younger locations. ## 15 Discussion ### 15.1 The Age-Metallicity Relation In Table4 we summarise the resulting ages and metallicities of the investigated clusters. In order to compare our results with the literature, we compiled a list of clusters with ages and metallicities according to various sources (all published after 1989). The data is tabulated in Table9. In addition, we used the compilation of Sagar & Pandey (sagar89 (1989)), from which only clusters have been selected with a limiting magnitude below $`V=21`$. This limit shall serve as a rough quality criterion that is comparable to the more recent data and explains why most (photographic) papers cited by Sagar & Pandey are excluded. These clusters are plotted together with the newly investigated clusters and our field AMR in Fig. 32. The solid line is the field AMR accompanied by two dotted lines which mark 1$`\sigma `$ borders. If the older clusters are excluded, a weak correlation appears for the clusters: clusters younger $`1`$ Gyr have a mean metallicity of $`[Fe/H]=0.34`$ with a standard deviation of 0.14 and in the age range $`12.5`$ Gyr the mean metallicity is $`[Fe/H]=0.71`$ with a standard deviation of $`0.17`$ (11 cluster). The mean metallicity of our young clusters ($`<10^{9.0}`$ yr) is $`0.57\pm 0.04`$ dex and thus lower than what we found using the newer cluster sample from the literature. The field of the same age has a mean metallicity of $`0.4\pm 0.2`$ dex and is in good agreement with spectroscopic measurements ($`0.38\pm 0.11`$ dex) of young field stars. The latter comparison is reasonable since a) most of the stars for which high resolution spectroscopy has been obtained are located at a similar radial distance and b) no radial gradient can be seen in the spectroscopic sample. We compiled a list of high resolution spectroscopic measurements of LMC stars in Table 8. Bica et al. (bica98 (1998)) obtained ages and metallicities of 13 outer clusters in the LMC using Washington photometry. The mean metallicity of all the surrounding field stars is $`[Fe/H]0.6\pm 0.1`$. The mean metallicities of our field populations seem to be systematically more metal poor than this value thus indicating a possible zero point difference of the order of 0.2 dex and comparable to the probable shift between our cluster metallicities and the ones taken from the literature. However, such a difference between cluster and field stars has not been seen in a study by Korn et al. (korn00 (2000)) who employed high resolution spectroscopy of supergiants. The AMR for stars older than $`3`$ Gyr is consistent with little or even no enrichment until 8 Gyr ago. The AMR in this age range agrees well with the $`4`$ Gyr old clusters studied by Sarajedini (sarajedini98 (1998)) and also ESO 121 SC03 is in agreement with the derived field star AMR, especially when taking the systematic underestimation of the metallicity for older stars on the order of $`0.050.1`$ dex into account (see Sect. 3). Therefore, we do not see the necessity that ESO 121 SC03 belongs to a dwarf galaxy that is in the process of merging with the LMC as proposed by Bica et al. bica98 (1998)). The field population of NGC 1651 and NGC 2257 is considerably different from that around the other clusters: around NGC 1651 we find two distinct field populations, around NGC 2257 only one, thus these fields cannot be compared to the other fields, where a mixture of populations have been detected. These fields contain a significantly larger fraction of old stars than the other fields, what is expected from their location in the LMC (e.g. Santos Jr. et al. santos99 (1999)). If cluster and field are compared it becomes apparent that our AMR does not argue for an extremely decoupled enrichment history between cluster and field stars, only hints can be seen that the younger clusters are slightly more metal poor than the surrounding field population of the same age. Bica et al. (bica98 (1998)) found the same behaviour for several of their (young) clusters and the surrounding field population. One has to consider the possibility, that these low mean cluster abundances are a result of the statistically larger effect of blending towards the cluster. This has been proposed by Bessell (bessell93 (1993)) to explain the low Strömgren metallicity of NGC 330 measured by Grebel & Richtler (grebel92 (1992)). In this work of Grebel & Richtler (grebel92 (1992)) the mean metallicity found for the surrounding field population ($`0.74`$ dex) agreed well with later on performed spectroscopic measurements ($`0.69`$ dex, Hill hillV99 (1999)) ( for a more comprehensive discussion on NGC 330 the reader is referred to the work by Gonzalez & Wallerstein gonzalez99 (1999)). Having this agreement in mind, one can estimate from the difference of mean field and cluster metallicity that the contamination has a minor effect on our derived metallicities, accounting possibly for a systematic deviation of less than $`<0.15`$ dex. Our AMR is inconsistent with a recent calculation presented by Pagel & Tautvais̆vienė (pagel98 (1998)) based on LMC clusters and on planetary nebulae observed by Dopita et al. (dopita97 (1997)), that predicts a steeper increase of the metallicity in earlier time, thus older stars should have a higher metallicity than what we observe (see Fig. 32). The AMR is more consistent with closed box model calculations performed by Geha et al. (geha98 (1998)). They present theoretical enrichment models for the two SFHs put forward by Holtzman et al. (1997) and by Vallenari et al. (1996a,b). These SFHs agree in the sense, that a long period of low star formation activity was followed by a sudden increase about 2 Gyr ago. With the Vallenari et al.- SFH, the metallicity increased by a factor of five during the last 2 Gyr, while a modest increase of a factor of three resulted from the Holtzman et al.-SFH. Dopita et al. (dopita97 (1997)) published an AMR for the LMC based on planetary nebulae and found that the metallicity only doubled in the last 2-3 Gyr which is seen in our AMR as well. Another common feature is that a distinct enrichment (if any) between 4 and 9 Gyr cannot be seen. A comparison of the Dopita et al.-values with ours is made difficult by the fact that they measured $`\alpha `$-element abundances instead of $`[Fe/H]`$, but they stated that ”there is no evidence in this sample of any ”halo” abundance object”. If we would apply a constant shift of $`0.35`$ on the \[O/H\]-abundance, to correct approximately the \[O/Fe\] overabundance in the LMC in comparison to the Milky Way the metallicity of the PNs with an age of $`10^{8.8}10^{9.8}`$ yr would nicely be in coincidence with our measurements. However, the $`[O/Fe]`$ variation in dependence of $`[Fe/O]`$ is still under discussion (see e.g. Russell & Dopita russell92 (1992) or Pagel & Tautvais̆vienė pagel98 (1998)). Judging from the field around NGC 2257 we find that no radial metallicity gradient can be seen, since the field stars are consistent with the AMR derived from the inner fields. Taking also NGC 1651 into account we find that in fields where no recent star formation happened the stellar population is dominated by a population with an age between $`2`$ and $`4`$ Gyr. Thus deriving a global SFH on a limited sample is quite uncertain. #### 15.1.1 The Star Formation History The manner in which we derived the field star SFH contains several points that may induce biases. One reason is that the isochrones have only a crude spacing in the parameters age and metallicity. Therefore, simulations are helpful for a discussion of the SFH of the field population as described above (Sect. 12). To derive a SFH from our data is more difficult than to derive an AMR, since one has to know not only the age of a star with a given metallicity, but the amount of stars with a given age has to be quite precise. As a result the AMR is quite robust against for example reddening variations compared to the AND. Two SFHs are shown that illustrate how to interpret the AND (Fig. 33). One SFH has a constant SFR during the whole LMC evolution and thus serves to give an impression how the selection effects behave (left two panels in Fig. 33).The SFR of the second SFH was constant until $`10^{9.7}`$ yr ago, then it increased by a factor of $`5`$ until $`10^{8.4}`$ yr ago, before the SFR dropped to its old low level. The AND & AMR resulting from this SFH is plotted in the right panel of Fig. 33. The constant SFR is marginally inconsistent with our data, which holds for a different reddening correction of $`E_{by}=\pm 0.02`$. This is not true for exact behaviour of the SFH: for example a decrease in the reddening of $`E_{BV}=0.02`$ results in a SFH in which a much larger increase (around a factor of $`10`$) is necessary to describe the observations. However, the general trend, namely, the increase of the SFR around $`10^{9.5\pm 0.2}`$ yr ($`25`$ Gyr) ago and the necessary declining SFR some $`10^{8.5}`$ yr ago in these fields is more robust. Since stars in the LMC should be mixed (at least azimuthally) after $`1`$ Gyr (Gallager et al. gallagher96 (1996)) the SFH of the older stars should be a measure for the average SFH of the LMC in the radial distance of the investigated clusters. As a rule of thumb an increase in the applied reddening correction of $`E_{by}=0.1`$ results for a single age population in a decrease in age by a factor of $`0.7`$. Vallenari et al. (1996a,b) proposed, on the basis of ground based observations, a SFH in which the SFR increased about a factor of ten 2 Gyr ago, thus only around 5 % of the stars should be older than 4 Gyr. This has recently also been found by Elson et al. (elson97 (1997)) with HST observations. A different SFH was advanced by Holtzman et al. (holtzman97 (1997)), Geha et al. (geha98 (1998)) and Holtzman et al. (holtzman99 (1999)) also based on HST observations. In their model, approximately half of the stars are older than 4 Gyr. In our data the fraction of stars older than 4 Gyr is $`40\pm 20`$%, but we note that already a small additional reddening of of $`E_{by}=0.015`$ leaves only $`15`$% of the stars older than 4 Gyr. Olsen (olsen99 (1999)) used Washington photometry of the LMC field population and derived a SFH which is compatible with the one proposed by Holtzman et al. (holtzman99 (1999)). Summarizing, despite the uncertainty in the amount of the increase, the SFH is consistent with an increased SFR that started roughly $`3\pm 1`$ Gyr ago. Interestingly the sparsely populated outer fields are tentatively populated by mainly a population with ages between $`2`$ and $`4`$ Gyr. If this result will hold for a larger sample of outlying fields this could mean that the stellar body of the inner part of the LMC contains more younger and older stars compared to these intermediate age stars than the more remote parts of this galaxy. However, spectroscopic studies of several of these candidate stars are needed, especially because the tilted red clump of NGC 1651 could be due to a He overabundance and thus isochrones might be misleading. The observed SFH is inconsistent with a starformation history in which no star has been born between $`4`$ and $`8`$ Gyr. The simulations showed that virtually no star should have been recovered with an age of more than $`10^{9.3}`$ yr, even if the differential reddening is as large as $`E_{BV}=0.07`$ and the binary fraction is $`70\%`$. However, a large amount of old ($`>10`$ Gyr) CN anomal stars could mimic starformation between $`4`$ and $`8`$ Gyr. A last remark on the cluster formation rate: it has been noted several times that there apparently was a long period in the LMC where no clusters (or a few) have been formed. Recently, Larsen & Richtler (1999) performed a search for bright star clusters in 21 face-on galaxies. They found a correlation of the specific cluster frequency with parameters indicating the SFR. The age gap of the LMC cluster thus could reflect the low SFR during this period, where the condition for cluster formation where not present. ## 16 Summary We tried to determine the Age Metallicity Relation (AMR) and the Star Formation History (SFH) in the LMC on the basis of metallicities and ages of red giants, measured by Strömgren photometry. Our stars are located both in star clusters and in the respective surrounding fields. While statements regarding the AMR are relatively robust, the SFH is much more difficult to evaluate because of the incompleteness effects, for which we can only approximately correct. Between $`8`$ to $`3`$ Gyr ago the metallicity of the LMC was constant or varied only very slow, after this period the speed of the rate of the enrichment grew: starting $`3`$ Gyr ago the metallicity of the stars increased by a factor of six. The cluster and field AMR during this time was coupled, however a possibility remains that the cluster have on the average a slightly smaller metallicity than the field stars of the same age. Our field star AMR is also consistent with the $`4`$ Gyr old clusters recently studied by Sarajedini (sarajedini98 (1998)) and with ESO 121 SC03. Thus for the latter there is no need to explain this cluster as a recent merger remnant. Good agreement can be found with the photometrically determined metallicity of the young stars and the abundances measured with high resolution spectroscopy. The star formation rate increased around $`3`$ Gyr ago, however it is not possible to constrain the SFH further due to uncertainties in reddening, the CN anomaly and the difficult completeness considerations. ###### Acknowledgements. The authors gratefully acknowledge observing time at La Silla and the aid by the staff of the European Southern Observatory. We are grateful to E. Grebel and J. Roberts for the opportunity to use their isochrones. We thank the referee for his/her comments and suggestions which greatly helped to improve the paper. We also thank Doug Geisler and Antonella Vallenari for their valuable comments. BD was supported through the DFG Graduiertenkolleg ”The Magellanic Clouds and other dwarf galaxies” (GRK 118). MH thanks Fondecyt Chile for support through ‘Proyecto FONDECYT 3980032’. WPG gratefully acknowledges support received by Fondecyt grant No. 1971076. TR thanks the Uttar Pradesh State Observatory, Nainital, for warm hospitality and financial support. This research has made use of the Simbad database, operated at CDS, Strasbourg, France. ## Appendix A Appendix
warning/0004/astro-ph0004060.html	ar5iv	text	# Radiation-Hydrodynamical Collapse of Pregalactic Clouds in the Ultraviolet Background ## 1 Introduction It is widely recognized that the Ultraviolet (UV) background radiation, inferred from the proximity effect of Ly$`\alpha `$ absorption lines in QSO spectra (e.g., Bajtlik, Duncan & Ostriker 1988; Bechtold 1994; Giallongo et al. 1996), is likely to exert a significant influence upon the collapse of pregalactic clouds and consequently the formation of galaxies. Several authors have argued that the formation of subgalactic objects is suppressed via photoionization and photoheating caused by the UV background (Umemura & Ikeuchi 1984; Ikeuchi 1986; Rees 1986; Bond, Szalay & Silk 1988; Efstathiou 1992; Babul & Rees 1992; Zhang, Anninos & Norman 1995; Thoul & Weinberg 1996). Final states of the photoionized clouds, however, are still unclear, because most of these studies assume optically-thin media, which precludes us from the correct assessment of self-shielding against the external UV fields. Self-shielding not only plays an important role in the thermal and dynamical evolution, but is also essential for the formation of hydrogen molecules (H<sub>2</sub>) which control the star formation in metal poor environments such as primordial galaxies. Several attempts have recently been made to take account of the radiative transfer of ionizing photons, adopting for instance a pure absorption approximation (Kepner, Babul & Spergel 1997; Kitayama & Ikeuchi 2000, hereafter KI), the photon conservation method (Abel, Norman & Madau 1999), and the full radiative transfer treatment (Tajiri & Umemura 1998, hereafter TU; Barkana & Loeb 1999; Susa & Umemura 2000). As for the dynamical states, Kepner, Babul & Spergel (1997) and Barkana & Loeb (1999) have considered hydrostatic equilibria of spherical clouds within virialized dark halos. KI have explored the hydrodynamical evolution of a spherical system composed of dark matter and baryons. The full radiative transfer treatment, however, has not hitherto been incorporated with the hydrodynamics of spherical pregalactic clouds. In this paper, we attempt to implement an accurate radiation-hydrodynamical (RHD) calculation on the evolution of spherical clouds exposed to the UV background, solving simultaneously the radiative transfer of photons and the gas hydrodynamics. Our goals are to predict at a physically reliable level final states of photoionized clouds and also to assess accurately self-shielding against the UV background that is essential for H<sub>2</sub> cooling and the subsequent star formation. Throughout the present paper, we assume the density parameter $`\mathrm{\Omega }_0=1`$, the Hubble constant $`h=H_0`$/(100 km s<sup>-1</sup> Mpc$`{}_{}{}^{1})=0.5`$, and the baryon density parameter $`\mathrm{\Omega }_\mathrm{b}=0.1`$. ## 2 Model A pregalactic cloud is supposed to be a mixture of baryonic gas and dark matter with the mass ratio of $`\mathrm{\Omega }_\mathrm{b}:\mathrm{\Omega }_0\mathrm{\Omega }_\mathrm{b}=1:9`$. The numerical scheme for the spherical Lagrangian dynamics of two-component matter follows the method described in Thoul & Weinberg (1995) and KI. At each time-step, the radiative transfer is solved with the method devised by TU, in which both absorption and emission (scattering) of ionizing photons are explicitly taken into account. We assume for simplicity that the baryonic gas is composed of pure hydrogen. Neglecting helium causes only a minor effect in the ionization degree less than about the order of 10% (Osterbrock 1989; Nakamoto, Susa & Umemura 1998). We further assume ionization equilibrium among photoionization, collisional ionization, and recombination. This assumption is well justified in the present analysis (see §2.2 of KI for discussion). The number of mass shells is $`N_\mathrm{b}=200`$ for baryonic gas and $`N_\mathrm{d}=2000`$ for dark matter. At each radial point, angular integration of the radiative transfer equation is done over at least 20 bins in $`\theta =0\pi `$, where $`\theta `$ is the angle between the light ray and the radial direction. This is achieved by handling $`300700`$ impact parameters for light rays. The radiation field and the ionization states in the cloud interior are solved iteratively until the HI fraction, $`X_{\mathrm{HI}}`$, in each mesh converges within an accuracy of 1%. The external UV field is presumed to be isotropic and to have a power-law spectrum: $$J(\nu )=J_{21}\left(\frac{\nu }{\nu _{\mathrm{HI}}}\right)^\alpha \times 10^{21}\text{erg s}\text{-1}\text{ cm}\text{-2}\text{ str}\text{-1}\text{ Hz}\text{-1},$$ (1) where $`J_{21}`$ is the intensity at the Lyman edge of hydrogen ($`h\nu _{\mathrm{HI}}=13.6`$eV) and $`\alpha `$ is the spectral index. We consider two typical cases for $`\alpha `$, i.e., $`\alpha =1`$ representative for black hole accretion and $`\alpha =5`$ for stellar UV sources. Observations of the proximity effect in the Ly$`\alpha `$ forest suggest $`J_{21}=10^{\pm 0.5}`$ at $`z=1.74.1`$ (e.g., Bajtlik, Duncan & Ostriker 1988; Bechtold 1994; Giallongo et al. 1996), but its value is still uncertain at other redshifts. In what follows we give the onset of the UV background to be at $`z_{\mathrm{UV}}=20`$ and study the following two cases for $`J_{21}`$: 1. Constant UV $$J_{21}=1zz_{\mathrm{UV}},$$ (2) 2. Evolving UV $$J_{21}=\{\begin{array}{cc}\mathrm{exp}[(z5)]\hfill & 5zz_{\mathrm{UV}}\hfill \\ 1\hfill & 3z5\hfill \\ \left(\frac{1+z}{4}\right)^4\hfill & 0z3.\hfill \end{array}$$ (3) The form of the UV evolution at $`z>5`$ in (ii) is roughly consistent with the results of recent models for the reionization of the universe (e.g., Ciardi et al 2000; Umemura, Nakamoto & Susa 2000). We start the simulations when the overdensity of a cloud is still in the linear regime, adopting the initial and boundary conditions described in KI. The initial overdensity profile is $`\delta _\mathrm{i}(r)=\delta _\mathrm{i}(0)\mathrm{sin}(kr)/kr`$, where $`k`$ is the comoving wave number, and the central overdensity $`\delta _\mathrm{i}(0)`$ is fixed at 0.2. The outer boundary is taken at the first minimum of $`\delta _\mathrm{i}(r)`$, i.e., $`kr=4.4934`$, within which the volume averaged overdensity $`\overline{\delta }(<r)`$ vanishes. Following Haiman, Thoul & Loeb (1996), the characteristic mass of a cloud $`M_{\mathrm{cloud}}`$ is defined as the baryon mass enclosed within the first zero of $`\delta _i(r)`$, i.e., $`kr=\pi `$. Collapse redshift $`z_\mathrm{c}`$ is defined as the epoch at which $`M_{\mathrm{cloud}}`$ would collapse to the center in the absence of thermal pressure. Circular velocity $`V_\mathrm{c}`$ and virial temperature $`T_{\mathrm{vir}}`$ are related to $`z_\mathrm{c}`$ and $`M_{\mathrm{cloud}}`$ via usual definitions: $$V_c=15.9\left(\frac{M_{\mathrm{cloud}}\mathrm{\Omega }_0/\mathrm{\Omega }_b}{10^9h^1\mathrm{M}_{}}\right)^{1/3}(1+z_c)^{1/2}\text{ km s}\text{-1},$$ (4) $$T_{\mathrm{vir}}=9.09\times 10^3\left(\frac{\mu }{0.59}\right)\left(\frac{M_{\mathrm{cloud}}\mathrm{\Omega }_0/\mathrm{\Omega }_b}{10^9h^1\mathrm{M}_{}}\right)^{2/3}(1+z_c)\text{ K},$$ (5) where $`\mu `$ is the mean molecular weight in units of the proton mass $`m_\mathrm{p}`$. The collapse of a gas shell is traced until it reaches the rotation radius specified by the dimensionless spin parameter; $$r_{\mathrm{rot}}=0.05\left(\frac{\mathrm{\Omega }_\mathrm{b}/\mathrm{\Omega }_0}{0.1}\right)^1\left(\frac{\lambda _{\mathrm{ta}}}{0.05}\right)^2r_{\mathrm{ta}},$$ (6) where $`r_{\mathrm{ta}}`$ is the turnaround radius of the gas shell, and we adopt a median for the spin parameter, $`\lambda _{\mathrm{ta}}=0.05`$ (Efstathiou & Jones 1979; Barns & Efstathiou 1987; Warren et al. 1992). Below the size given by (6), the system would attain rotational balance and forms a disk eventually. ## 3 Results ### 3.1 Significance of radiative transfer To demonstrate the significance of coupling the radiative transfer of ionizing photons with hydrodynamics, we first present in Fig. 1 estimations for relevant timescales at the center of a uniform static cloud with circular velocity $`V_\mathrm{c}`$ and collapse epoch $`z_\mathrm{c}`$, where the external UV is specified by $`J_{21}`$ and a spectral index $`\alpha `$. The photoionization timescale $`t_{\mathrm{ion}}`$ and the photoheating timescale $`t_{\mathrm{heat}}`$ are compared to the dynamical timescale $`t_{\mathrm{dyn}}`$, either by solving the radiative transfer of UV photons through the cloud or by just assuming the optically thin medium. Wherever necessary, we have adopted the temperature $`T=10^4`$K, the proton number density $`n_\mathrm{H}=n_\mathrm{H}^{\mathrm{vir}}=5.0\times 10^5(\mathrm{\Omega }_\mathrm{b}h^2/0.025)(1+z_\mathrm{c})^3`$ cm<sup>-3</sup>, the mass density $`\rho =m_pn_\mathrm{H}\mathrm{\Omega }_0/\mathrm{\Omega }_b`$, and the radius $`R=R_{\mathrm{vir}}=64(V_\mathrm{c}/30\text{km s}\text{-1})(\mathrm{\Omega }_0h^2/0.25)^{1/2}(1+z_\mathrm{c})^{3/2}`$ kpc. In the low-density limit, the contours of $`t_{\mathrm{ion}}=t_{\mathrm{dyn}}`$ and $`t_{\mathrm{heat}}=t_{\mathrm{dyn}}`$ both approach asymptotically the optically thin case $`J_{21}\sqrt{n_\mathrm{H}^{\mathrm{vir}}}(1+z_c)^{3/2}`$. The differences originated from the cloud sizes are rather small compared to those arisen when one incorporates the radiative transfer or not. Fig. 1 predicts that the significance of radiative transfer effects increases with increasing redshift (i.e., increasing cloud density) and decreasing $`J_{21}`$ for a given $`\alpha `$. Under the UV evolution given in equation (3), for example, $`t_{\mathrm{dyn}}<t_{\mathrm{ion}}`$ and $`t_{\mathrm{dyn}}<t_{\mathrm{heat}}`$ are satisfied at the center of a cloud with $`V_c=30`$km s<sup>-1</sup> at $`z_\mathrm{c}\stackrel{>}{}10`$, and the cloud is shielded against the external radiation in terms of photoionization as well as photoheating. At $`6\stackrel{<}{}z_\mathrm{c}\stackrel{<}{}10`$, $`t_{\mathrm{heat}}<t_{\mathrm{dyn}}<t_{\mathrm{ion}}`$ is achieved, indicating that the cloud center is heated but not ionized by the UV radiation (Gnedin & Ostriker 1997; KI). Contrastively, both $`t_{\mathrm{ion}}`$ and $`t_{\mathrm{heat}}`$ become shorter than $`t_{\mathrm{dyn}}`$ at $`z_\mathrm{c}<6`$ and the cloud is likely to evolve in a similar fashion to the optically-thin case. We can see qualitatively the same relations for $`\alpha =5`$, except that shielding of ionization and heating occurs almost simultaneously at lower redshift. It should be noted that the above estimations are made for a virialized cloud using the UV intensity only at $`z_c`$. As shown in forthcoming sections, the central density of a collapsing cloud actually continues to ascend to above $`n_\mathrm{H}^{\mathrm{vir}}`$ and the radiative transfer effects can become important even at low redshifts. In addition, possible changes of the UV intensity during the dynamical growth of a cloud can also affect its ionization structure. ### 3.2 Dynamical evolution Fig. 2 shows the dynamical evolution of a cloud with $`V_\mathrm{c}=29`$km s<sup>-1</sup> under the constant UV background at $`z<z_{\mathrm{UV}}=20`$. This cloud would collapse at $`z_\mathrm{c}=3`$ if there were no UV background. In practice, the cloud turns around and contracts in the inner parts, while it continues to expand in the outer envelope. For a hard UV spectrum with $`(J_{21},\alpha )=(1,1)`$, the gas is ionized and heated up to $`T10^4`$K promptly at the onset of the UV background. For a soft UV spectrum with $`(J_{21},\alpha )=(1,5)`$, the cloud center is kept self-shielded against the external field and the temperature ascends more gradually. For comparison, results of the optically-thin and pure absorption calculations are also presented in Fig. 2. The pure absorption case is based on the analytical formalism described in KI. The cloud evolution is altered in no small way by different treatments of the UV radiation. Under the optically-thin assumption, the cloud is completely prohibited from collapsing for $`(J_{21},\alpha )=(1,1)`$, and from being self-shielded for $`(J_{21},\alpha )=(1,5)`$. The pure absorption approximation leads to accelerating the collapse of the central core and to underestimating photoionization and photoheating. Fig. 3 exhibits the radial profiles of the same clouds at $`z=3`$ for $`(J_{21},\alpha )=(1,5)`$. The inner part of a cloud does not settle into hydrostatic equilibrium, but rather undergoes isothermal run-away collapse and the density profile follows the self-similar solution of Bertschinger (1985). The trend is regardless of the treatment of the UV radiation transfer, because the temperature of thermal equilibrium in optically-thin media comes close to $`10^4`$K in the high density limit. The present results demonstrate that the approximation of hydrostatic equilibrium employed in previous analyses (Kepner, Babul & Spergel 1997; Barkana & Loeb 1999) breaks down for the final states. The ionization structure is quite different depending on whether one incorporates the radiative transfer or not. In an optically-thin cloud, the ionization degree changes merely according to the ionization parameter, $`n_\mathrm{H}/J_{21}`$, whereas the radiative transfer effects produce a self-shielded neutral core (Fig. 3c). Fig. 3(d) further indicates that the UV heating rate is reduced significantly in the self-shielded region, with an increasing contribution of scattered photons to the total heating rate. Implications of the present results on H<sub>2</sub> cooling \[thin lines in panel (d)\] will be discussed in detail in Sec 4. Effects of the evolution of the UV background are illustrated in Fig. 4. Due to very low UV intensity at high redshift, the whole cloud is kept self-shielded in terms of both photoionization and photoheating at $`z\stackrel{>}{}12`$. For a cloud with relatively high $`z_c`$ ($`z_c=4.8`$, Fig. 4a), the cloud center begins to contract before the penetration of the external UV. Hence, the dynamical evolution closely coincides with that without the UV background. As the virial temperature of the cloud is $`T_{\mathrm{vir}}3\times 10^4`$K, the could center is shock-heated to above $`10^4`$K and cools via atomic cooling. On the other hand, a cloud with low $`z_c`$ ($`z_c=0.5`$, Fig. 4b) is once photoionized to the similar level to the constant UV case at $`3<z<5`$, but is able to collapse as the UV intensity drops at lower redshifts. ### 3.3 Criteria for collapse Figs 5 and 6 summarize the results of present calculations for a variety of initial conditions on a $`V_\mathrm{c}z_\mathrm{c}`$ plane. As our simulations assume ionization equilibrium and only incorporate atomic cooling, they can be most securely applied to a cloud once heated to above $`10^4`$K either by shock or by UV photons during the course of its evolution. In contrast, evolution of lower temperature systems is still tentative and may be altered once cooling by molecular hydrogen is explicitly taken into account. These figures thus distinguish “high temperature clouds” (circles) defined as those photoheated to $`>10^4`$K or those with $`T_{\mathrm{vir}}>10^4`$K, and the other “low temperature clouds” (triangles). Each of these populations are further classified by open and filled symbols depending on whether or not they collapse to the rotation barrier within the present age of the universe. It is obvious that the UV background prohibits small clouds from collapsing even if the transfer of ionizing photons is considered. Under a constant UV flux (Fig. 5), threshold circular velocity for the collapse gradually increases with decreasing redshift, i.e., decreasing cloud density. The threshold velocity also decreases with increasing photon spectral index $`\alpha `$ because of the smaller number of high energy photons. At $`z_c\stackrel{>}{}5`$, the threshold falls below $`10^4`$K, because the cloud center starts to contract before the onset of the UV background and remains impervious to the external photons. In the presence of the UV evolution (Fig. 6), the threshold velocity increases sharply at $`z_c\stackrel{<}{}3`$ and drops slightly at $`z_c\stackrel{<}{}1`$. The central temperature of a cloud denoted by an open triangle can reach $`10^4`$K by adiabatic compression and the collapse is promoted by atomic cooling. A cloud denoted by a filled triangle fails to collapse because of the lack of coolant at $`T<10^4`$K in our simulations. As will be discussed in Sec 4, however, evolution of these “low temperature clouds” may be modified by H<sub>2</sub> cooling. KI have suggested that the threshold for collapse is roughly determined by the balance between the gravitational force and the thermal pressure gradient when the gas is maximally exposed to the external UV flux. To confirm this, we plot in the same figures a relation $`T_{\mathrm{vir}}=T_{\mathrm{eq}}^{\mathrm{max}}`$, where $`T_{\mathrm{eq}}^{\mathrm{max}}`$ is the maximum equality temperature defined semi-analytically as follows. Firstly, given the initial overdensity profile, the collapse of a spherical perturbation is approximated by the self-similar solution of Bertschinger (1985). Secondly, at an arbitrary stage of the collapse, one can compute the equality temperature at which radiative cooling balances photoheating in the optically thin limit, using the central density deduced from the self-similar solution and the UV intensity assumed in the simulation. Finally, $`T_{\mathrm{eq}}^{\mathrm{max}}`$ is set equal to the maximum of such temperature. For the gas exposed to a constant UV flux from the linear regime, $`T_{\mathrm{eq}}^{\mathrm{max}}`$ is attained essentially at turn-around. If a cloud has $`T_{\mathrm{vir}}`$ above $`T_{\mathrm{eq}}^{\mathrm{max}}`$, it is likely to be gravitationally unstable against gas pressure. Figs 5 and 6 show that the above criteria agree reasonably well with our simulation results based on the accurate treatment of the radiative transfer. This is because the dynamical evolution is basically regulated by the Jeans criterion when a cloud is heated up to $`10^4`$K by the UV background before contraction. At high redshifts ($`z_c\stackrel{>}{}6`$ in Fig. 5 and $`z_c\stackrel{>}{}3`$ in Fig. 6), however, the cloud center is strongly self-shielded from early stages and the approximation of the optically thin limit breaks down. The Jeans scale at these epochs is reduced below the relation $`T_{\mathrm{vir}}=T_{\mathrm{eq}}^{\mathrm{max}}`$. ## 4 Implications for Galaxy Formation in the UV Background The present RHD calculations give an accurate prediction for the suppression of pregalactic collapse due to the UV background, which has been one of the primary concerns from the viewpoint of galaxy formation (Umemura & Ikeuchi 1984; Ikeuchi 1986; Rees 1986; Bond, Szalay & Silk 1988; Efstathiou 1992; Babul & Rees 1992; Zhang, Anninos & Norman 1995; Thoul & Weinberg 1996; KI). What is of additional significance in this context is the subsequent formation of stars in collapsing clouds. In order for stars to form, a cloud needs to be cooled down to below $`10^4`$K by hydrogen molecules, because they are the only coolant in metal-deficient gas (e.g., Peebles & Dicke 1968; Matsuda, Sato & Takeda 1969; Tegmark et al. 1997). H<sub>2</sub> cooling is a two-body collision process (e.g., Hollenbach & Mckee 1989; Galli & Palla 1998), while the photoheating rate is in proportion to the density. The potentiality of H<sub>2</sub> cooling is therefore enhanced with increasing density. In this respect, run-away collapse should provide favorable situations for H<sub>2</sub> cooling. Based on our simulation results presented in the previous section, we further investigate the possibility of star formation in a collapsing core. We suppose that cooling below $`10^4`$K becomes efficient if the following two conditions are both satisfied; 1) photodissociation of molecular hydrogen by the UV photons in the Lyman-Werner bands at $`11.2613.6`$keV (e.g., Stecher & Williams 1967) is less important than other H<sub>2</sub> destruction processes, and 2) H<sub>2</sub> cooling overtakes UV photoheating. The first condition is fulfilled for clouds heated up to $`>10^4`$K, because destruction of H<sub>2</sub> is dominated by collisional dissociation insofar as $`T\begin{array}{c}>\hfill \\ \hfill \end{array}2000`$K (Corbelli, Galli & Palla 1997). More specifically, requiring that the timescale of H<sub>2</sub> dissociation via collisions with H<sup>+</sup> (reaction 12 of Shapiro & Kang 1987) is shorter than that of photodissociation (Draine & Bertoldi 1996; Omukai & Nishi 1999) in the conservative optically thin limit, we have for the electron density $$n_\mathrm{e}>4.7\times 10^5\left(\frac{13.6}{12.4}\right)^\alpha J_{21}\mathrm{exp}(21200\text{ K}/T)\text{ cm}\text{-3},$$ (7) which is amply satisfied at the collapsing core of our simulated clouds. The second condition depends on the amount of H<sub>2</sub> formed. In the absence of the external UV field, the H<sub>2</sub> abundance in the metal-deficient postshock layer converges roughly to $`X_{\mathrm{H2}}10^3`$ (e.g., Shapiro & Kang 1987; Ferrara 1998; Susa et al. 1998). Under the UV field, $`X_{\mathrm{H2}}10^3`$ is also achieved if photoheating is strongly attenuated by self-shielding, while the abundance is reduced down to $`X_{\mathrm{H2}}10^6`$ in the case of weak attenuation of photoheating (Kang & Shapiro 1992; Corbelli et al. 1997; Susa & Umemura 2000; Susa & Kitayama 2000). It is also likely that the H<sub>2</sub> abundance depends on the spectrum of the impinging radiation and the details of the ionizing structure of the medium (e.g., Kang & Shapiro 1992; Ciardi et al. 2000). Leaving such complexities elsewhere (Kitayama et al. in preparation), we make crude estimations for the cooling rate using the H<sub>2</sub> cooling function of Hollenbach & Mckee (1989) and assuming $`X_{\mathrm{H2}}=10^3`$ in the full transfer and pure absorption cases, and $`X_{\mathrm{H2}}=10^6`$ in the optically thin case. Fig. 3(d) shows that H<sub>2</sub> cooling greatly overwhelms photoheating in the self-shielded neutral core for $`X_{\mathrm{H2}}=10^3`$. It turns out that what enables H<sub>2</sub> cooling to be effective is essentially the attenuation of photoheating rather than the final abundance of H<sub>2</sub> molecules. In fact, in the strongly self-shielded core, H<sub>2</sub> cooling remains to be dominant even if the abundance is reduced to a level similar to the optically-thin case, i.e., $`X_{\mathrm{H2}}=10^6`$ (the H<sub>2</sub> cooling rate is roughly proportional to $`X_{\mathrm{H2}}`$ at $`T\stackrel{<}{}10^4`$ K). We have further checked that this is the case with every collapsed cloud plotted as an open circle in Figs 5 and 6. Under the optically-thin approximation, on the other hand, UV heating is much stronger than H<sub>2</sub> cooling, prohibiting the cooling below $`10^4`$K. While the above arguments support the possibility of H<sub>2</sub> cooling during the collapse of “high temperature clouds” (open circles in Figs 5 and 6), the situation is rather intricate for “low temperature clouds” (triangles in the same figures). Once incorporated formation and destruction of molecular hydrogen explicitly in our calculations, these clouds would be able to collapse as long as the external UV flux is negligible or very weak. For somewhat stronger UV flux, photodissociation may operate efficiently to suppress H<sub>2</sub> formation and to prohibit the collapse (Haiman, Rees & Loeb 1996, 1997; Haiman, Abel & Rees 2000; Ciardi et al. 2000). Alternatively, self-shielding induced by run-away collapse may still enable a sufficient amount of H<sub>2</sub> to be formed for the system to cool down to $`100`$K. These points will be investigated thoroughly in our future work (Kitayama et al. in preparation). In summary, self-shielding against the UV background is indispensable for H<sub>2</sub> cooling and subsequent star formation to proceed. Run-away collapse is likely to promote clouds with $`T_{\mathrm{vir}}>10^4`$K to cool down before collapsing to the rotation barrier given by equation (6). These results are closely relevant to the formation of dwarf galaxies at high redshifts as well as to metal injection into intergalactic space (e.g., Nishi & Susa 1999). Note that the above conclusions do not conflict with those of several recent works (Haiman Rees & Loeb 1996, 1997; Ciardi et al. 2000; Haiman et al. 2000) that the so called ”radiative feedback” is operating in the early universe and suppresses the collapse of some objects via photodissociation of H<sub>2</sub> molecules. Current simulations are mainly aimed to study the clouds once heated to above $`10^4`$K either by shock or by UV photons, for which photodissociation has minor impacts on H<sub>2</sub> formation compared to collisional dissociation (eq. ). In addition, run-away collapse yields a highly self-shielded core which is particularly favorable for H<sub>2</sub> formation. Photodissociation may be of greater significance in smaller clouds and lower density regions. Finally, we discuss the effects of the geometry of cloud evolution upon H<sub>2</sub> cooling. Near and above the Jeans scale, clouds are expected to contract spherically to a first order approximation and result in run-away collapse as shown in the present paper. Clouds far above the Jeans scale, on the other hand, are likely to undergo pancake collapse. The pancakes would not end up with run-away collapse, because the gravity of the sheet is readily overwhelmed by the pressure force. Susa & Umemura (2000) have explored pregalactic pancake collapse including the UV transfer and H<sub>2</sub> formation, and shown that the collapsing pancakes bifurcate, depending upon the initial masses, into less massive UV-heated ones and more massive cooled ones. This is because the self-shielding in sheet collapse is governed by the column density of a pregalactic cloud. To conclude, star formation in pregalactic clouds in the UV background is strongly regulated by the behaviour of collapse and the manner of the radiation transfer. ## Acknowledgments We thank Andrea Ferrara for helpful comments, Taishi Nakamoto for discussions, and Susumu Inoue for careful reading of the manuscript. This work is supported in part by Research Fellowships of the Japan Society for the Promotion of Science for Young Scientists, No. 7202 (TK) and 2370 (HS).
warning/0004/astro-ph0004393.html	ar5iv	text	# New CMB constraints on the cosmic matter budget: trouble for nucleosynthesis? \[ ## Abstract We compute the joint constraints on ten cosmological parameters from the latest CMB measurements. The lack of a significant second acoustic peak in the latest Boomerang and Maxima data favors models with more baryons than Big Bang nucleosynthesis predicts, almost independently of what prior information is included. The simplest flat inflation models with purely scalar scale-invariant fluctuations prefer a baryon density $`0.022<h^2\mathrm{\Omega }_\mathrm{b}<0.040`$ and a total nonbaryonic (hot + cold) dark matter density $`0.14<h^2\mathrm{\Omega }_{dm}<0.32`$ at 95% confidence, and allow reionization no earlier than $`z30`$. preprint: IASSNS-AST 97/666 \] One of the main challenges in modern cosmology is to refine and test the standard model of structure formation by precision measurements of its free parameters. The cosmic matter budget involves at least the four parameters $`\mathrm{\Omega }_\mathrm{b}`$, $`\mathrm{\Omega }_{\mathrm{cdm}}`$, $`\mathrm{\Omega }_\nu `$ and $`\mathrm{\Omega }_\mathrm{\Lambda }`$, which give the percentages of critical density corresponding to baryons, cold dark matter, massive neutrinos and vacuum energy. A “budget deficit” $`\mathrm{\Omega }_\mathrm{k}1\mathrm{\Omega }_\mathrm{b}\mathrm{\Omega }_{\mathrm{cdm}}\mathrm{\Omega }_\nu \mathrm{\Omega }_\mathrm{\Lambda }`$ manifests itself as spatial curvature. The description of the initial seed fluctuations predicted by inflation requires at least four parameters, the amplitudes $`A_s`$ & $`A_t`$ and slopes $`n_s`$ & $`n_t`$ of scalar and tensor fluctuations, respectively. Finally, the optical depth parameter $`\tau `$ quantifies when the first stars or quasars reionized the Universe and the Hubble parameter $`h`$ gives its current expansion rate. During the past year or so, a number of papers have used the measured cosmic microwave background (CMB) fluctuations to constrain subsets of these parameters. CMB data has improved dramatically since fluctuations were first detected . The measurement of a first acoustic peak at the degree scale , suggesting that the Universe is flat ($`\mathrm{\Omega }_\mathrm{k}=0`$), has now been beautifully confirmed and improved by using the ground-breaking high fidelity maps of the Boomerang and Maxima experiments. As can be seen in Figure 1, perhaps the most important new information from Boomerang and Maxima is their accurate measurements of the angular power spectrum $`C_{\mathrm{}}`$ on even smaller scales, out to multipole $`ł600800`$. The striking lack of a significant second acoustic peak places strong constraints on the cosmological parameters , making a new full-fledged analysis of all the CMB data very timely. In this Letter, we jointly constrain the following 10 cosmological parameters: $`\tau `$, $`\mathrm{\Omega }_\mathrm{k}`$, $`\mathrm{\Omega }_\mathrm{\Lambda }`$, $`n_s`$, $`n_t`$, $`A_s`$, the tensor-to scalar ratio $`rA_t/A_s`$, and the physical matter densities $`\omega _\mathrm{b}h^2\mathrm{\Omega }_\mathrm{b}`$, $`\omega _{\mathrm{cdm}}h^2\mathrm{\Omega }_{\mathrm{cdm}}`$ and $`\omega _\nu h^2\mathrm{\Omega }_\nu `$. The identity $`h=\sqrt{(\omega _{\mathrm{cdm}}+\omega _\mathrm{b}+\omega _\nu )/(1\mathrm{\Omega }_\mathrm{k}\mathrm{\Omega }_\mathrm{\Lambda })}`$ fixes the Hubble parameter. We use the 10-dimensional grid method described in . In essence, this utilizes a technique for accelerating the the CMBfast package by a factor around $`10^3`$ to compute theoretical power spectra on a grid in the 10-dimensional parameter space, fitting these models to the data and then using cubic interpolation of the resulting 10-dimensional likelihood function to marginalize it down to constraints on individual or pairs of parameters. We use the 87 data points shown in Figure 1, combining the 65 tabulated in with the 12 new Boomerang points and the 10 new Maxima points . Our 95% confidence limits on the best constrained parameters are summarized in Table 1. Figure 2 shows that CMB alone suggests that the Universe is either flat (near the diagonal line $`\mathrm{\Omega }_\mathrm{m}+\mathrm{\Omega }_\mathrm{\Lambda }=1`$, where $`\mathrm{\Omega }_\mathrm{m}\mathrm{\Omega }_\mathrm{b}+\mathrm{\Omega }_{\mathrm{cdm}}+\mathrm{\Omega }_\nu `$) or closed (upper right). These constraints come largely from the location of the first peak, which is well-known to move to the right if the curvature $`\mathrm{\Omega }_\mathrm{k}`$ is increased . Very closed models work only because the first acoustic peak can also be moved to the right by increasing the tilt $`n_s`$ or decreasing the matter density and bringing the large-scale COBE signal back up with tensor fluctuations (gravity waves) . Galaxy clustering constraints disfavor such strong blue-tilting, and Figures 2 and 3 show that closing this loophole by barring gravity waves ($`r=0`$) favors curvature near zero and $`n_s`$ near unity. This is a striking success for the oldest and simplest inflation models, which make the three predictions $`r0`$, $`\mathrm{\Omega }_\mathrm{k}0`$ and $`n_s1`$ . Another important success for inflation is that the first peak is so narrow — if the data had revealed the type of broad peak expected in many topological defect scenarios, none of the models in our grid would have provided an acceptable fit. Because of these tantalizing hints that “back to basics” inflation is correct, Table 1 and Figure 3 include results assuming this inflation prior $`r=\mathrm{\Omega }_\mathrm{k}=0,n_s=1`$. Table 1 – Maximum-likelihood values and 95% confidence limits. The “inflation prior” for each parameter is indicated in boxes in Figure 3. $`\mathrm{\Omega }_{dm}\mathrm{\Omega }_{\mathrm{cdm}}+\mathrm{\Omega }_\nu `$. A dash indicates that no limit was found, with the likelihood still above $`e^2`$ at the edge of our grid. Extrapolation would suggest a limit $`n_s\mathrm{}<1.75`$. \| \| 10 free parameters \| \| \| Inflation prior \| \| \| \| --- \| --- \| --- \| --- \| --- \| --- \| --- \| \| Quantity \| Min \| Best \| Max \| Min \| Best \| Max \| \| $`\tau `$ \| 0.0 \| 0.0 \| 0.33 \| 0.0 \| 0.0 \| 0.28 \| \| $`h^2\mathrm{\Omega }_\mathrm{b}`$ \| .017 \| .05 \| $``$ \| .022 \| .03 \| .040 \| \| $`h^2\mathrm{\Omega }_{dm}`$ \| 0.02 \| $`0.08`$ \| $``$ \| 0.14 \| 0.20 \| 0.32 \| \| $`\mathrm{\Omega }_\mathrm{\Lambda }`$ \| $``$ \| 0.2 \| 0.80 \| $``$0.16 \| .43 \| 0.65 \| \| $`\mathrm{\Omega }_\mathrm{k}`$ \| $``$ \| $``$0.6 \| 0.13 \| $``$0.13 \| 0 \| 0.10 \| \| $`n_s`$ \| 0.8 \| 1.5 \| $``$ \| 0.84 \| 1.0 \| 1.17 \| The constraints in Table 1 are seen to be much more interesting than those before Boomerang and Maxima , thanks to new information on the scale of the second peak and beyond. Cold dark matter and neutrinos have indistinguishable effects on the CMB except for very light neutrinos (small $`\omega _\nu `$), and the current data still lacks the precision to detect this subtle difference. The predicted height ratio of the first two peaks therefore depends essentially on only three parameters : $`n_s`$, $`\omega _\mathrm{b}`$ and $`\omega _{dm}`$, where the total dark matter density $`\omega _{dm}\omega _{\mathrm{cdm}}+\omega _\nu `$. Let us focus on the constraints on these parameters. Increasing $`\omega _\mathrm{b}`$ tends to boost the odd-numbered peaks (1, 3, etc.) at the expense of even ones (2, 4, etc.) , whereas increasing $`\omega _{dm}`$ suppresses all peaks (see the CMB movies at $`www.hep.upenn.edu/max`$ or $`www.ias.edu/whu`$). The low second peak can therefore be fit by either decreasing the tilt $`n_s`$ or by increasing the baryon density $`\omega _\mathrm{b}`$ compared to the usually assumed values $`n_s1`$, $`\omega _\mathrm{b}0.02`$. As illustrated in Figure 4, this conclusion is essentially independent of what priors are assumed. However, reducing $`n_s`$ below 0.9 is seen to make things worse again, as the first peak becomes too low relative to the COBE-normalization. In short, there are two very simple ways of explaining the lack of a prominent second acoustic peak: more baryons or a red-tilted spectrum . However, as we will now discuss, both of these solutions have problems of their own. Figure 5 shows that when more baryons are added, more dark matter is needed to keep the first peak height constant. When the tilt $`n_s`$ is fixed by the inflation prior, the constraints on the remaining two parameters $`\omega _{dm}`$ and $`\omega _\mathrm{b}`$ are seen to become quite tight. Intriguingly, the preferred baryon fraction is of the same order as preferred by Big Bang nucleosynthesis, but nonetheless higher than the tight nucleosynthesis error bars $`\omega _\mathrm{b}=0.019\pm 0.0024`$ allow. Even if the nucleosynthesis error bars have somehow been underestimated so that $`\omega _\mathrm{b}\mathrm{}>0.023`$ as required by the CMB data plus simple inflation is allowed, this solution may conflict with other astrophysical constraints. For instance, X-ray observations of clusters of galaxies can be used to determine the ratio of baryons to dark matter , and $`\omega _\mathrm{b}=0.03`$ can only be reconciled with these observations by having $`\mathrm{\Omega }_m\mathrm{}>0.7`$ which would conflict with the supernova 1a results and other estimates of the dark matter density . On the other hand, the tilt solution is no panacea either. In a class of popular inflationary models known as power law inflation, the amplitude of the tensor component is approximately related to the tilt of the scalar spectrum, $`r7(1n_s)`$ . If we choose to fit the data by lowering the tilt to $`n_s=0.9`$, this would raise the COBE-normalization by $`70\%`$. Models that match the COBE normalization therefore make the first peak too low by a factor of 1.7 in power, which is ruled out by the data. In other words, imposing $`r7(1n_s)`$ (which we have not done in our analysis) would exclude $`\omega _\mathrm{b}`$ as low as 0.02. Thus the simple tilt solution does not work for all inflation models. Could the apparent problem be a mere statistical fluke? It would certainly be premature to claim a rock-solid discrepancy between CMB and nucleosynthesis plus power law inflation. The $`\chi ^2`$-value for the best fit inflation model with $`\omega _\mathrm{b}=0.02`$ is still statistically acceptable ($`\chi ^281`$ for 87 degrees of freedom reduced by about 5 effective parameters). However, serious discrepancies in peak heights tend to get statistically diluted by the swarm of points with large error bars at lower $`ł`$ that agree with most anything reasonable (indeed, $`\chi ^2`$ drops down to 71 for $`\omega _\mathrm{b}=0.03`$ and as low as 68 without any priors), and the relative likelihood rises sharply with $`\omega _\mathrm{b}`$ regardless of what priors are imposed. To assess the sensitivity of the results to the choice of data, we therefore repeated our entire analysis for the following cases: (a) using all the data except Maxima and (b) using only COBE and the new Boomerang data. Omitting Maxima removed the “CMB only” and “CMB+$`h`$” exclusion regions that are seen to protrude in from the left in Figure 4. This is because the Maxima points place an upper limit on the height of (the left part of) the third peak, effectively giving an upper limit on the baryon density. Dropping Maxima also loosened the upper limit on $`\omega _{dm}`$ somewhat and marginally weakened the bounds on $`\mathrm{\Omega }_\mathrm{k}`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }`$. The other constraints were essentially unaffected. Most importantly, the lower bound on $`\omega _\mathrm{b}`$ seen in Figure 4 remained unchanged, since it comes from the low ratio of the 2nd to 1st peak heights . Although our inclusion of the 10% uncertainly in the Boomerang’s calibration (20% in power) was not very important in our full analysis, as the fitting procedure de facto calibrated Boomerang off of other experiments, this substantially degraded the results of our COBE + Boomerang analysis. We therefore repeated it three more times, without the calibration error but multiplying the Boomerang points by 0.9, 1.0 and 1.1, respectively. The results were quite similar to those using all the data, as expected from the experimental concordance seen in Figure 1. However, most constraints got slightly tighter, consistent with the above-mentioned $`\chi ^2`$ dilution hypothesis. Rather than go away, the baryon problem became exacerbated: the 95% inflationary lower limit on $`\omega _\mathrm{b}`$ was tightened from $`0.024`$ to 0.027 with $`\chi ^2=12`$ (with a total of 20 Boomerang + COBE points and 4 free parameters). In contrast, the tilt solution gave $`\chi ^2=22`$, and higher still when the Boomerang normalization was raised or lowered by 10%. Most strikingly, in the COBE+Boomerang version of Figure 4, $`\omega _\mathrm{b}`$ is not permitted to be low enough to agree with nucleosynthesis for any value of the tilt $`n_s`$, so the tilt solution may have worked using all the data merely because of the above-mentioned dilution effect. Can the baryon problem be explained by inaccuracies in our numerical method? The correlations between the Boomerang points (which we could not include since the have not yet been made public) are reportedly very small . Although a range of approximations are involved as detailed in , for instance in the likelihood calculation, it appears unlikely that such inaccuracies are large enough to have a major impact on the lower bound on $`\omega _\mathrm{b}`$. Perhaps the best indication of this is that a number of independent analyses have been made available since this paper was originally submitted, using a wide range of computational techniques, and they all favor baryon fractions in excess of the current nucleosynthesis prediction. More baryons also solve some older problems . A number of other ways out have been proposed , ranging from mechanisms for delaying standard recombination to time-variation of physical constants and severe mis-estimates of the Boomerang beam width. However, these explanations are all of a highly speculative nature. An excellent way to clear up this mystery will be to search for a third acoustic peak, which is boosted by more baryons but suppressed by most of the other proposed remedies. Apart from the matter budget, Table 1 and Figure 3 also show that the CMB data provides perhaps the first meaningful upper limit on $`\tau `$, the optical depth due to reionization (compare ). Since $`\tau h\mathrm{\Omega }_\mathrm{b}z_{\mathrm{ion}}^{3/2}\mathrm{\Omega }_\mathrm{m}^{1/2}`$ if the redshift of reionization $`z_{\mathrm{ion}}1`$, our lower limit $`\omega _\mathrm{b}>0.024`$ combined with our upper limit $`\tau <0.35`$ give the constraint $`z_{\mathrm{ion}}\mathrm{}<49h^{2/3}\mathrm{\Omega }_\mathrm{m}^{1/3}`$, or $`z_{\mathrm{ion}}\mathrm{}<28`$ for $`h=\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$. This is compatible with the range $`z_{\mathrm{ion}}=820`$ favored by numerical simulations, but challenges more extreme models. In conclusion, the new Boomerang results look like a triumph for the simplest possible inflationary model but for one rather large fly in the ointment: the lack of a significant second acoustic peak suggests that we may need to abandon either a popular version of inflation, the current nucleosynthesis constraints, or some even more cherished assumption. In answering one question, Boomerang has raised another. Its answer is likely to lie in the third peak, and the race to reach it has now begun. The authors wish to thank John Beacom, Kevin Cahill, Angélica de Oliveira-Costa, Mark Devlin, Andrew Hamilton, David Hogg, Wayne Hu, Lam Hui, William Kinney, Andrew Liddle, Dominik Schwarz, Paul Steinhardt and Ned Wright for helpful comments and discussions. Support for this work was provided by NSF grant AST00-71213, NASA grant NAG5-9194 and Hubble Fellowship HF-01116.01-98A from STScI, operated by AURA, Inc. under NASA contract NAS5-26555.
warning/0004/hep-ph0004074.html	ar5iv	text	# On the Applicability of Weak-Coupling Results in High Density QCD ## I Introduction The starting point for a description of matter at high baryon density and low temperature is a Fermi sea of quarks. The important degrees of freedom — those whose fluctuations cost little free energy — are those involving quarks near the Fermi surface. We know from the work of Bardeen, Cooper, and Schrieffer that any attractive interaction between the quarks, regardless how weak, makes the Fermi sea unstable to the formation of a condensate of Cooper pairs. In QCD, the interaction of two quarks whose colors are antisymmetric (the color $`\overline{\mathrm{𝟑}}_A`$ channel) is attractive. (The attractiveness of this interaction can be seen from single-gluon exchange, as is relevant at short distances, or via counting strings or analyzing the instanton induced coupling, as may be relevant at longer distances.) We therefore expect that under any circumstance in which cold dense quark matter is present, it will be in a color superconducting phase . The one caveat is that this conclusion is known to be false if the number of colors is $`N_c=\mathrm{}`$. Recent work indicates that quark matter is in a color superconducting phase for $`N_c`$ less than of order thousands, and in this paper we only discuss QCD with $`N_c=3`$. We now know much about the symmetries and physical properties of color superconducting quark matter. The dominant condensate in QCD with two flavors of quarks is in the color $`\overline{\mathrm{𝟑}}_A`$ channel, breaking $`SU(3)_{\mathrm{color}}SU(2)`$, and is a flavor singlet. Quarks with two of three colors have a gap in this 2SC phase, and five of eight gluons get a mass via the Meissner effect. In QCD with three flavors of quarks, the Cooper pairs cannot be flavor singlets, and flavor symmetries are necessarily broken. The symmetries of the phase which results have been analyzed in Ref. , and were in fact first analyzed in a different (zero density) context in Ref. . The dominant condensate locks color and flavor symmetries, leaving an unbroken global symmetry under simultaneous $`SU(3)`$ transformations of color, left-flavor, and right-flavor. In this CFL phase, all nine quarks have a gap and all eight gluons have a mass. Chiral symmetry is spontaneously broken, as is baryon number, and there are consequently nine massless Goldstone bosons . Matter in the CFL phase is therefore similar in many respects to superfluid hypernuclear matter . The fact that color superconducting phases always feature either chiral symmetry breaking (as in the CFL phase) or some quarks which remain gapless (as in the 2SC phase) may be understood as a consequence of imposing ’t Hooft’s anomaly matching criterion . The first order phase transition between the CFL and 2SC phases has been analyzed in detail, but all that will concern us below is that any finite strange quark mass is unimportant at large enough $`\mu `$, and quark matter is therefore in the CFL phase at asymptotically large $`\mu `$. Much recent work has resulted in two classes of estimates of the magnitude of $`\mathrm{\Delta }`$, the gap in the density of quasiparticle states in the superconducting phase. The first class of estimates are done within the context of models whose parameters are chosen to give reasonable vacuum physics. Examples include analyses in which the interaction between quarks is replaced simply by four-fermion interactions with the quantum numbers of the instanton interaction or of single-gluon exchange and more sophisticated analyses done using instanton liquid models. Renormalization group analyses have also been used to explore the space of all possible four-fermion interactions allowed by the symmetries of QCD. These methods yield results which are in qualitative agreement: the gaps range from several tens of MeV up to as much as about 100 MeV and the corresponding critical temperatures, above which the superconducting condensates vanish, can be as large as about 50 MeV. The second class of estimates uses $`\mu \mathrm{}`$ physics as a guide. At asymptotically large $`\mu `$, models with short range interactions are bound to fail, because the dominant interaction is due to the long-range magnetic interaction coming from single-gluon exchange. The collinear infrared divergence in small-angle scattering via single-gluon exchange results in a gap which is parametrically larger at $`\mu \mathrm{}`$ than it would be for any point-like four-fermion interaction. Son showed that this collinear divergence is regulated by Landau damping (dynamical screening) and that as a consequence, the parametric dependence of the gap in the limit in which the QCD coupling $`g0`$ is $$\frac{\mathrm{\Delta }}{\mu }\frac{1}{g^5}\mathrm{exp}\left(\frac{3\pi ^2}{\sqrt{2}g}\right),$$ (1) which is more easily seen as an expansion in $`g`$ when rewritten as $$\mathrm{ln}\left(\frac{\mathrm{\Delta }}{\mu }\right)=\frac{3\pi ^2}{\sqrt{2}}\frac{1}{g}5\mathrm{ln}g+f(g).$$ (2) This equation should be viewed as a definition of $`f(g)`$, which will include a term which is constant for $`g0`$ and terms which vanish for $`g0`$. The result (1) has been confirmed using a variety of methods, and several estimates of $`lim_{g0}f(g)`$ exist in the literature. For example, Schaefer and Wilczek find $$\underset{g0}{lim}f(g)\mathrm{ln}\left[2^{1/3}256\pi ^4\left(\frac{2}{3}\right)^{5/2}\right]=8.88$$ (3) in the CFL phase (see also Ref. ), and Brown, Liu, and Ren find a result for $`lim_{g0}f(g)`$ which is smaller by $`(\pi ^2+4)/8\mathrm{ln}2=1.04`$. If this asymptotic expression is applied by taking $`g=g(\mu )`$ from the perturbative QCD $`\beta `$-function (with $`\mathrm{\Lambda }_{\mathrm{QCD}}=200`$ MeV), evaluating $`\mathrm{\Delta }`$ at $`\mu 500`$ MeV yields gaps in rough agreement with the estimates based on zero-density phenomenology. The central purpose of this paper is to demonstrate that this nice agreement must at present be seen as coincidental, because present estimates for $`f`$ are demonstrably uncontrolled for $`g>g_c0.8`$, corresponding to $`\mu <\mu _c`$ with $`\mu _c10^8`$ or higher. The weak-coupling calculations are derived from analyses (done using varying approximations) of the one-loop Schwinger-Dyson equation without vertex correction, and (with one exception) yield gauge dependent results. However, Schaefer and Wilczek argue that the result for $`lim_{g0}f(g)`$ in such a calculation is gauge invariant. The one calculation which is gauge invariant throughout is the calculation of $`T_c`$ (and hence $`\mathrm{\Delta }`$ since the BCS relation $`T_c=0.57\mathrm{\Delta }`$ holds ) done by Brown, Liu, and Ren. As in other calculations, however, these authors neglect vertex corrections. Our purpose is to use the fact that our calculation (like most) is gauge dependent, and only gauge invariant for $`g0`$, to estimate the $`g`$ above which vertex corrections, left out of all calculations, cannot be neglected. We begin by sketching the derivation of the one-loop Schwinger-Dyson equation for $`\mathrm{\Delta }`$, making as few approximations as we can. We solve the resulting gap equation numerically in several different gauges. Our results are (yet one more) confirmation of (1). Furthermore, we do find evidence that the gauge dependence of $`f`$ decreases for $`g0`$. However, this decrease only begins to set in for $`g0.8`$. This implies that the contributions to $`\mathrm{\Delta }`$ which have been neglected — like those arising from vertex corrections — only become subleading for $`gg_c0.8`$. If we translate $`g_c`$ to $`\mu _c`$ by assuming $`g`$ should be taken as $`g(\mu )`$, this corresponds to $`\mu _c10^8`$ MeV. Recent work shows that $`g`$ should be evaluated at a much lower ($`g`$-dependent) scale than $`\mu `$. This means that the condition $`g<g_c0.8`$ would translate into $`\mu >\mu _c`$ with $`\mu _c`$ orders of magnitude larger than $`\mu _c10^8`$ MeV. The original purpose of our investigation was to do a self-consistent calculation of the influence of the Meissner effect on the magnitude of the gap in the CFL phase. In the presence of a condensate, the gluon propagator is modified: some gluons get a mass. In the CFL phase, all gluons get a mass, and this makes a calculation based on perturbative single-gluon exchange a self-consistent and complete description of the physics at asymptotically large $`\mu `$, with no remaining infrared problems. (In the 2SC phase, in contrast, the calculation of $`\mathrm{\Delta }`$ leaves unanswered any questions about the non-Abelian infrared physics of the three gluons left unscreened by the condensate.) We felt that this motivation warranted a self-consistent calculation in which we calculate the gap using a Schwinger-Dyson equation in which the gluon propagator is modified not only by the presence of the Fermi sea (Debye mass, Landau damping) but is also affected by the condensate (the Meissner effect). We set this calculation up in an appendix. Previous work, beginning with that of Ref. , shows that the form of Eq. (1) is unmodified by including the Meissner effect, but $`f(g)`$ is affected. Our preliminary results suggest that the changes in $`f(g)`$ are small, as anticipated in Refs. . Indeed, the effects of physics left out of the present analysis, which we have diagnosed via the gauge dependence of $`f(g)`$, are much larger than those introduced by the Meissner effect at any $`g`$ we have investigated. ## II Deriving the gap equation In this section, we derive the gap equation for QCD with three massless flavors which is valid at asymptotically high densities. We follow Ref. , but make fewer approximations. Because our point is to stress the importance of effects which we do not calculate, we will make our assumptions and approximations very clear as we proceed. In other words, since the lesson we learn from our results is that they cannot yet be trusted, it is important to detail carefully all points at which we leave something out. We use the standard Nambu-Gorkov formalism by defining an eight-component field $`\mathrm{\Psi }=(\psi ,\overline{\psi }^T)`$. In this basis, the inverse quark propagator takes the form $$S^1(k)=\left(\begin{array}{cc}k/+\mu \gamma _0& \overline{\mathrm{\Delta }}\\ \mathrm{\Delta }& (k/\mu \gamma _0)^T\end{array}\right)$$ (4) where $`\overline{\mathrm{\Delta }}=\gamma _0\mathrm{\Delta }^{}\gamma _0`$. The color, flavor, and Dirac indices are suppressed in the above expression. The diagonal blocks correspond to ordinary propagation and the off-diagonal blocks reflect the possibility for “anomalous propagation” in the presence of a diquark condensate. We make the following ansatz for the form of the gap matrix: $`\mathrm{\Delta }_{ij}^{ab}(k)`$ $`=(\lambda _I^A)^{ab}(\lambda _I^A)_{ij}C\gamma _5\left(\mathrm{\Delta }_1^A(k_0)P_+(k)+\mathrm{\Delta }_2^A(k_0)P_{}(k)\right)`$ (5) $`+(\lambda _J^S)^{ab}(\lambda _J^S)_{ij}C\gamma _5\left(\mathrm{\Delta }_1^S(k_0)P_+(k)+\mathrm{\Delta }_2^S(k_0)P_{}(k)\right)`$ (6) Here, $`a,b=1,2,3`$ are color indices, $`i,j=1,2,3`$ are flavor indices, $`\lambda _I^A`$ are antisymmetric $`U(3)`$ color or flavor matrices with $`I=1,2,3`$, and $`\lambda _J^S`$ are symmetric $`U(3)`$ color or flavor with $`J=1,\mathrm{},6`$, and the projection operators $`P_\pm `$ are defined as $`P_+(k)={\displaystyle \frac{1+\stackrel{}{\alpha }\widehat{k}}{2}}`$ (7) $`P_{}(k)={\displaystyle \frac{1\stackrel{}{\alpha }\widehat{k}}{2}}`$ (8) with $`\stackrel{}{\alpha }=\gamma _0\stackrel{}{\gamma }`$. By making this ansatz, we are making several assumptions: * First, we have taken $`\mathrm{\Delta }_1^A`$, $`\mathrm{\Delta }_2^A`$, $`\mathrm{\Delta }_1^S`$, and $`\mathrm{\Delta }_2^S`$ to be functions of $`k_0`$ only. All are in principle functions of both $`k_0`$ and $`\|\stackrel{}{k}\|`$, but we assume that they are dominated by $`\|\stackrel{}{k}\|\mu `$. This is a standard assumption, and although we do not expect that relaxing this assumption would resolve the problems which we diagnose below, this does belong on the list of potential cures. * Second, we have explicitly separated the gaps which are antisymmetric $`\overline{\mathrm{𝟑}}_A`$ in color and flavor from those which are symmetric $`\mathrm{𝟔}_S`$ in color and flavor and, in both cases, we have assumed that the favored channel is the one in which color and flavor rotations are locked. The color and flavor structure of our ansatz is thus precisely that first explored in Ref. , which allows quarks of all three colors and all three flavors to pair. Subsequent work confirms that this is the favored condensate, and we will not attempt to further generalize it here. * Third, we have assumed that the Cooper pairs in the condensate have zero spin and orbital angular momentum. This seems a safe assumption in the CFL phase, where the dominant condensate, made of Cooper pairs with zero spin and orbital angular momentum, leaves no quarks ungapped. * Fourth, we neglect $`\overline{\psi }\psi `$ condensates. Since chiral symmetry is broken in the CFL phase, these must be nonzero . This applies to both color singlet and color octet $`\overline{\psi }\psi `$ condensates. Such condensates are small, however, and we expect that neglecting them results in only a very small error in the magnitude of the dominant diquark condensate. The most important assumption we make is that we obtain the gap by solving the one-loop Schwinger-Dyson equation of the form $$S^1(k)S_0^1(k)=ig^2\frac{d^4q}{(2\pi )^4}\mathrm{\Gamma }_\mu ^aS(q)\mathrm{\Gamma }_\nu ^bD_{ab}^{\mu \nu }(kq),$$ (9) using a medium-modified gluon propagator described below and unmodified vertices $$\mathrm{\Gamma }_\mu ^a=\left(\begin{array}{cc}\gamma _\mu \lambda ^a/2& 0\\ 0& (\gamma _\mu \lambda ^a/2)^T\end{array}\right).$$ (10) Here, $`S_0`$ is the bare fermion propagator with $`\mathrm{\Delta }=0`$. Note that we use a Minkowski metric unless stated otherwise. We will demonstrate that our results are completely uncontrolled for $`g>g_c0.8`$. This breakdown could in principle reflect a failure of any of our assumptions. We expect, however, that it arises because contributions which have been truncated in writing (9) are large for $`g>g_c`$. That is, we expect that this truncation (and not any of the simplifications introduced by our ansatz for $`\mathrm{\Delta }`$) is the most significant assumption we are making. We obtain four coupled gap equations $`\mathrm{\Delta }_{1,2}^A(k_0)={\displaystyle \frac{i}{6}}g^2{\displaystyle \frac{d^4q}{(2\pi )^4}\mathrm{Tr}\left[P_\pm (k)\gamma _\mu \left(P_+(q)a_+(q)+P_{}(q)a_{}(q)\right)\gamma _\nu \right]D^{\mu \nu }(kq)}`$ (11) $`\mathrm{\Delta }_{1,2}^S(k_0)={\displaystyle \frac{i}{6}}g^2{\displaystyle \frac{d^4q}{(2\pi )^4}\mathrm{Tr}\left[P_\pm (k)\gamma _\mu \left(P_+(q)b_+(q)+P_{}(q)b_{}(q)\right)\gamma _\nu \right]D^{\mu \nu }(kq)}`$ (12) where $`P_\pm `$ means $`P_+`$ in the $`\mathrm{\Delta }_1`$ equation and $`P_{}`$ in the $`\mathrm{\Delta }_2`$ equation and where $`a_+(q)`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Delta }_2^S(q_0)\mathrm{\Delta }_2^A(q_0)}{q_0^2(\|\stackrel{}{q}\|+\mu )^24\left[\mathrm{\Delta }_2^A(q_0)+2\mathrm{\Delta }_2^S(q_0)\right]^2}}`$ (13) $`+{\displaystyle \frac{\left[\mathrm{\Delta }_2^A(q_0)\mathrm{\Delta }_2^S(q_0)\right]\left[q_0^2+(\|\stackrel{}{q}\|+\mu )^2+(5\mathrm{\Delta }_2^A(q_0)+7\mathrm{\Delta }_2^S(q_0))(\mathrm{\Delta }_2^A(q_0)+3\mathrm{\Delta }_2^S(q_0))\right]}{\left[q_0^2(\|\stackrel{}{q}\|+\mu )^2(\mathrm{\Delta }_2^A(q_0)\mathrm{\Delta }_2^S(q_0))^2\right]\left[q_0^2(\|\stackrel{}{q}\|+\mu )^24(\mathrm{\Delta }_2^A(q_0)+2\mathrm{\Delta }_2^S(q_0))^2\right]}}`$ (14) $`a_{}(q)`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Delta }_1^S(q_0)\mathrm{\Delta }_1^A(q_0)}{q_0^2(\|\stackrel{}{q}\|\mu )^24\left[\mathrm{\Delta }_1^A(q_0)+2\mathrm{\Delta }_1^S(q_0)\right]^2}}`$ (15) $`+{\displaystyle \frac{\left[\mathrm{\Delta }_1^A(q_0)\mathrm{\Delta }_1^S(q_0)\right]\left[q_0^2+(\|\stackrel{}{q}\|\mu )^2+(5\mathrm{\Delta }_1^A(q_0)+7\mathrm{\Delta }_1^S(q_0))(\mathrm{\Delta }_1^A(q_0)+3\mathrm{\Delta }_1^S(q_0))\right]}{\left[q_0^2(\|\stackrel{}{q}\|\mu )^2(\mathrm{\Delta }_1^A(q_0)\mathrm{\Delta }_1^S(q_0))^2\right]\left[q_0^2(\|\stackrel{}{q}\|\mu )^24(\mathrm{\Delta }_1^A(q_0)+2\mathrm{\Delta }_1^S(q_0))^2\right]}}`$ (16) $`b_+(q)`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Delta }_2^S(q_0)}{q_0^2(\|\stackrel{}{q}\|+\mu )^24\left[\mathrm{\Delta }_2^A(q_0)+2\mathrm{\Delta }_2^S(q_0)\right]^2}}`$ (17) $`+{\displaystyle \frac{\left[\mathrm{\Delta }_2^A(q_0)\mathrm{\Delta }_2^S(q_0)\right]\left[\mathrm{\Delta }_2^A(q_0)+\mathrm{\Delta }_2^S(q_0)\right]\left[\mathrm{\Delta }_2^A(q_0)+5\mathrm{\Delta }_2^S(q_0)\right]}{\left[q_0^2(\|\stackrel{}{q}\|+\mu )^2(\mathrm{\Delta }_2^A(q_0)\mathrm{\Delta }_2^S(q_0))^2\right]\left[q_0^2(\|\stackrel{}{q}\|+\mu )^24(\mathrm{\Delta }_2^A(q_0)+2\mathrm{\Delta }_2^S(q_0))^2\right]}}`$ (18) $`b_{}(q)`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Delta }_1^S(q_0)}{q_0^2(\|\stackrel{}{q}\|\mu )^24\left[\mathrm{\Delta }_1^A(q_0)+2\mathrm{\Delta }_1^S(q_0)\right]^2}}`$ (19) $`+{\displaystyle \frac{\left[\mathrm{\Delta }_1^A(q_0)\mathrm{\Delta }_1^S(q_0)\right]\left[\mathrm{\Delta }_1^A(q_0)+\mathrm{\Delta }_1^S(q_0)\right]\left[\mathrm{\Delta }_1^A(q_0)+5\mathrm{\Delta }_1^S(q_0)\right]}{\left[q_0^2(\|\stackrel{}{q}\|\mu )^2(\mathrm{\Delta }_1^A(q_0)\mathrm{\Delta }_1^S(q_0))^2\right]\left[q_0^2(\|\stackrel{}{q}\|\mu )^24(\mathrm{\Delta }_1^A(q_0)+2\mathrm{\Delta }_1^S(q_0))^2\right]}}.`$ (20) In a general covariant gauge, the resummed gluon propagator is given by $$D_{\mu \nu }(q)=\frac{P_{\mu \nu }^T}{q^2G(q)}+\frac{P_{\mu \nu }^L}{q^2F(q)}\xi \frac{q_\mu q_\nu }{q^4}$$ (21) where $`G(q)`$ and $`F(q)`$ are functions of $`q_0`$ and $`\|\stackrel{}{q}\|`$ and the projectors $`P_{\mu \nu }^{T,L}`$ are defined as follows: $$P_{ij}^T=\delta _{ij}\widehat{q_i}\widehat{q_j},P_{00}^T=P_{0i}^T=0,P_{\mu \nu }^L=g_{\mu \nu }+\frac{q_\mu q_\nu }{q^2}P_{\mu \nu }^T.$$ (22) The functions $`F`$ and $`G`$ describe the effects of the medium on the gluon propagator. If we neglect the Meissner effect (that is, if we neglect the modification of $`F(q)`$ and $`G(q)`$ due to the gap $`\mathrm{\Delta }`$ in the fermion propagator) then $`F(q)`$ describes Thomas-Fermi screening and $`G(q)`$ describes Landau damping and they are given in the HDL approximation by $`F(q)`$ $`=`$ $`m^2{\displaystyle \frac{q^2}{\|\stackrel{}{q}\|}}\left(1{\displaystyle \frac{iq_0}{\|\stackrel{}{q}\|}}Q_0\left({\displaystyle \frac{iq_0}{\|\stackrel{}{q}\|}}\right)\right),\mathrm{with}Q_0(x)={\displaystyle \frac{1}{2}}\mathrm{log}\left({\displaystyle \frac{x+1}{x1}}\right)`$ (23) $`G(q)`$ $`=`$ $`{\displaystyle \frac{1}{2}}m^2{\displaystyle \frac{iq_0}{\|\stackrel{}{q}\|}}\left[\left(1\left({\displaystyle \frac{iq_0}{\|\stackrel{}{q}\|}}\right)^2\right)Q_0\left({\displaystyle \frac{iq_0}{\|\stackrel{}{q}\|}}\right)+{\displaystyle \frac{iq_0}{\|\stackrel{}{q}\|}}\right],`$ (24) where $`m^2=3g^2\mu ^2/2\pi ^2`$ is the Debye screening mass for $`N_f=3`$. We discuss the modifications of $`F(q)`$ and $`G(q)`$ due to the Meissner effect in an Appendix. In order to obtain the final form of the gap equation, we need the following trace: $$\begin{array}{c}\mathrm{Tr}\left[P_\pm (k)\gamma _\mu \left(P_+(q)a_+(q)+P_{}(q)a_{}(q)\right)\gamma _\nu \right]D^{\mu \nu }(kq)\hfill \\ \\ \begin{array}{ccc}& =a_+(q)\hfill & [2\frac{1\widehat{(kq)}\widehat{k}\widehat{(kq)}\widehat{q}}{(kq)^2G(kq)}+\frac{1\widehat{k}\widehat{q}\frac{(kq)_0^2+(\stackrel{}{k}\stackrel{}{q})^2}{(kq)^2}\pm 2\widehat{(kq)}\widehat{k}\widehat{(kq)}\widehat{q}\frac{(kq)_0^2}{(kq)^2}}{(kq)^2F(kq)}\hfill \\ & & +\frac{\xi }{(kq)^2}(1\widehat{k}\widehat{q}\frac{(kq)_0^2+(\stackrel{}{k}\stackrel{}{q})^2}{(kq)^2}\pm 2\widehat{(kq)}\widehat{k}\widehat{(kq)}\widehat{q}\frac{(\stackrel{}{k}\stackrel{}{q})^2}{(kq)^2})]\hfill \\ & & \\ & +a_{}(q)\hfill & [2\frac{1\pm \widehat{(kq)}\widehat{k}\widehat{(kq)}\widehat{q}}{(kq)^2G(kq)}+\frac{1\pm \widehat{k}\widehat{q}\frac{(kq)_0^2+(\stackrel{}{k}\stackrel{}{q})^2}{(kq)^2}2\widehat{(kq)}\widehat{k}\widehat{(kq)}\widehat{q}\frac{(kq)_0^2}{(kq)^2}}{(kq)^2F(kq)}\hfill \\ & & +\frac{\xi }{(kq)^2}(1\pm \widehat{k}\widehat{q}\frac{(kq)_0^2+(\stackrel{}{k}\stackrel{}{q})^2}{(kq)^2}2\widehat{(kq)}\widehat{k}\widehat{(kq)}\widehat{q}\frac{(\stackrel{}{k}\stackrel{}{q})^2}{(kq)^2})].\hfill \end{array}\hfill \end{array}$$ (25) This allows us to recast Eq. (11) into the following form: Δ1A(k0)=i6g2 d4q(2π)4[a+(q)(21(kq)^k^(kq)^q^(kq)2G(kq)+1k^q^(kq)02+(kq)2(kq)2+2(kq)^k^(kq)^q^(kq)02(kq)2(kq)2F(kq)+ξ(kq)2(1k^q^(kq)02+(kq)2(kq)2+2(kq)^k^(kq)^q^(kq)2(kq)2))+a(q)(21+(kq)^k^(kq)^q^(kq)2G(kq)+1+k^q^(kq)02+(kq)2(kq)22(kq)^k^(kq)^q^(kq)02(kq)2(kq)2F(kq)+ξ(kq)2(1+k^q^(kq)02+(kq)2(kq)22(kq)^k^(kq)^q^(kq)2(kq)2))]Δ2A(k0)=i6g2 d4q(2π)4[a+(q)(21+(kq)^k^(kq)^q^(kq)2G(kq)+1+k^q^(kq)02+(kq)2(kq)22(kq)^k^(kq)^q^(kq)02(kq)2(kq)2F(kq)+ξ(kq)2(1+k^q^(kq)02+(kq)2(kq)22(kq)^k^(kq)^q^(kq)2(kq)2))+a(q)(21(kq)^k^(kq)^q^(kq)2G(kq)+1k^q^(kq)02+(kq)2(kq)2+2(kq)^k^(kq)^q^(kq)02(kq)2(kq)2F(kq)+ξ(kq)2(1k^q^(kq)02+(kq)2(kq)2+2(kq)^k^(kq)^q^(kq)2(kq)2))]\begin{array}[]{ll}\Delta_{1}^{A}(k_{0})=-{i\over 6}g^{2}\begin{minipage}[l][0.5in][c]{7.22743pt}\par\vskip 0.0pt\vskip 8.5pt plus 3.0pt minus 4.0pt\hbox to469.75499pt{\hskip 0.0pt\hskip 0.0pt plus 1000.0pt$\displaystyle\int$\hskip 0.0pt plus 1000.0pt}\vskip 8.5pt plus 3.0pt minus 4.0pt\noindent\ignorespaces\end{minipage}\frac{d^{4}q}{(2\pi)^{4}}&\!\left[a_{+}(q)\!\left(\!2{-1-\widehat{(k-q)}\cdot\hat{k}\widehat{(k-q)}\cdot\hat{q}\over(k-q)^{2}-G(k-q)}\!+\!{-1-\hat{k}\cdot\hat{q}{(k-q)_{0}^{2}+(\vec{k}-\vec{q})^{2}\over(k-q)^{2}}+2\widehat{(k-q)}\cdot\hat{k}\widehat{(k-q)}\cdot\hat{q}{(k-q)_{0}^{2}\over(k-q)^{2}}\over(k-q)^{2}-F(k-q)}\right.\right.\\ &+\!\left.{\xi\over(k-q)^{2}}\!\left(\!1\!-\!\hat{k}\cdot\hat{q}{(k-q)_{0}^{2}+(\vec{k}-\vec{q})^{2}\over(k-q)^{2}}\!+\!2\widehat{(k-q)}\cdot\hat{k}\widehat{(k-q)}\cdot\hat{q}{(\vec{k}-\vec{q})^{2}\over(k-q)^{2}}\right)\right)\\ &+a_{-}(q)\!\!\left(\!2{-1+\widehat{(k-q)}\cdot\hat{k}\widehat{(k-q)}\cdot\hat{q}\over(k-q)^{2}-G(k-q)}\!+\!{-1+\hat{k}\cdot\hat{q}{(k-q)_{0}^{2}+(\vec{k}-\vec{q})^{2}\over(k-q)^{2}}-2\widehat{(k-q)}\cdot\hat{k}\widehat{(k-q)}\cdot\hat{q}{(k-q)_{0}^{2}\over(k-q)^{2}}\over(k-q)^{2}-F(k-q)}\right.\\ &+\!\left.\left.{\xi\over(k-q)^{2}}\!\left(\!1\!+\!\hat{k}\cdot\hat{q}{(k-q)_{0}^{2}+(\vec{k}-\vec{q})^{2}\over(k-q)^{2}}\!-\!2\widehat{(k-q)}\cdot\hat{k}\widehat{(k-q)}\cdot\hat{q}{(\vec{k}-\vec{q})^{2}\over(k-q)^{2}}\right)\right)\right]\\ &\\ \Delta_{2}^{A}(k_{0})=-{i\over 6}g^{2}\begin{minipage}[l][0.5in][c]{7.22743pt}\par\vskip 0.0pt\vskip 8.5pt plus 3.0pt minus 4.0pt\hbox to469.75499pt{\hskip 0.0pt\hskip 0.0pt plus 1000.0pt$\displaystyle\int$\hskip 0.0pt plus 1000.0pt}\vskip 8.5pt plus 3.0pt minus 4.0pt\noindent\ignorespaces\end{minipage}\frac{d^{4}q}{(2\pi)^{4}}&\!\left[a_{+}(q)\!\left(\!2{-1+\widehat{(k-q)}\cdot\hat{k}\widehat{(k-q)}\cdot\hat{q}\over(k-q)^{2}-G(k-q)}\!+\!{-1+\hat{k}\cdot\hat{q}{(k-q)_{0}^{2}+(\vec{k}-\vec{q})^{2}\over(k-q)^{2}}-2\widehat{(k-q)}\cdot\hat{k}\widehat{(k-q)}\cdot\hat{q}{(k-q)_{0}^{2}\over(k-q)^{2}}\over(k-q)^{2}-F(k-q)}\right.\right.\\ &+\!\left.{\xi\over(k-q)^{2}}\!\left(\!1\!+\!\hat{k}\cdot\hat{q}{(k-q)_{0}^{2}+(\vec{k}-\vec{q})^{2}\over(k-q)^{2}}\!-\!2\widehat{(k-q)}\cdot\hat{k}\widehat{(k-q)}\cdot\hat{q}{(\vec{k}-\vec{q})^{2}\over(k-q)^{2}}\right)\right)\\ &+a_{-}(q)\!\!\left(2{-1-\widehat{(k-q)}\cdot\hat{k}\widehat{(k-q)}\cdot\hat{q}\over(k-q)^{2}-G(k-q)}\!+\!{-1-\hat{k}\cdot\hat{q}{(k-q)_{0}^{2}+(\vec{k}-\vec{q})^{2}\over(k-q)^{2}}+2\widehat{(k-q)}\cdot\hat{k}\widehat{(k-q)}\cdot\hat{q}{(k-q)_{0}^{2}\over(k-q)^{2}}\over(k-q)^{2}-F(k-q)}\right.\\ &+\!\left.\left.{\xi\over(k-q)^{2}}\!\left(\!1\!-\!\hat{k}\cdot\hat{q}{(k-q)_{0}^{2}+(\vec{k}-\vec{q})^{2}\over(k-q)^{2}}\!+\!2\widehat{(k-q)}\cdot\hat{k}\widehat{(k-q)}\cdot\hat{q}{(\vec{k}-\vec{q})^{2}\over(k-q)^{2}}\right)\right)\right]\end{array} (26) Δ1S(k0)=i6g2 d4q(2π)4[b+(q)(21(kq)^k^(kq)^q^(kq)2G(kq)+1k^q^(kq)02+(kq)2(kq)2+2(kq)^k^(kq)^q^(kq)02(kq)2(kq)2F(kq)+ξ(kq)2(1k^q^(kq)02+(kq)2(kq)2+2(kq)^k^(kq)^q^(kq)2(kq)2))+b(q)(21+(kq)^k^(kq)^q^(kq)2G(kq)+1+k^q^(kq)02+(kq)2(kq)22(kq)^k^(kq)^q^(kq)02(kq)2(kq)2F(kq)+ξ(kq)2(1+k^q^(kq)02+(kq)2(kq)22(kq)^k^(kq)^q^(kq)2(kq)2))]Δ2S(k0)=i6g2 d4q(2π)4[b+(q)(21+(kq)^k^(kq)^q^(kq)2G(kq)+1+k^q^(kq)02+(kq)2(kq)22(kq)^k^(kq)^q^(kq)02(kq)2(kq)2F(kq)+ξ(kq)2(1+k^q^(kq)02+(kq)2(kq)22(kq)^k^(kq)^q^(kq)2(kq)2))+b(q)(21(kq)^k^(kq)^q^(kq)2G(kq)+1k^q^(kq)02+(kq)2(kq)2+2(kq)^k^(kq)^q^(kq)02(kq)2(kq)2F(kq)+ξ(kq)2(1k^q^(kq)02+(kq)2(kq)2+2(kq)^k^(kq)^q^(kq)2(kq)2))].\begin{array}[]{ll}\Delta_{1}^{S}(k_{0})=-{i\over 6}g^{2}\begin{minipage}[l][0.5in][c]{7.22743pt}\par\vskip 0.0pt\vskip 8.5pt plus 3.0pt minus 4.0pt\hbox to469.75499pt{\hskip 0.0pt\hskip 0.0pt plus 1000.0pt$\displaystyle\int$\hskip 0.0pt plus 1000.0pt}\vskip 8.5pt plus 3.0pt minus 4.0pt\noindent\ignorespaces\end{minipage}\frac{d^{4}q}{(2\pi)^{4}}&\!\left[b_{+}(q)\!\left(\!2{-1-\widehat{(k-q)}\cdot\hat{k}\widehat{(k-q)}\cdot\hat{q}\over(k-q)^{2}-G(k-q)}\!+\!{-1-\hat{k}\cdot\hat{q}{(k-q)_{0}^{2}+(\vec{k}-\vec{q})^{2}\over(k-q)^{2}}+2\widehat{(k-q)}\cdot\hat{k}\widehat{(k-q)}\cdot\hat{q}{(k-q)_{0}^{2}\over(k-q)^{2}}\over(k-q)^{2}-F(k-q)}\right.\right.\\ &+\!\left.{\xi\over(k-q)^{2}}\!\left(\!1\!-\!\hat{k}\cdot\hat{q}{(k-q)_{0}^{2}+(\vec{k}-\vec{q})^{2}\over(k-q)^{2}}\!+\!2\widehat{(k-q)}\cdot\hat{k}\widehat{(k-q)}\cdot\hat{q}{(\vec{k}-\vec{q})^{2}\over(k-q)^{2}}\right)\right)\\ &+b_{-}(q)\!\!\left(\!2{-1+\widehat{(k-q)}\cdot\hat{k}\widehat{(k-q)}\cdot\hat{q}\over(k-q)^{2}-G(k-q)}\!+\!{-1+\hat{k}\cdot\hat{q}{(k-q)_{0}^{2}+(\vec{k}-\vec{q})^{2}\over(k-q)^{2}}-2\widehat{(k-q)}\cdot\hat{k}\widehat{(k-q)}\cdot\hat{q}{(k-q)_{0}^{2}\over(k-q)^{2}}\over(k-q)^{2}-F(k-q)}\right.\\ &+\!\left.\left.{\xi\over(k-q)^{2}}\!\left(\!1\!+\!\hat{k}\cdot\hat{q}{(k-q)_{0}^{2}+(\vec{k}-\vec{q})^{2}\over(k-q)^{2}}\!-\!2\widehat{(k-q)}\cdot\hat{k}\widehat{(k-q)}\cdot\hat{q}{(\vec{k}-\vec{q})^{2}\over(k-q)^{2}}\right)\right)\right]\\ &\\ \Delta_{2}^{S}(k_{0})=-{i\over 6}g^{2}\begin{minipage}[l][0.5in][c]{7.22743pt}\par\vskip 0.0pt\vskip 8.5pt plus 3.0pt minus 4.0pt\hbox to469.75499pt{\hskip 0.0pt\hskip 0.0pt plus 1000.0pt$\displaystyle\int$\hskip 0.0pt plus 1000.0pt}\vskip 8.5pt plus 3.0pt minus 4.0pt\noindent\ignorespaces\end{minipage}\frac{d^{4}q}{(2\pi)^{4}}&\!\left[b_{+}(q)\!\left(\!2{-1+\widehat{(k-q)}\cdot\hat{k}\widehat{(k-q)}\cdot\hat{q}\over(k-q)^{2}-G(k-q)}\!+\!{-1+\hat{k}\cdot\hat{q}{(k-q)_{0}^{2}+(\vec{k}-\vec{q})^{2}\over(k-q)^{2}}-2\widehat{(k-q)}\cdot\hat{k}\widehat{(k-q)}\cdot\hat{q}{(k-q)_{0}^{2}\over(k-q)^{2}}\over(k-q)^{2}-F(k-q)}\right.\right.\\ &+\!\left.{\xi\over(k-q)^{2}}\!\left(\!1\!+\!\hat{k}\cdot\hat{q}{(k-q)_{0}^{2}+(\vec{k}-\vec{q})^{2}\over(k-q)^{2}}\!-\!2\widehat{(k-q)}\cdot\hat{k}\widehat{(k-q)}\cdot\hat{q}{(\vec{k}-\vec{q})^{2}\over(k-q)^{2}}\right)\right)\\ &+b_{-}(q)\!\!\left(2{-1-\widehat{(k-q)}\cdot\hat{k}\widehat{(k-q)}\cdot\hat{q}\over(k-q)^{2}-G(k-q)}\!+\!{-1-\hat{k}\cdot\hat{q}{(k-q)_{0}^{2}+(\vec{k}-\vec{q})^{2}\over(k-q)^{2}}+2\widehat{(k-q)}\cdot\hat{k}\widehat{(k-q)}\cdot\hat{q}{(k-q)_{0}^{2}\over(k-q)^{2}}\over(k-q)^{2}-F(k-q)}\right.\\ &+\!\left.\left.{\xi\over(k-q)^{2}}\!\left(\!1\!-\!\hat{k}\cdot\hat{q}{(k-q)_{0}^{2}+(\vec{k}-\vec{q})^{2}\over(k-q)^{2}}\!+\!2\widehat{(k-q)}\cdot\hat{k}\widehat{(k-q)}\cdot\hat{q}{(\vec{k}-\vec{q})^{2}\over(k-q)^{2}}\right)\right)\right].\end{array} (26) ## III Solving the gap equation In order to obtain a tractable numerical problem, we make two further simplifying assumptions: * First, at weak coupling we expect the physics to be dominated by particles and holes near the Fermi surface. This manifests itself in Eq. (26) in the fact that $`a_{}`$ and $`b_{}`$ have singularities on the Fermi surface while $`a_+`$ and $`b_+`$ are regular there, and we therefore expect that at weak coupling we can neglect $`a_+`$ and $`b_+`$. Upon doing this, we have equations for $`\mathrm{\Delta }_1^{A,S}`$ which do not involve $`\mathrm{\Delta }_2^{A,S}`$. We are only interested in $`\mathrm{\Delta }_1^{A,S}`$, since $`\mathrm{\Delta }_2^{A,S}`$ describe the propagation of antiparticles far from the Fermi surface. If we assume that we are at weak enough coupling that $`a_+`$ and $`b_+`$ can be neglected (that is if we assume that $`\mathrm{\Delta }_1^{A,S}\mu `$) then we can ignore $`\mathrm{\Delta }_2^{A,S}`$ in our calculation of $`\mathrm{\Delta }_1^{A,S}`$. (Note that we are not assuming that $`\mathrm{\Delta }_2^{A,S}`$ is any smaller than $`\mathrm{\Delta }_1^{A,S}`$; there is no reason for this to be true.) We will see that our results break down for $`g0.8`$, at which $`\mathrm{\Delta }<10^7\mu `$. Because $`\mathrm{\Delta }\mu `$, neglecting the effects of $`\mathrm{\Delta }_2^{A,S}`$ on $`\mathrm{\Delta }_1^{A,S}`$ should be a good approximation, and we do not expect that including these effects would cure the problems we discover. This should, however, be investigated further. * Second, we set $`\mathrm{\Delta }_1^S=0`$, and solve an equation for $`\mathrm{\Delta }_1^A`$ alone. This assumption is in fact inconsistent, as the gap in the symmetric channel must be nonzero. This is clear from explicit examination of the gap equations Eq. (26) (and indeed of the gap equations of Ref. ). In fact, this result is manifest on symmetry grounds : in the presence of $`\mathrm{\Delta }_1^A0`$, a nonzero $`\mathrm{\Delta }_1^S`$ breaks no new global symmetries and there is therefore no symmetry to keep it zero. Because single-gluon exchange is repulsive in the symmetric channel, this condensate can only exist in the presence of condensation in the antisymmetric channel. Explicit calculation shows that the symmetric condensates are much smaller than those in the antisymmetric channels. We are therefore confident that keeping $`\mathrm{\Delta }_1^S`$ would yield only a very small correction to $`\mathrm{\Delta }_1^A`$. We must now solve a single gap equation for $`\mathrm{\Delta }_1^A(k_0)`$, which henceforth we denote simply as $`\mathrm{\Delta }(k_0)`$. The reader will see below that this equation is still rather involved. Most authors have made further approximations, valid for $`g0`$. Because we make no further approximations, our results cannot be gauge invariant. This allows us to test the claim that the results become gauge invariant in the limit $`g0`$, and to use the rapidity of the disappearance of gauge dependence as this limit is approached to evaluate at what $`g`$ the contributions we have truncated can legitimately be ignored. In order to obtain numerical solutions, it is convenient to do a Wick rotation $`q_0iq_0`$ to Euclidean space, yielding the gap equation $`\mathrm{\Delta }(k_0)`$ $`=`$ $`{\displaystyle \frac{g^2}{6}}{\displaystyle }{\displaystyle \frac{d^4q}{(2\pi )^4}}[{\displaystyle \frac{\mathrm{\Delta }(q_0)}{q_0^2+(\|\stackrel{}{q}\|\mu )^2+4\mathrm{\Delta }^2(q_0)}}`$ (27) $`+{\displaystyle \frac{\mathrm{\Delta }(q_0)\left(q_0^2+(\|\stackrel{}{q}\|\mu )^2+5\mathrm{\Delta }^2(q_0)\right)}{\left(q_0^2+(\|\stackrel{}{q}\|\mu )^2+\mathrm{\Delta }^2(q_0)\right)\left(q_0^2+(\|\stackrel{}{q}\|\mu )^2+4\mathrm{\Delta }^2(q_0)\right)}}]`$ (28) $`[2{\displaystyle \frac{1\widehat{(kq)}\widehat{k}\widehat{(kq)}\widehat{q}}{(kq)_0^2+(\stackrel{}{k}\stackrel{}{q})^2+G(k_0q_0,\|\stackrel{}{k}\stackrel{}{q}\|)}}`$ (29) $`+{\displaystyle \frac{1+\widehat{k}\widehat{q}\frac{(kq)_0^2+(\stackrel{}{k}\stackrel{}{q})^2}{(kq)_0^2+(\stackrel{}{k}\stackrel{}{q})^2}+2\widehat{(kq)}\widehat{k}\widehat{(kq)}\widehat{q}\frac{(kq)_0^2}{(kq)_0^2+(\stackrel{}{k}\stackrel{}{q})^2}}{(kq)_0^2+(\stackrel{}{k}\stackrel{}{q})^2+F(k_0q_0,\|\stackrel{}{k}\stackrel{}{q}\|)}}`$ (30) $`+\xi {\displaystyle \frac{1+\widehat{k}\widehat{q}\frac{(kq)_0^2+(\stackrel{}{k}\stackrel{}{q})^2}{(kq)_0^2+(\stackrel{}{k}\stackrel{}{q})^2}2\widehat{(kq)}\widehat{k}\widehat{(kq)}\widehat{q}\frac{(\stackrel{}{k}\stackrel{}{q})^2}{(kq)_0^2+(\stackrel{}{k}\stackrel{}{q})^2}}{(kq)_0^2+(\stackrel{}{k}\stackrel{}{q})^2}}].`$ (31) The integral over the azimuthal angle $`\varphi `$ is trivial, and we therefore have three integrals to do. We do the remaining angular integral analytically, after making a change of variables. We define $`\stackrel{}{q^{}}=\stackrel{}{k}\stackrel{}{q}`$because the integration over the polar angle $`\theta `$ is simpler when the momentum integration is done over $`\stackrel{}{q^{}}`$. The simplification arises because there is no longer any angular dependence in the functions $`F`$ and $`G`$: $$F(k_0q_0,\|\stackrel{}{k}\stackrel{}{q}\|)=F(k_0q_0,\|\stackrel{}{q^{}}\|)$$ and similarly for $`G`$. After doing the angular integral, the gap equation reduces to a double integral equation with integration variables $`\|\stackrel{}{q^{}}\|`$ (which we henceforth denote $`q`$) and $`q_0`$: $`\mathrm{\Delta }(k_0)`$ $`=`$ $`{\displaystyle \frac{g^2}{48\pi ^3}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}dq_0{\displaystyle _0^{\mathrm{}}}dq[{\displaystyle \frac{\mathrm{\Delta }(q_0)}{(k_0q_0)^2+q^2+G(k_0q_0,q)}}I_G(q_0,q)`$ (32) $`+{\displaystyle \frac{\mathrm{\Delta }(q_0)}{(k_0q_0)^2+q^2+F(k_0q_0,q)}}I_F(k_0,q_0,q)+\xi {\displaystyle \frac{q\mathrm{\Delta }(q_0)}{(k_0q_0)^2+q^2}}I_\xi (k_0,q_0,q)]`$ (33) $`\mathrm{where}`$ (34) $`I_G(q_0,q<\mu )`$ $`=`$ $`{\displaystyle \frac{2(q_0^2+4\mathrm{\Delta }^2(q_0)+q^2)(q^2+4\mu ^2q_0^24\mathrm{\Delta }^2(q_0))}{3q\mu ^2\sqrt{q_0^2+4\mathrm{\Delta }^2(q_0)}}}\mathrm{arctan}{\displaystyle \frac{q}{\sqrt{q_0^2+4\mathrm{\Delta }^2(q_0)}}}`$ (35) $`+`$ $`{\displaystyle \frac{4(q_0^2+\mathrm{\Delta }^2(q_0)+q^2)(q^2+4\mu ^2q_0^2\mathrm{\Delta }^2(q_0))}{3q\mu ^2\sqrt{q_0^2+\mathrm{\Delta }^2(q_0)}}}\mathrm{arctan}{\displaystyle \frac{q}{\sqrt{q_0^2+\mathrm{\Delta }^2(q_0)}}}`$ (36) $`+`$ $`{\displaystyle \frac{12\mathrm{\Delta }^2(q_0)+6q_0^22q^224\mu ^2}{3\mu ^2}}`$ (37) $`I_F(k_0,q_0,q<\mu )=`$ (38) $`{\displaystyle \frac{2((q_0^2+4\mathrm{\Delta }^2(q_0))(k_0q_0)^2q^4)(q^24\mu ^2+q_0^2+4\mathrm{\Delta }^2(q_0))}{3q\mu ^2\sqrt{q_0^2+4\mathrm{\Delta }^2(q_0)}((k_0q_0)^2+q^2)}}\mathrm{arctan}{\displaystyle \frac{q}{\sqrt{q_0^2+4\mathrm{\Delta }^2(q_0)}}}`$ (39) $`+{\displaystyle \frac{4((q_0^2+\mathrm{\Delta }^2(q_0))(k_0q_0)^2q^4)(q^24\mu ^2+q_0^2+\mathrm{\Delta }^2(q_0))}{3q\mu ^2\sqrt{q_0^2+\mathrm{\Delta }^2(q_0)}((k_0q_0)^2+q^2)}}\mathrm{arctan}{\displaystyle \frac{q}{\sqrt{q_0^2+\mathrm{\Delta }^2(q_0)}}}`$ (40) $`+{\displaystyle \frac{6q^4+2(k_0q_0)^2(2q^2+12\mu ^23q_0^26\mathrm{\Delta }^2(q_0))}{3\mu ^2((k_0q_0)^2+q^2)}}`$ (41) $`I_\xi (k_0,q_0,q<\mu )=`$ (42) $`{\displaystyle \frac{2(q_0^2+4\mathrm{\Delta }^2(q_0)(k_0q_0)^2)(q^24\mu ^2+q_0^2+4\mathrm{\Delta }^2(q_0))}{3\mu ^2\sqrt{q_0^2+4\mathrm{\Delta }^2(q_0)}((k_0q_0)^2+q^2)}}\mathrm{arctan}{\displaystyle \frac{q}{\sqrt{q_0^2+4\mathrm{\Delta }^2(q_0)}}}`$ (43) $`{\displaystyle \frac{4(q_0^2+\mathrm{\Delta }^2(q_0)(k_0q_0)^2)(q^24\mu ^2+q_0^2+\mathrm{\Delta }^2(q_0))}{3\mu ^2\sqrt{q_0^2+\mathrm{\Delta }^2(q_0)}((k_0q_0)^2+q^2)}}\mathrm{arctan}{\displaystyle \frac{q}{\sqrt{q_0^2+\mathrm{\Delta }^2(q_0)}}}`$ (44) $`+{\displaystyle \frac{2q(2q^23(k_0q_0)^212\mu ^2+3q_0^2+6\mathrm{\Delta }^2(q_0))}{3\mu ^2((k_0q_0)^2+q^2)}}`$ (45) $`I_G(q_0,q\mu )=`$ (46) $`{\displaystyle \frac{(q_0^2+4\mathrm{\Delta }^2(q_0)+q^2)(q^2+4\mu ^2q_0^24\mathrm{\Delta }^2(q_0))}{3q\mu ^2\sqrt{q_0^2+4\mathrm{\Delta }^2(q_0)}}}(\mathrm{arctan}{\displaystyle \frac{q}{\sqrt{q_0^2+4\mathrm{\Delta }^2(q_0)}}}`$ (47) $`\mathrm{arctan}{\displaystyle \frac{q2\mu }{\sqrt{q_0^2+4\mathrm{\Delta }^2(q_0)}}})`$ (48) $`+{\displaystyle \frac{2(q_0^2+\mathrm{\Delta }^2(q_0)+q^2)(q^2+4\mu ^2q_0^2\mathrm{\Delta }^2(q_0))}{3q\mu ^2\sqrt{q_0^2+\mathrm{\Delta }^2(q_0)}}}(\mathrm{arctan}{\displaystyle \frac{q}{\sqrt{q_0^2+\mathrm{\Delta }^2(q_0)}}}`$ (49) $`\mathrm{arctan}{\displaystyle \frac{q2\mu }{\sqrt{q_0^2+\mathrm{\Delta }^2(q_0)}}})`$ (50) $`+{\displaystyle \frac{4(q_0^2+\mathrm{\Delta }^2(q_0)+q^2)}{3q\mu }}\mathrm{ln}{\displaystyle \frac{q_0^2+\mathrm{\Delta }^2(q_0)+q^2}{q_0^2+\mathrm{\Delta }^2(q_0)+(q2\mu )^2}}`$ (51) $`+{\displaystyle \frac{2(q_0^2+4\mathrm{\Delta }^2(q_0)+q^2)}{3q\mu }}\mathrm{ln}{\displaystyle \frac{q_0^2+4\mathrm{\Delta }^2(q_0)+q^2}{q_0^2+4\mathrm{\Delta }^2(q_0)+(q2\mu )^2}}`$ (52) $`+{\displaystyle \frac{12\mathrm{\Delta }^2(q_0)+6q_0^26q^28\mu ^212\mu q}{3\mu q}}`$ (53) $`I_F(k_0,q_0,q\mu )=`$ (54) $`{\displaystyle \frac{((q_0^2+4\mathrm{\Delta }^2(q_0))(k_0q_0)^2q^4)(q^24\mu ^2+q_0^2+4\mathrm{\Delta }^2(q_0))}{3q\mu ^2\sqrt{q_0^2+4\mathrm{\Delta }^2(q_0)}((k_0q_0)^2+q^2)}}(\mathrm{arctan}{\displaystyle \frac{q}{\sqrt{q_0^2+4\mathrm{\Delta }^2(q_0)}}}`$ (55) $`\mathrm{arctan}{\displaystyle \frac{q2\mu }{\sqrt{q_0^2+4\mathrm{\Delta }^2(q_0)}}})`$ (56) $`+{\displaystyle \frac{2((q_0^2+\mathrm{\Delta }^2(q_0))(k_0q_0)^2q^4)(q^24\mu ^2+q_0^2+\mathrm{\Delta }^2(q_0))}{3q\mu ^2\sqrt{q_0^2+\mathrm{\Delta }^2(q_0)}((k_0q_0)^2+q^2)}}(\mathrm{arctan}{\displaystyle \frac{q}{\sqrt{q_0^2+\mathrm{\Delta }^2(q_0)}}}`$ (57) $`\mathrm{arctan}{\displaystyle \frac{q2\mu }{\sqrt{q_0^2+\mathrm{\Delta }^2(q_0)}}})`$ (58) $`+{\displaystyle \frac{4(q^4(q_0^2+\mathrm{\Delta }^2(q_0))(k_0q_0)^2)}{3q\mu (q^2+(k_0q_0)^2)}}\mathrm{ln}{\displaystyle \frac{q_0^2+\mathrm{\Delta }^2(q_0)+q^2}{q_0^2+\mathrm{\Delta }^2(q_0)+(q2\mu )^2}}`$ (59) $`+{\displaystyle \frac{2(q^4(q_0^2+4\mathrm{\Delta }^2(q_0))(k_0q_0)^2)}{3q\mu (q^2+(k_0q_0)^2)}}\mathrm{ln}{\displaystyle \frac{q_0^2+4\mathrm{\Delta }^2(q_0)+q^2}{q_0^2+4\mathrm{\Delta }^2(q_0)+(q2\mu )^2}}`$ (60) $`+{\displaystyle \frac{6q^4+2(k_0q_0)^2(6\mu q+4\mu ^23q_0^26\mathrm{\Delta }^2(q_0))}{3\mu q((k_0q_0)^2+q^2)}}`$ (61) $`I_\xi (k_0,q_0,q\mu )=`$ (62) $`{\displaystyle \frac{(q_0^2+4\mathrm{\Delta }^2(q_0)(k_0q_0)^2)(q^24\mu ^2+q_0^2+4\mathrm{\Delta }^2(q_0))}{3\mu ^2\sqrt{q_0^2+4\mathrm{\Delta }^2(q_0)}((k_0q_0)^2+q^2)}}(\mathrm{arctan}{\displaystyle \frac{q}{\sqrt{q_0^2+4\mathrm{\Delta }^2(q_0)}}}`$ (63) $`\mathrm{arctan}{\displaystyle \frac{q2\mu }{\sqrt{q_0^2+4\mathrm{\Delta }^2(q_0)}}})`$ (64) $`{\displaystyle \frac{2(q_0^2+\mathrm{\Delta }^2(q_0)(k_0q_0)^2)(q^24\mu ^2+q_0^2+\mathrm{\Delta }^2(q_0))}{3\mu ^2\sqrt{q_0^2+\mathrm{\Delta }^2(q_0)}((k_0q_0)^2+q^2)}}(\mathrm{arctan}{\displaystyle \frac{q}{\sqrt{q_0^2+\mathrm{\Delta }^2(q_0)}}}`$ (65) $`\mathrm{arctan}{\displaystyle \frac{q2\mu }{\sqrt{q_0^2+\mathrm{\Delta }^2(q_0)}}})`$ (66) $`+{\displaystyle \frac{4(q_0^2+\mathrm{\Delta }^2(q_0)(k_0q_0)^2)}{3\mu (q^2+(k_0q_0)^2)}}\mathrm{ln}{\displaystyle \frac{q_0^2+\mathrm{\Delta }^2(q_0)+q^2}{q_0^2+\mathrm{\Delta }^2(q_0)+(q2\mu )^2}}`$ (67) $`+{\displaystyle \frac{2(q_0^2+4\mathrm{\Delta }^2(q_0)(k_0q_0)^2)}{3\mu (q^2+(k_0q_0)^2)}}\mathrm{ln}{\displaystyle \frac{q_0^2+4\mathrm{\Delta }^2(q_0)+q^2}{q_0^2+4\mathrm{\Delta }^2(q_0)+(q2\mu )^2}}`$ (68) $`+{\displaystyle \frac{6q_0^2+12\mathrm{\Delta }^2(q_0)6(k_0q_0)^28\mu ^212\mu q)}{3\mu ((k_0q_0)^2+q^2)}}`$ (69) We have solved the gap equation (32) numerically for several different values of $`g`$ and several different values of $`\xi `$. It is convenient to change integration variables from $`q_0`$ to $`\mathrm{ln}q_0`$ and from $`q`$ to $`\mathrm{ln}q`$. We evaluate the $`q`$ integral over a range $`q_{\mathrm{min}}<q<10^4\mu `$ with $`q_{\mathrm{min}}/\mu `$ chosen differently for each $`g`$ in such a way that it is less than $`10^5\mathrm{\Delta }(0)`$ in all cases. The $`q_0`$ integral is made even in $`q_0`$ (by taking the average of the integrand at $`q_0`$ and $`q_0`$) and then evaluated over a range $`q_{0\mathrm{m}\mathrm{i}\mathrm{n}}<q_0<100\mu `$, where we chose $`q_{0\mathrm{m}\mathrm{i}\mathrm{n}}=q_{\mathrm{min}}`$. We have checked that our results are insensitive to the choice of upper and lower cutoffs of the integration region. It was probably not necessary to choose $`q_{\mathrm{min}}`$ and $`q_{0\mathrm{m}\mathrm{i}\mathrm{n}}`$ quite as small as we did. It is, however, quite important to extend the upper limit of the $`q_0`$ and $`q`$ integrals to well above $`\mu `$ in order to avoid sensitivity to the ultraviolet cutoff.The one exception, in which we do find some sensitivity to one of our limits of integration, is at $`g=3.5576`$. With $`g`$ this large, we should perhaps have extended the upper cutoff of the $`q_0`$ integration to 1000 $`\mu `$, as the results shown in Fig. 1 below make clear. We use an iterative method, in which an initial guess for $`\mathrm{\Delta }(k_0)`$ is used on the right-hand side of (32), the integrals are done yielding a new $`\mathrm{\Delta }(k_0)`$, which is in turn used on the right-hand side. The solution converges well after about ten iterations. All results we show were iterated at least fifteen times. Our results are shown in Fig. 1. Note that the output of our calculation is a plot of $`\mathrm{\Delta }(q_0)/\mu `$ as a function of $`q_0/\mu `$ for some choice of $`g`$ and $`\xi `$. The only way in which $`\mu `$ enters the calculation is to set the units of energy. The values of $`\mu `$ shown in Fig. 1 corresponding to each value of $`g`$ do not come from the calculation. They are obtained by assuming that the running coupling $`g`$ should be evaluated at the scale $`\mu `$ and using the one-loop beta function with $`\mathrm{\Lambda }_{\mathrm{QCD}}=200`$ MeV. We include these values of $`\mu `$ to make comparison with the results of Refs. easier. If, as seems quite reasonable, $`g`$ should in fact be evaluated at a $`g`$-dependent scale which is lower than $`\mu `$, then the values of $`g`$ at which we have done our calculations correspond to larger values of $`\mu `$ than shown in Fig. 1 . Evans, Hormuzdiar, Hsu, and Schwetz have obtained numerical solutions to simplified gap equations describing the gap in the CFL phase . Their results agree reasonably well with the results of our calculation done in $`\xi =0`$ gauge but disagree qualitatively with ours in any other gauge. Simply setting $`\xi =0`$, as in Ref. , is not a valid approximation at the values of $`g`$ at which we (and these authors) work. How should one interpret the results of a gauge dependent calculation, given that at any fixed $`g`$ one can obtain any result one likes if one is willing to explore gauge parameters $`\mathrm{}<\xi <\mathrm{}`$? In the present circumstance, the idea is that we expect this calculation to give a gauge invariant result in the $`g0`$ limit. More precisely, if we define $$f(g)\mathrm{ln}\left[\frac{\mathrm{\Delta }(0)}{\mu }\right]+\frac{3\pi ^2}{\sqrt{2}g}+5\mathrm{ln}g$$ (70) then we expect $`f`$ to go to a $`\xi `$-independent constant in the $`g0`$ limit. In Fig. 2, we plot $`f(g)`$ in five different gauges. From this figure we learn: * For any $`\xi `$, $`f(g)`$ is a reasonably slowly varying function of $`g`$. This confirms Son’s result (1) and justifies an analysis in terms of $`f(g)`$. * It does appear that $`lim_{g0}f(g)`$ is a $`\xi `$-independent constant, perhaps not far from the estimate of Ref. , namely $`lim_{g0}f(g)=8.88`$, or that of Ref. , namely $`lim_{g0}f(g)=7.84`$. * If we do a calculation in some fixed gauge, we expect that at small enough $`g`$ this calculation yields a good estimate of the true gauge invariant result. By doing calculations in several gauges, we can bound the regime of applicability of this estimate. We can only trust our calculation of $`f(g)`$ in the regime in which the $`\xi `$-dependence of $`f`$ decreases with decreasing $`g`$. Our calculation of $`f(g)`$ is completely meaningless unless $`g`$ is small enough that the curves for different values of $`\xi `$ are converging. Fig. 2 shows that the gauge dependence of our result for $`f`$ is about the same for all $`g0.8`$. It is only for $`g0.8`$ that $`f(g)`$ calculated in different gauges begins to converge. At larger values of $`g`$ our calculation provides no guide whatsoever as to the value of $`f`$ that would be obtained in a complete, gauge invariant calculation including all the physics neglected in the present calculation. Even at $`g=0.8`$ the values of $`\mathrm{\Delta }(0)`$ differ by a factor of about 400 for gaps with $`\xi =4`$ and $`\xi =4`$. We could make the gauge dependence look even larger by choosing larger values of $`\|\xi \|`$. Our result does not guarantee that the calculation is under control for $`g<0.8`$, but it does guarantee that the result is uncontrolled and completely meaningless for $`g>0.8`$. ## IV Conclusion We have detailed our assumptions and approximations as we made them. Let us now ask which of them should be improved upon if we wish to include those contributions whose neglect we have diagnosed via the gauge dependence of our results. Note that $`g=0.8`$ corresponds to $`\mathrm{\Delta }/\mu 10^7`$. Thus, those contributions to $`f`$ which we have neglected which are controlled when $`\mathrm{\Delta }\mu `$ are not responsible for the breakdown of our calculation around $`g0.8`$. We believe that the assumptions we made in writing the ansatz (5) and the assumptions we made in neglecting $`\mathrm{\Delta }_1^S`$ and $`\mathrm{\Delta }_2^{A,S}`$ all introduce errors which are small when $`\mathrm{\Delta }\mu `$. (For example, even though neglecting $`\mathrm{\Delta }_2`$ is a source of gauge dependence, we do not expect that remedying this neglect would change $`f(g)`$ appreciably in any gauge at $`g0.8`$, where $`\mathrm{\Delta }/\mu `$ is so small.) Hence, we believe that it is the assumptions made in writing the truncated gap equation (9) that are at fault. One obvious possible explanation is the absence of vertex corrections, although there are other missing skeleton diagrams which should also be investigated. The gap $`\mathrm{\Delta }`$ is of course a gauge invariant observable. A complete calculation would yield a gauge invariant expression for the function $`f`$, which could be expanded as a power series in $`g`$. We learn three things from our (incomplete and gauge dependent) calculation. First, our results obtained in different gauges appear to converge at small $`g`$ and support previous estimates of $`lim_{g0}f(g)`$, namely the $`g^0`$ term in the expansion of $`f`$. Second, because the results we obtain in different gauges only begin to converge for $`g<g_c0.8`$, we learn that contributions to our gauge dependent function $`f`$ which are of order $`g^1`$ and higher must have gauge dependent parts which are numerically large at $`gg_c`$. Although we have simply evaluated $`f(g)`$ and not expanded it in $`g`$, we learn that such an expansion is uncontrolled for $`g>g_c`$. This suggests that if we knew the complete, gauge invariant function $`f`$, the $`g^1`$ and higher terms in that expansion would also become uncontrolled for $`g>g_c`$. It may be that the vertex corrections are the dominant contribution to the missing physics which is responsible for this breakdown: this hypothesis is supported by the arguments of Ref. that these effects contribute to $`f`$ at order $`g^1`$. Regardless of whether the vertex corrections turn out to be the most important effect left out of the truncated gap equation (9), our calculation demonstrates that some contribution which is formally subleading is in fact large enough to render the calculation uncontrolled at $`gg_c`$. The third thing we learn is that although present calculations do yield reasonable estimates of $`lim_{g0}f(g)`$, if one is interested in using these calculations to estimate the value of $`\mathrm{\Delta }`$ to within a factor of two, this can only be done for $`gg_c0.8`$. In the CFL phase, all eight gluons get a mass. This means that in the CFL phase there are no gapless fermionic excitations, and no massless gluonic excitations, and therefore no non-Abelian physics in the infrared to obstruct weak-coupling calculations. The lesson we have learned is that even though everything is in principle under control, present weak-coupling calculations break down for $`g>g_c0.8`$, corresponding to $`\mu <\mu _c`$ with $`\mu _c10^8`$ MeV (or higher ). This break down occurs even though $`\mathrm{\Delta }\mu `$ at $`gg_c`$. It should be noted that what breaks down is the weak-coupling calculation of the magnitude of the gap $`\mathrm{\Delta }`$. Estimates based on models normalized to give reasonable zero density phenomenology can still be used as a guide, albeit a qualitative one. Furthermore, regardless of the fact that a controlled calculation of $`\mathrm{\Delta }`$ has not yet been done at $`\mu <10^8`$ MeV, it is possible to construct a controlled effective field theory which describes the infrared physics of the CFL phase on length scales long compared to $`1/\mathrm{\Delta }`$, since in such an effective theory $`\mathrm{\Delta }`$ is simply a parameter determined by physics outside the effective theory. This infrared physics is dominated by the massless Abelian gauge bosons , the Nambu-Goldstone boson arising from spontaneously broken $`U(1)_B`$ , and the pseudo-Nambu-Goldstone bosons arising from spontaneously broken chiral symmetry which have small masses due to the nonvanishing quark masses . ###### Acknowledgements. We thank I. Shovkovy for suggesting that gauge dependence could be used as a diagnostic device and thank T. Schaefer for very helpful discussions. We are grateful to the Department of Energy’s Institute for Nuclear Theory at the University of Washington for generous hospitality and support during the completion of this work. This research is also supported in part by the Department of Energy under cooperative research agreement DF-FC02-94ER40818. The work of KR is supported in part by by a DOE OJI grant and by the Alfred P. Sloan Foundation. ## A The Meissner Effect In this appendix, we set up the calculation of the Meissner effect. That is, we investigate the effect of the presence of a gap $`\mathrm{\Delta }`$ on the functions $`F`$ and $`G`$ which describe the screening of the gluon propagator. In order to establish some necessary notation, we must begin by filling in some details in the derivation of Eq. (11) from Eq. (9). We work in a color-flavor basis $`(\{i,a\},\{j,b\})`$. In this basis, we define the following two $`9\times 9`$ matrices: $$Q_{ij}^{ab}=(\lambda _I^A)^{ab}(\lambda _I^A)_{ij}=\left(\begin{array}{ccccccccc}0& 1& 1& & & & & & \\ 1& 0& 1& & & & & & \\ 1& 1& 0& & & & & & \\ & & & 0& 1& & & & \\ & & & 1& 0& & & & \\ & & & & & 0& 1& & \\ & & & & & 1& 0& & \\ & & & & & & & 0& 1\\ & & & & & & & 1& 0\end{array}\right)$$ (A1) $$R_{ij}^{ab}=(\lambda _J^S)^{ab}(\lambda _J^S)_{ij}=\left(\begin{array}{ccccccccc}2& 1& 1& & & & & & \\ 1& 2& 1& & & & & & \\ 1& 1& 2& & & & & & \\ & & & 0& 1& & & & \\ & & & 1& 0& & & & \\ & & & & & 0& 1& & \\ & & & & & 1& 0& & \\ & & & & & & & 0& 1\\ & & & & & & & 1& 0\end{array}\right)$$ (A2) which represent the antisymmetric color and flavor $`\overline{\mathrm{𝟑}}_A`$ and the symmetric color and flavor $`\mathrm{𝟔}_S`$ channels respectively in this basis. In the derivation of the gap equation, we were only interested in the off-diagonal lower left component of the Nambu-Gorkov fermion propagator $`S`$. However, the calculation of the Meissner effect involves all components of the fermion propagator. Obtaining the fermion propagator by inverting the inverse propagator (4) is straightforward but tedious. After a lot of algebra and using the ansatz (5) for the gap matrix, we find: $`S(q)=\left(\begin{array}{cc}S_{11}(q)& S_{12}(q)\\ S_{21}(q)& S_{22}(q)\end{array}\right)`$ (A5) where $`S_{11}(q)=\left(\begin{array}{ccccccccc}A(q)& B(q)& B(q)& & & & & & \\ B(q)& A(q)& B(q)& & & & & & \\ B(q)& B(q)& A(q)& & & & & & \\ & & & C(q)& & & & & \\ & & & & C(q)& & & & \\ & & & & & C(q)& & & \\ & & & & & & C(q)& & \\ & & & & & & & C(q)& \\ & & & & & & & & C(q)\end{array}\right)`$ (A15) $`S_{22}(q)=\left(\begin{array}{ccccccccc}E(q)& H(q)& H(q)& & & & & & \\ H(q)& E(q)& H(q)& & & & & & \\ H(q)& H(q)& E(q)& & & & & & \\ & & & D(q)& & & & & \\ & & & & D(q)& & & & \\ & & & & & D(q)& & & \\ & & & & & & D(q)& & \\ & & & & & & & D(q)& \\ & & & & & & & & D(q)\end{array}\right)`$ (A25) $`S_{21}(q)=\overline{S_{12}}(q)=\left(\begin{array}{ccccccccc}K(q)& L(q)& L(q)& & & & & & \\ L(q)& K(q)& L(q)& & & & & & \\ L(q)& L(q)& K(q)& & & & & & \\ & & & 0& M(q)& & & & \\ & & & M(q)& 0& & & & \\ & & & & & 0& M(q)& & \\ & & & & & M(q)& 0& & \\ & & & & & & & 0& M(q)\\ & & & & & & & M(q)& 0\end{array}\right)`$ (A35) and where the above functions are defined as follows: $$\begin{array}{cc}A(q)\hfill & =\gamma ^0[P_+(q)\frac{q_0\mu \|\stackrel{}{q}\|}{q_0^2(\|\stackrel{}{q}\|+\mu )^24(\mathrm{\Delta }_2^A(q_0)+2\mathrm{\Delta }_2^S(q_0))^2}\frac{q_0^2(\|\stackrel{}{q}\|+\mu )^23\left(\mathrm{\Delta }_2^A(q_0)\right)^211\left(\mathrm{\Delta }_2^S(q_0)\right)^210\mathrm{\Delta }_2^A(q_0)\mathrm{\Delta }_2^S(q_0)}{q_0^2(\|\stackrel{}{q}\|+\mu )^2(\mathrm{\Delta }_2^A(q_0)\mathrm{\Delta }_2^S(q_0))^2}\hfill \\ & +P_{}(q)\frac{q_0\mu +\|\stackrel{}{q}\|}{q_0^2(\|\stackrel{}{q}\|\mu )^24(\mathrm{\Delta }_1^A(q_0)+2\mathrm{\Delta }_1^S(q_0))^2}\frac{q_0^2(\|\stackrel{}{q}\|\mu )^23\left(\mathrm{\Delta }_1^A(q_0)\right)^211\left(\mathrm{\Delta }_1^S(q_0)\right)^210\mathrm{\Delta }_1^A(q_0)\mathrm{\Delta }_1^S(q_0)}{q_0^2(\|\stackrel{}{q}\|\mu )^2(\mathrm{\Delta }_1^A(q_0)\mathrm{\Delta }_1^S(q_0))^2}]\hfill \\ & \\ B(q)\hfill & =\gamma ^0[P_+(q)\frac{q_0\mu \|\stackrel{}{q}\|}{q_0^2(\|\stackrel{}{q}\|+\mu )^24(\mathrm{\Delta }_2^A(q_0)+2\mathrm{\Delta }_2^S(q_0))^2}\frac{\left(\mathrm{\Delta }_2^A(q_0)+5\mathrm{\Delta }_2^S(q_0)\right)\left(\mathrm{\Delta }_2^A(q_0)+\mathrm{\Delta }_2^S(q_0)\right)}{q_0^2(\|\stackrel{}{q}\|+\mu )^2(\mathrm{\Delta }_2^A(q_0)\mathrm{\Delta }_2^S(q_0))^2}\hfill \\ & +P_{}(q)\frac{q_0\mu +\|\stackrel{}{q}\|}{q_0^2(\|\stackrel{}{q}\|\mu )^24(\mathrm{\Delta }_1^A(q_0)+2\mathrm{\Delta }_1^S(q_0))^2}\frac{\left(\mathrm{\Delta }_1^A(q_0)+5\mathrm{\Delta }_1^S(q_0)\right)\left(\mathrm{\Delta }_1^A(q_0)+\mathrm{\Delta }_1^S(q_0)\right)}{q_0^2(\|\stackrel{}{q}\|\mu )^2(\mathrm{\Delta }_1^A(q_0)\mathrm{\Delta }_1^S(q_0))^2}]\hfill \end{array}$$ (A36) $$\begin{array}{cc}C(q)\hfill & =\gamma ^0\left[P_+(q)\frac{q_0\mu \|\stackrel{}{q}\|}{q_0^2(\|\stackrel{}{q}\|+\mu )^2(\mathrm{\Delta }_2^A(q_0)\mathrm{\Delta }_2^S(q_0))^2}+P_{}(q)\frac{q_0\mu +\|\stackrel{}{q}\|}{q_0^2(\|\stackrel{}{q}\|\mu )^2(\mathrm{\Delta }_1^A(q_0)\mathrm{\Delta }_1^S(q_0))^2}\right]\hfill \\ & \\ D(q)\hfill & =C\gamma ^0\left[P_{}(q)\frac{q_0+\mu +\|\stackrel{}{q}\|}{q_0^2(\|\stackrel{}{q}\|+\mu )^2(\mathrm{\Delta }_2^A(q_0)\mathrm{\Delta }_2^S(q_0))^2}+P_+(q)\frac{q_0+\mu \|\stackrel{}{q}\|}{q_0^2(\|\stackrel{}{q}\|\mu )^2(\mathrm{\Delta }_1^A(q_0)\mathrm{\Delta }_1^S(q_0))^2}\right]C\hfill \\ & \\ E(q)\hfill & =C\gamma ^0[P_{}(q)\frac{q_0+\mu +\|\stackrel{}{q}\|}{q_0^2(\|\stackrel{}{q}\|+\mu )^24(\mathrm{\Delta }_2^A(q_0)+2\mathrm{\Delta }_2^S(q_0))^2}\frac{q_0^2(\|\stackrel{}{q}\|+\mu )^23\left(\mathrm{\Delta }_2^A(q_0)\right)^211\left(\mathrm{\Delta }_2^S(q_0)\right)^210\mathrm{\Delta }_2^A(q_0)\mathrm{\Delta }_2^S(q_0)}{q_0^2(\|\stackrel{}{q}\|+\mu )^2(\mathrm{\Delta }_2^A(q_0)\mathrm{\Delta }_2^S(q_0))^2}\hfill \\ & +P_+(q)\frac{q_0+\mu \|\stackrel{}{q}\|}{q_0^2(\|\stackrel{}{q}\|\mu )^24(\mathrm{\Delta }_1^A(q_0)+2\mathrm{\Delta }_1^S(q_0))^2}\frac{q_0^2(\|\stackrel{}{q}\|\mu )^23\left(\mathrm{\Delta }_1^A(q_0)\right)^211\left(\mathrm{\Delta }_1^S(q_0)\right)^210\mathrm{\Delta }_1^A(q_0)\mathrm{\Delta }_1^S(q_0)}{q_0^2(\|\stackrel{}{q}\|\mu )^2(\mathrm{\Delta }_1^A(q_0)\mathrm{\Delta }_1^S(q_0))^2}]C\hfill \\ & \\ H(q)\hfill & =C\gamma ^0[P_{}(q)\frac{q_0+\mu +\|\stackrel{}{q}\|}{q_0^2(\|\stackrel{}{q}\|+\mu )^24(\mathrm{\Delta }_2^A(q_0)+2\mathrm{\Delta }_2^S(q_0))^2}\frac{\left(\mathrm{\Delta }_2^A(q_0)+5\mathrm{\Delta }_2^S(q_0)\right)\left(\mathrm{\Delta }_2^A(q_0)+\mathrm{\Delta }_2^S(q_0)\right)}{q_0^2(\|\stackrel{}{q}\|+\mu )^2(\mathrm{\Delta }_2^A(q_0)\mathrm{\Delta }_2^S(q_0))^2}\hfill \\ & +P_+(q)\frac{q_0+\mu \|\stackrel{}{q}\|}{q_0^2(\|\stackrel{}{q}\|\mu )^24(\mathrm{\Delta }_1^A(q_0)+2\mathrm{\Delta }_1^S(q_0))^2}\frac{\left(\mathrm{\Delta }_1^A(q_0)+5\mathrm{\Delta }_1^S(q_0)\right)\left(\mathrm{\Delta }_1^A(q_0)+\mathrm{\Delta }_1^S(q_0)\right)}{q_0^2(\|\stackrel{}{q}\|\mu )^2(\mathrm{\Delta }_1^A(q_0)\mathrm{\Delta }_1^S(q_0))^2}]C\hfill \\ & \\ K(q)\hfill & =2C\gamma ^5[P_+(q)(\frac{\mathrm{\Delta }_2^S(q_0)}{q_0^2(\|\stackrel{}{q}\|+\mu )^24(\mathrm{\Delta }_2^A(q_0)+2\mathrm{\Delta }_2^S(q_0))^2}\hfill \\ & +\frac{\mathrm{\Delta }_2^A(q_0)\mathrm{\Delta }_2^S(q_0)}{q_0^2(\|\stackrel{}{q}\|+\mu )^24(\mathrm{\Delta }_2^A(q_0)+2\mathrm{\Delta }_2^S(q_0))^2}\frac{\left(\mathrm{\Delta }_2^A(q_0)+5\mathrm{\Delta }_2^S(q_0)\right)\left(\mathrm{\Delta }_2^A(q_0)+\mathrm{\Delta }_2^S(q_0)\right)}{q_0^2(\|\stackrel{}{q}\|+\mu )^2(\mathrm{\Delta }_2^A(q_0)\mathrm{\Delta }_2^S(q_0))^2})\hfill \\ & +P_{}(q)(\frac{\mathrm{\Delta }_1^S(q_0)}{q_0^2(\|\stackrel{}{q}\|\mu )^24(\mathrm{\Delta }_1^A(q_0)2\mathrm{\Delta }_1^S(q_0))^2}\hfill \\ & +\frac{\mathrm{\Delta }_1^A(q_0)\mathrm{\Delta }_1^S(q_0)}{q_0^2(\|\stackrel{}{q}\|\mu )^24(\mathrm{\Delta }_1^A(q_0)+2\mathrm{\Delta }_1^S(q_0))^2}\frac{\left(\mathrm{\Delta }_1^A(q_0)+5\mathrm{\Delta }_1^S(q_0)\right)\left(\mathrm{\Delta }_1^A(q_0)+\mathrm{\Delta }_1^S(q_0)\right)}{q_0^2(\|\stackrel{}{q}\|\mu )^2(\mathrm{\Delta }_1^A(q_0)\mathrm{\Delta }_1^S(q_0))^2})]\hfill \\ & \\ L(q)\hfill & =C\gamma ^5[P_+(q)(\frac{\mathrm{\Delta }_2^S(q_0)+\mathrm{\Delta }_2^A(q_0)}{q_0^2(\|\stackrel{}{q}\|+\mu )^24(\mathrm{\Delta }_2^A(q_0)+2\mathrm{\Delta }_2^S(q_0))^2}\hfill \\ & +\frac{\mathrm{\Delta }_2^A(q_0)+\mathrm{\Delta }_2^S(q_0)}{q_0^2(\|\stackrel{}{q}\|+\mu )^24(\mathrm{\Delta }_2^A(q_0)+2\mathrm{\Delta }_2^S(q_0))^2}\frac{\left(\mathrm{\Delta }_2^A(q_0)+5\mathrm{\Delta }_2^S(q_0)\right)\left(\mathrm{\Delta }_2^A(q_0)+\mathrm{\Delta }_2^S(q_0)\right)}{q_0^2(\|\stackrel{}{q}\|+\mu )^2(\mathrm{\Delta }_2^A(q_0)\mathrm{\Delta }_2^S(q_0))^2})\hfill \\ & +P_{}(q)(\frac{\mathrm{\Delta }_1^S(q_0)+\mathrm{\Delta }_1^A(q_0)}{q_0^2(\|\stackrel{}{q}\|\mu )^24(\mathrm{\Delta }_1^A(q_0)+2\mathrm{\Delta }_1^S(q_0))^2}\hfill \\ & +\frac{\mathrm{\Delta }_1^A(q_0)+\mathrm{\Delta }_1^S(q_0)}{q_0^2(\|\stackrel{}{q}\|\mu )^24(\mathrm{\Delta }_1^A(q_0)+2\mathrm{\Delta }_1^S(q_0))^2}\frac{\left(\mathrm{\Delta }_1^A(q_0)+5\mathrm{\Delta }_1^S(q_0)\right)\left(\mathrm{\Delta }_1^A(q_0)+\mathrm{\Delta }_1^S(q_0)\right)}{q_0^2(\|\stackrel{}{q}\|\mu )^2(\mathrm{\Delta }_1^A(q_0)\mathrm{\Delta }_1^S(q_0))^2})]\hfill \\ & \\ M(q)\hfill & =C\gamma ^5\left[P_+(q)\frac{\mathrm{\Delta }_2^A(q_0)+\mathrm{\Delta }_2^S(q_0)}{q_0^2(\|\stackrel{}{q}\|+\mu )^2(\mathrm{\Delta }_2^A(q_0)\mathrm{\Delta }_2^S(q_0))^2}+P_{}(q)\frac{\mathrm{\Delta }_1^A(q_0)+\mathrm{\Delta }_1^S(q_0)}{q_0^2(\|\stackrel{}{q}\|\mu )^2(\mathrm{\Delta }_1^A(q_0)\mathrm{\Delta }_1^S(q_0))^2}\right].\hfill \end{array}$$ (A36) Note that $`S_{21}(q)=\overline{S_{12}}(q)`$ is a general property of the Fermion propagator $`S`$ and can be proved for an arbitrary number of colors and flavors using only the definition of the inverse Fermion propagator, Eq. (4), and properties of the Dirac gamma matrices. Whereas only $`K`$, $`L`$ and $`M`$ were used in the derivation of the gap equation, all these functions are required in evaluating the Meissner effect. The Meissner effect is the change in the screening of the gluon propagator induced by the presence of a gap. To one loop order, we need to evaluate the gluon propagator of Fig. 3 using the full fermion propagator including the gap. The result can still be written in the form (21) but now $$F(q)=F_0(q)+\delta F(q)\mathrm{and}G(q)=G_0(q)+\delta G(q)$$ (A37) where $`F_0`$ and $`G_0`$ are the $`\mathrm{\Delta }=0`$ functions written as $`F`$ and $`G`$ in (23). Recall that $`G_0`$, which describes Landau damping, vanishes for $`q_00`$. Because $`\delta G`$ is nonzero in the $`q_00`$ limit, the Meissner effect can be described as giving a mass to the gluons. Previous analyses of the Meissner effect have either been done for two-flavor QCD or have used simplified estimates . Our goal is to formulate the correct calculation of $`\delta F(q)`$ and $`\delta G(q)`$ in the CFL phase. Recent work along the same lines can be found in Ref. . From the diagram of Fig. 3, we obtain the gluon polarization $$\begin{array}{ccc}\mathrm{\Pi }_{ab}^{\mu \nu }\hfill & =ig^2\frac{d^4k}{(2\pi )^4}\mathrm{Tr}\hfill & \left[\mathrm{\Gamma }_a^\mu S(k+q)\mathrm{\Gamma }_b^\nu S(k)\right]\hfill \\ & =ig^2\frac{d^4k}{(2\pi )^4}\mathrm{Tr}\hfill & [\gamma ^\mu \frac{\lambda _a}{2}S_{11}(k+q)\gamma ^\nu \frac{\lambda _b}{2}S_{11}(k)+\left(\gamma ^\mu \frac{\lambda _a}{2}\right)^TS_{22}(k+q)\left(\gamma ^\nu \frac{\lambda _b}{2}\right)^TS_{22}(k)\hfill \\ & & \gamma ^\mu \frac{\lambda _a}{2}S_{12}(k+q)\left(\gamma ^\nu \frac{\lambda _b}{2}\right)^TS_{21}(k)\left(\gamma ^\mu \frac{\lambda _a}{2}\right)^TS_{21}(k+q)\gamma ^\nu \frac{\lambda _b}{2}S_{12}(k)],\hfill \end{array}$$ (A38) where the trace is taken over color, flavor, and Dirac indices and all four elements of the fermion propagator, $`S(q)`$, have been defined previously in Eqs. (A15) – (A36). This polarization amplitude contains all the one loop contributions to the gluon propagator including the gap independent contributions, $`F_0(q)`$ and $`G_0(q)`$. $`\mathrm{\Pi }_{ab}^{\mu \nu }`$ can be written in terms of $`F`$ and $`G`$ in a simple fashion: $$\mathrm{\Pi }_{ab}^{\mu \nu }=\delta _{ab}\left[\left(G_0(q)+\delta G(q)\right)P^{\mu \nu T}+\left(F_0(q)+\delta F(q)\right)P^{\mu \nu L}\right].$$ (A39) Hence, we only need to compute two components of $`\mathrm{\Pi }_{ab}^{\mu \nu }`$ in order to obtain the functions $`\delta F(q)`$ and $`\delta G(q)`$, for example, $`\mathrm{\Pi }_{33}^{00}`$ and $`\mathrm{\Pi }_{33}^{11}`$. Because we already know $`F_0(q)`$ and $`G_0(q)`$, our goal is to extract $`\delta F(q)`$ and $`\delta G(q)`$. We are therefore only interested in the difference $`\mathrm{\Pi }_{ab}^{\mu \nu }(\mathrm{\Delta }0)\mathrm{\Pi }_{ab}^{\mu \nu }(\mathrm{\Delta }=0)`$. Finally, because $`\delta F(q)`$ and $`\delta G(q)`$ depend only on $`q_0`$ and $`\|\stackrel{}{q}\|`$, we can choose $`\stackrel{}{q}`$ to lie along the $`z`$-axis for simplicity. Keeping all this in mind, we find that (in Euclidean space) $`\delta F(q)={\displaystyle \frac{q_0^2+\|\stackrel{}{q}\|^2}{\|\stackrel{}{q}\|^2}}\left(\mathrm{\Pi }_{33}^{00}(\mathrm{\Delta }0)\mathrm{\Pi }_{33}^{00}(\mathrm{\Delta }=0)\right)`$ (A40) $`\delta G(q)=\mathrm{\Pi }_{33}^{11}(\mathrm{\Delta }0)\mathrm{\Pi }_{33}^{11}(\mathrm{\Delta }=0).`$ (A41) Note that (unlike the integrals which arise on the right hand side of the gap equation) the integrals which must be done in evaluating $`\mathrm{\Pi }(q)`$ are ultraviolet divergent, and therefore sensitive to how they are cutoff at large $`k_0`$ and $`k`$. This ultraviolet divergence has nothing to do with $`\mathrm{\Delta }`$, and is canceled in our calculation of $`\delta F`$ and $`\delta G`$ by subtracting the $`\mathrm{\Delta }=0`$ result for $`\mathrm{\Pi }(q)`$. We have checked that our results for $`\delta F`$ and $`\delta G`$ are insensitive to the ultraviolet cutoffs in the integrals. Looking back at the definition of $`\mathrm{\Pi }_{ab}^{\mu \nu }`$, we can see that it depends on $`\mathrm{\Delta }_1^{A,S}(k_0)`$ and $`\mathrm{\Delta }_2^{A,S}(k_0)`$. We make the same assumptions here as in our solution of the gap equation, namely that the antiparticle and sextet contributions can be neglected if $`\mathrm{\Delta }\mu `$ and if one is interested in physics dominated by particles and holes near the Fermi surface. Before we proceed, let us define the following notation for the functions $`A(q)`$ through $`M(q)`$ defined in Eq. (A36): identify the scalar functions multiplying the $`P_\pm `$ projectors with the appropriate $`\pm `$ signs, e.g. $`A_+(q)`$. With this notation, the dominant contributions to the two polarization amplitudes we are interested in are: $$\begin{array}{cc}\mathrm{\Pi }_{33}^{00}=\frac{i}{2}g^2\frac{d^4k}{(2\pi )^4}\hfill & (1+\widehat{(k+q)}\widehat{k})[A_{}(k+q)A_{}(k)B_{}(k+q)B_{}(k)\hfill \\ & +2C_{}(k+q)C_{}(k)+E_+(k+q)E_+(k)H_+(k+q)H_+(k)\hfill \\ & +2D_+(k+q)D_+(k)2K_{}(k+q)K_{}(k)\hfill \\ & +2L_{}(k+q)L_{}(k)+2M_{}(k+q)M_{}(k)]\hfill \\ \mathrm{\Pi }_{33}^{11}=\frac{i}{2}g^2\frac{d^4k}{(2\pi )^4}\hfill & (1+2\widehat{(k+q)}^1\widehat{k}^1\widehat{(k+q)}\widehat{k})[A_{}(k+q)A_{}(k)\hfill \\ & B_{}(k+q)B_{}(k)+2C_{}(k+q)C_{}(k)+E_+(k+q)E_+(k)\hfill \\ & H_+(k+q)H_+(k)+2D_+(k+q)D_+(k)+2K_{}(k+q)K_{}(k)\hfill \\ & 2L_{}(k+q)L_{}(k)2M_{}(k+q)M_{}(k)].\hfill \end{array}$$ (A42) In any one gauge, i.e. for a particular choice of $`\xi `$, our task is now clear. We first calculate $`\mathrm{\Delta }(k_0)`$ with $`\delta F(q)=\delta G(q)=0`$, as described in the body of the paper. We must then use (A42) to evaluate $`\delta F(q)`$ and $`\delta G(q)`$ given by Eq. (A40). As in the calculation of $`\mathrm{\Delta }`$, we can do all angular integrals analytically and evaluate the double integral over $`k_0`$ and $`\|\stackrel{}{k}\|`$ numerically. We must then re-evaluate $`\mathrm{\Delta }(k_0)`$ with the new gluon propagator, modified by the addition of $`\delta F(q)`$ and $`\delta G(q)`$. We must then iterate this procedure, calculating $`\delta F(q)`$ and $`\delta G(q)`$ and then recalculating $`\mathrm{\Delta }(k_0)`$ repeatedly, until all results have converged. We have not carried this program to completion. However, preliminary numerical investigation suggests that, in agreement with arguments and estimates made by others , the change in $`\mathrm{\Delta }`$ arising from the inclusion of $`\delta F`$ and $`\delta G`$ is small. In particular, it appears to be much smaller than the change in $`\mathrm{\Delta }`$ which arises if one changes gauge from $`\xi =1`$ to $`\xi =0`$ to $`\xi =1`$. Perhaps at some extremely small $`g`$, the influence of the Meissner effect on the gap could be larger than the influence of the neglected physics whose absence we diagnose via the gauge dependence of our results. At any $`g`$ at which we have been able to obtain numerical results, however, the Meissner effect is insignificant relative to that which is missing.
warning/0004/cond-mat0004395.html	ar5iv	text	# The influence of single magnetic impurities on the conductance of quantum microconstrictions. ## Abstract The nonlinear ballistic conductance of three-dimensional quantum microconstrictions, which contain magnetic impurities, is investigated. The nonlinear part of the conductance, which is due to the interaction of electrons with magnetic impurities is obtained. The analytical results have been analyzed numerically. It is shown that intensity of the Kondo anomaly in the conductance as the function of the applied voltage depends on the diameter of the constriction and the positions of impurities. The impurity-electron interaction in Kondo systems can be effectively studied by using the point contacts (PC). In first measurements of the differential PC resistance $`R(V)`$ in metals with magnetic impurities the zero-bias Kondo anomaly had been observed . These experiments were explained by quasiclassical theory of Kondo effect in PC’s . It was shown that in second-order Born approximation the magnetic impurity contribution to the PC resistance includes the logarithmic dependence $`R(V)\mathrm{ln}\left(V\right)`$ for $`eVT_K`$ and the saturation for $`eVT_K`$ ($`T_K`$ is the Kondo temperature, $`V`$ is the voltage applied to the PC). In accordance to the theory , the nonlinear correction to the ballistic PC resistance is proportional to the contact diameter. But in the experiments the size dependence of the PC current was not investigated due to the limited range of contact diameters, which were accessible. The development of the technique of mechanically controllable break junctions (MCBJ) has made it possible to create the stable PC’s, with the diameter adjustable over broad range, down to a single atom . In the MCBJ experiments authors had studied the resistance of ultrasmall contacts with magnetic impurities as function of the PC diameter $`d.`$ In the contrast to the prediction of the quasiclassical theory Yanson at el. observed that Kondo scattering contribution to the contact resistance is nearly independent from the contact diameter $`d`$ for small $`d`$. Such behavior authors had explained by the increasing of Kondo impurity scattering cross-section with decreasing of contact size. In theoretical works it was shown that in very small contacts the discreteness of impurity positions must be taken into account, and experiments may be explained by the ”classical” mesoscopic effect of the dependence of the point contact conductance on the spatial distribution of impurities. This effect is essential in the ”short” contacts and in the quasiclassical approximation it disappears with the increasing of the contact length. Zarand and Udvardi had considered the contact in the in the form of a long channel and suggested that the Kondo temperature $`T_K`$ is changed due to the strong the local density of states fluctuations generated by the reflections of conduction electrons at the surface of the contact. As a result of that, the effective cross section of electrons has the maximum, if the position of the impurity inside the contact corresponds to the maximum in the local electron density of states. But the mesoscopic effect of the spatial distribution of impurities in quantum contacts was not analyzed in the paper . In ultrasmall contacts the quantum phenomena, which are known as quantum size effect, occur. The effect of the $`2e^2/h`$ conductance quantization has been observed in experiments on contacts in the two-dimensional electron gas and in the ultrasmall three-dimensional constrictions, which is created by using the scanning tunnel microscopy and mechanically controllable break junctions . The defects produce the backscattering of electrons, and thus break the quantization of the conductance. From the other hand, the impurities situated inside the quantum microconstriction produce the nonlinear dependence of the conductance on the applied voltage . This dependence is the result of the interference of electron waves reflected by these defects . In this paper we present the theoretical solution of this problem for the conductance of a quantum microconstriction in the form of the long ballistic channel , which contains single magnetic impurities. The study is made of the first and second order corrections to the conductance of the ballistic microconstriction in the Born approximation. The effect of impurity positions is taken into account. Within the framework of the model of the long channel the quantum formula for the conductance $`G`$ is obtained. By using the model of the cylindrical microconstriction, the nonlinear conductance as a function of voltage $`V`$ and the width of constriction $`d`$ is analyzed numerically for different positions of a single impurity. Let us consider the quantum microconstriction in the form of a long and perfectly clean channel with smooth boundaries and a diameter $`d`$ comparable with the Fermi wave length $`\lambda _F=h/\sqrt{2m\epsilon _F},`$ where $`\epsilon _F`$ is the Fermi energy. We assume that this channel is smoothly (over Fermi length scale) connected with a bulk metal banks. As it was shown , in such constriction in the zeroth approximation on the adiabatic parameter $`\left\|d\right\|1`$ accurate quantization can be obtained. The corrections to the tunneling and reflection coefficients of electrons due to deviation from the adiabatic constriction are exponentially small, except near the points where the modes are switched on and off . When a voltage $`V`$ is applied to the constriction, a net current start to flow. In the limit $`V0,`$ the ballistic conductance of the quantum microconstriction is given by the formula $$G=\frac{dJ}{dV}=G_0\underset{\beta }{}f_F\left(\epsilon _\beta \right),$$ (1) where $`f_F`$ is the Fermi function, $`\epsilon _\beta `$ is the minimal energy of the transverse electron mode, $`\beta `$ is the full set of transverse discrete quantum numbers. The ballistic quantum PC displays the specific nonlinear properties, such as the conductance jumps $`e^2/h.`$ For the two dimensional PC these effects was considered in the papers . The aim of this study is to analyze the zero bias Kondo minimum in the PC conductance. We assume that the bias $`eV`$ is much smaller not only the Fermi energy $`\epsilon _F`$, but also the distances between the energies $`\epsilon _\beta `$ of quantum modes. In this case the effect of the influence of the applied bias to the transmission is negligibly small. Impurities and defects scatter the electrons that leads to the decreasing of the transmission probability. In accordance with the standard procedure the decreasing of the electrical current $`\mathrm{\Delta }I`$ due to the electron-impurity interaction connects with the velocity of the energy $`E`$ dissipation by the relation: $$\mathrm{\Delta }IV=\frac{dE}{dt}=\frac{dH_1}{dt};$$ (2) The Hamiltonian of the electrons $`H`$ contains the following terms:. $$H=H_0+H_1+H_{int},$$ (3) where $$H_0=\underset{k,\sigma }{}\epsilon _kc_{k\sigma }^{}c_{k\sigma }$$ (4) is Hamiltonian of free electrons, $$H_1=\frac{eV}{2}\underset{k,\sigma }{}sign(v_z)c_{k\sigma }^{}c_{k\sigma }$$ (5) describes the interaction of electrons with electric field. The Hamiltonian of the interaction of electrons with magnetic impurities $`H_{int}`$ can be written as $$H_{int}=\underset{j,k,k^{}}{}𝐉_{j,k,k^{^{}}}[S_z(c_k^{^{}}^{}c_kc_k^{^{}}^{}c_k)+S^{}c_k^{^{}}^{}c_k++S^+c_k^{^{}}^{}c_k].$$ (6) Here the operator $`c_{k\sigma }^+\left(c_{k\sigma }\right)`$ creates (annihilates) a conduction electron with spin $`\sigma ,`$ wave function $`\phi _k,`$ and energy $`\epsilon _k;`$ $`𝐒`$ denotes the spin of impurity; $`v_z`$ is the electron velocity along the channel; $`𝐉_{j,k,k^{^{}}}`$ is the matrix element of the exchange interaction of electron with impurity in the point $`𝐫_j`$; $`k\sigma `$ is the full set of quantum numbers; $$𝐉_{j,k,k^{}}=𝑑𝐫J(𝐫,𝐫_j)\phi _k(𝐫)\phi _k^{}^{}(𝐫).$$ (7) The electron wave functions and eigenvalues in the long channel in the adiabatic approximation are $`\phi _k(𝐫)`$ $`=\psi _\beta (𝐑)\mathrm{exp}\left({\displaystyle \frac{i}{\mathrm{}}}p_zz\right);`$ (8) $`\epsilon _k`$ $`=\epsilon _\beta +{\displaystyle \frac{p_z^2}{2m_e}};`$ (9) where $`k=(\beta ,p_z),`$ $`\beta `$ is the set of discrete transverse quantum numbers; $`p_z`$ is the momentum of an electron along the contact axis; $`m_e`$is an electron mass; $`𝐫=(𝐑,z),`$ $`𝐑`$ is a coordinate in the plain, perpendicular to the $`z`$ axis. Differentiating $`H_1`$ over the time $`t`$ we obtain the equation for the changing $`\mathrm{\Delta }I`$ of the current as a result of the interaction of electrons with magnetic impurities: $$V\mathrm{\Delta }I=\frac{1}{i\mathrm{}}[H_1\left(t\right),H_{int}\left(t\right)],$$ (10) where $$\mathrm{}=Tr\left(\rho \left(t\right)\mathrm{}\right).$$ (11) All operators are in the representation of interaction. The statistical operator $`\rho \left(t\right)`$ satisfies to equation $$i\mathrm{}\frac{\rho }{t}=[H_{int}\left(t\right),\rho \left(t\right)],$$ (12) which can be solved using the perturbation theory: $$\rho \left(t\right)=\rho _0+\frac{1}{i\mathrm{}}\underset{\mathrm{}}{\overset{t}{}}𝑑t^{}[H_{int}\left(t^{}\right),\rho _0]+\frac{1}{\left(i\mathrm{}\right)^2}\underset{\mathrm{}}{\overset{t}{}}𝑑t^{}\underset{\mathrm{}}{\overset{t^{}}{}}𝑑t^{\prime \prime }[H_{int}\left(t^{}\right),[H_{int}\left(t^{\prime \prime }\right),\rho _0]]\mathrm{}$$ (13) By means of Eq.(13) the changing in the electric current due to magnetic impurities can be determined $$\mathrm{\Delta }I=I_1+I_2+\mathrm{}=$$ (14) $$\frac{1}{\mathrm{}^2V}\underset{\mathrm{}}{\overset{t}{}}𝑑t^{}Tr\left(\rho _0[[H_1,H_{int}(t)],H_{int}(t^{})]\right)$$ (15) $$\frac{1}{i\mathrm{}^3V}\underset{\mathrm{}}{\overset{t^{}}{}}𝑑t^{\prime \prime }\underset{\mathrm{}}{\overset{t}{}}𝑑t^{}Tr\left(\rho _0[[[H_1,H_{int}(t)],H_{int}(t^{})],H_{int}(t^{\prime \prime })]\right)+\mathrm{}$$ (16) After the simple, but cumbersome calculations we find the first and second order corrections to the PC current $$I_1=\frac{e\pi }{\mathrm{}}s(s+1)\underset{n,m}{}\underset{i,j}{}(\mathrm{𝑠𝑖𝑔𝑛}v_{z_m}\mathrm{𝑠𝑖𝑔𝑛}v_{z_n})(f_mf_n)\delta \left(\epsilon _n\epsilon _m\right)𝐉_{j,n,m}𝐉_{i,m,n};$$ (17) $$I_2=\frac{\pi e}{\mathrm{}}s(s+1)\underset{n,m,k}{}\underset{i,j,l}{}(\mathrm{𝑠𝑖𝑔𝑛}v_{z_k}\mathrm{𝑠𝑖𝑔𝑛}v_{z_n})$$ (18) $$[\delta (\epsilon _n\epsilon _k)\mathrm{Pr}\frac{1}{\epsilon _m\epsilon _k}+\delta (\epsilon _m\epsilon _k)\mathrm{Pr}\frac{1}{\epsilon _n\epsilon _k}]$$ (19) $$[𝐉_{j,n,k}𝐉_{i,m,n}𝐉_{l,k,m}+𝐉_{j,k,n}𝐉_{i,n,m}𝐉_{l,m,k}]$$ (20) $$[2f_n(f_kf_m)+(f_mf_k)],$$ (21) where $`f_n=f_F\left(\epsilon _n+\frac{eV}{2}signv_z\right).`$ The first addition $`I_1`$ to the PC current describes the small spin-depended correction (of the order $`(J/\epsilon _F)^2`$ ) to the changing of the current due to the usual scattering. The second addition $`I_2`$ is also small too, but contains the Kondo logarithmic dependence on the voltage, and it is most important for the analysis of the nonlinear conductance of constrictions with magnetic impurities. The expressions (15) and (16) can be further simplified in the case of $`\delta `$potential of impurities $$J\left(𝐫\right)=J\delta \left(𝐫\right)$$ (22) In this case the addition $`I_2`$ to the ballistic current has the form: $$I_2=\frac{2J^3\pi e}{\mathrm{}}s(s+1)\underset{n,m,k}{}\underset{i,j,l}{}(\mathrm{𝑠𝑖𝑔𝑛}v_{z_k}\mathrm{𝑠𝑖𝑔𝑛}v_{z_n})$$ (23) $$[\delta (\epsilon _n\epsilon _k)\mathrm{Pr}\frac{1}{\epsilon _m\epsilon _k}+\delta (\epsilon _m\epsilon _k)\mathrm{Pr}\frac{1}{\epsilon _n\epsilon _k}]$$ (24) $$Re[\phi _k^{}(𝐫_j)\phi _n^{}(𝐫_i)\phi _m^{}(𝐫_l)\phi _k(𝐫_l)\phi _m(𝐫_i)\phi _n(𝐫_j)]$$ (25) $$[2f_n(f_kf_m)+(f_mf_k)],$$ (26) As it follows from the Eqs.(16), (18), the current $`I_2`$ depends from the positions of impurities. Two effects influence by value $`I_2:`$ the effect of quantum interference of scattered electron waves, which depends from the distances between impurities, and effect of the electron density of states in the points, where the impurities are situated. The nonlinear part of the conductance can be easy obtained after differentiation the Eq.18 over the voltage $`G_2=dI_2/dV.`$ In the case of a single impurity and at zero temperature $`T=0`$ this equation can be analytically integrated over momentum $`p_z`$ and the conductance $`G_2`$ takes the following form: $`G_2`$ $`=`$ $`{\displaystyle \frac{\pi e^2m_e^3}{\mathrm{}^4}}J^3s(s+1){\displaystyle \underset{\alpha ,\beta ,\gamma }{}}{\displaystyle \underset{\varkappa =\pm }{}}\left\|\psi _\alpha (𝐑)\|^2\right\|\psi _\beta (𝐑)\|^2\|\psi _\beta (𝐑)\|^2\left[p_\alpha ^{\left(\varkappa \right)}p_\beta ^{\left(\varkappa \right)}p_\gamma ^{\left(\varkappa \right)}\right]^1`$ (29) $`\left[\mathrm{ln}\right\|{\displaystyle \frac{p_\gamma ^{\left(\varkappa \right)}p_\gamma ^{\left(\varkappa \right)}}{p_\gamma ^{\left(\varkappa \right)}+p_\gamma ^{\left(\varkappa \right)}}}\left({\displaystyle \frac{p_\alpha ^{\left(\varkappa \right)}}{p_\gamma ^{\left(\varkappa \right)}}}\right)\|+(1\delta _{\alpha \beta })\mathrm{ln}\|{\displaystyle \frac{p_\alpha ^{\left(\varkappa \right)}p_\beta ^{\left(\varkappa \right)}p_\alpha ^{\left(\varkappa \right)}p_\beta ^{\left(\varkappa \right)}}{p_\alpha ^{\left(\varkappa \right)}p_\beta ^{\left(\varkappa \right)}+p_\alpha ^{\left(\varkappa \right)}p_\beta ^{\left(\varkappa \right)}}}\|+`$ $`\delta _{\alpha \beta }\mathrm{ln}\left\|{\displaystyle \frac{\left(p_\alpha ^{\left(\varkappa \right)}\right)^2\left(p_\alpha ^{\left(\varkappa \right)}\right)^2}{\left(p_\alpha ^{\left(\varkappa \right)}\right)^2}}\right\|];`$ where $$p_\alpha ^{(\pm )}=\sqrt{2m_e\left(\epsilon _F\pm \frac{eV}{2}\epsilon _\alpha \right)},$$ (30) and the transverse parts of the wavefunction $`\psi _\alpha (𝐑)`$ and the electron energy $`\epsilon _\alpha `$ are defined by Eqs. (8), (9). Carrying out the numerical calculations we use the free electron model of a point contact consisting of two infinite half-spaces connected by a long ballistic cylinder of a radius $`R`$ and a length $`L`$ (Fig.1). In a limit $`L\mathrm{}`$ the electron wave functions $`\phi _k\left(𝐫\right)`$ and eigenstates $`\epsilon _k`$ can be written as $$\phi _k\left(𝐫\right)=\frac{1}{\sqrt{\mathrm{\Omega }}J_{m+1}\left(\gamma _{mn}\right)}J_m\left(\gamma _{mn}\frac{\rho }{R}\right)\mathrm{exp}\left(im\phi +\frac{i}{\mathrm{}}p_zz\right);$$ (31) $$\epsilon _k=\epsilon _{mn}+\frac{p_z^2}{2m_e};\epsilon _{mn}=\frac{\mathrm{}^2}{2m_eR^2}\gamma _{mn}^2$$ (32) and cylindrical coordinates $`𝐫=(\rho ,\phi ,z)`$ with the axis $`z`$ along the channel axis have been used. Here $`k=(n,m,p_z)`$ are the quantum numbers, $`\mathrm{\Omega }=\pi R^2L`$ is the volume of the channel, $`\gamma _{mn}`$ are the n-th zero of the Bessel function $`J_m.`$ Because the degeneration of the electron energy on azimuthal quantum number $`m`$ ( as a result of the symmetry of the model), quantum modes with $`\pm m`$ give the same contribution to the conductance. In this model the ballistic conductance (1) has not only steps $`G_0,`$ but also steps $`2G_0`$ . In Fig.2 the dependence of the nonlinear conductance on the applied bias is shown for the different positions of a single magnetic impurity inside the channel. The results obtained confirm that the nonlinear effect is strongly depend on the position of impurity. If the impurity is situated near the surface of the constriction $`𝐫=𝐑`$, where the square module of the electron wave function is small, its influence to the conductivity is negligible. This conclusion is confirmed by the calculations of the dependence $`G_2`$ on the position of the impurity for different number of quantum modes (Fig.3). Results indicate that the mesoscopic effect of the impurity position is more essential for ultrasmall contacts, which contain only few conducting modes, and $`G_2`$ has a maximum. The similar results is obtained for the dependence of $`G_2`$ on the radius $`R`$ of the constriction (Figs.4,5). In the single-mode constriction (Fig.4) the conductance $`G_2`$ displays much more stronger dependence on $`R,`$ than in the contact with five conducting modes (Fig.5). Thus, we have shown that in the long quantum microconstrictions the spatial distribution of magnetic impurities influences to the nonlinear dependence of the conductance on the applied voltage. This mesoscopic effect is due to the strong dependence the amplitude of an electron scattering on the positions of impurities. As a result of the reflection from the boundaries of the constriction the electron wave functions, which correspond to the finite electron motion in the transverse to the contact axis direction, are the standing waves. If the impurity is situated near the point, in which the electron wave function is equal to zero (near the surface of the constriction or, for quantum modes with numbers $`n>1,`$ in some points inside), its scattering of electrons is small. The fact the amplitude of the Kondo minimum of the conductance of the quantum contact display the mesoscopic effect of the dependence on the positions of single impurities. This effect is most important in the case, when only few quantum modes are responsible on the conductivity of the constriction. Figure captions. Fig. 1. Schematic representation of a ballistic microconstriction in the form of a long channel, adiabatically connected to large metallic reservoirs. Magnetic impurities inside the constriction are shown. Fig. 2. The voltage dependence of the nonlinear part of conductance $`G_2`$ (19) from the distance of the impurity from the contact axis ($`2\pi R=5.2\lambda _F;`$ $`T=0;`$ 1 - $`2\pi \rho =1.5\lambda _F`$ ; 2 -$`2\pi \rho =2.5\lambda _F`$ ; 3 -$`2\pi \rho =3.0\lambda _F`$ ; 4 - $`2\pi \rho =3.5\lambda _F`$ ). Fig. 3. The dependence of $`G_2`$ (19) on the position of the impurity for the different quantum modes in the constriction ($`V=0.02\epsilon _F;`$ $`T=0;`$ 1 - one mode ($`2\pi R=3\lambda _F`$); 2 - three modes($`2\pi R=4\lambda _F`$); 3 - five modes ($`2\pi R=5.3\lambda _F`$); 4 - six modes ($`2\pi R=6\lambda _F`$)). Fig. 4. The dependence of $`G_2`$ (19) on the radius of the constriction for the single mode channel and different positions of the impurity ($`V=0.02\epsilon _F;T=0;`$1 - $`2\pi \rho =0.5\lambda _F`$ ; 2 -$`2\pi \rho =1.0\lambda _F`$ ; 3 -$`2\pi \rho =1.5\lambda _F`$ ; 4 - $`2\pi \rho =2.0\lambda _F`$ ) Fig. 5. The dependence of $`G_2`$ on the radius for the microconstriction with five quantum modes and different positions of the impurity ($`V=0.02\epsilon _F;T=0;`$1 - $`2\pi \rho =0.5\lambda _F`$ ; 2 -$`2\pi \rho =1.5\lambda _F`$ ; 3 -$`2\pi \rho =2.5\lambda _F`$ ; 4 - $`2\pi \rho =4.5\lambda _F`$ )
warning/0004/nlin0004004.html	ar5iv	text	# A method to find unstable periodic orbits for the diamagnetic Kepler problem ## I Introduction Unstable periodic orbits (UPOs) as the skeleton of chaotic systems provide a powerful tool to analyze the property of low-dimensional dynamics. As a simple physical system, the diamagnetic Kepler problem (DKP) displays important roles of UPOs in the classical and quantum chaos. For the chaotic system, the UPOs are filled in the phase space and their total number roughly increases with the periodic length. Many UPOs can be found by searching along three immediate symmetry lines in the configuration space. In view of the scattering of a particle in the four-disk system, symbolic dynamics has been established and a method to find UPOs by contracting a rectangle in terms of the symbolic dynamics has been given. By considering the stretching and wrapping in the lifted space, symbolic dynamics has been established. In this paper, we will present a method to find UPOs based on the ordering of stable and unstable manifolds and describe a one to one correspondence between the UPOs and their corresponding symbolic sequences for the DKP. The paper is organized as follows. In Sec. II, we describe the method to locate initial values of UPOs in terms of symbolic dynamics. In Sec. III, by investigating the UPOs up to length 6, a one to one correspondence between the UPOs and their corresponding symbolic sequences is shown under the system symmetry decomposition. Finally, in Sec. IV, some conclusions are given. ## II Method to locate the initial values of UPOs The classical dynamics of a hydrogen atom with zero angular momentum in a uniform magnetic field $`B`$ along the $`z`$-axis is described by the “pseudo” Hamiltonian (using the semi-parabolic coordinates): $$h=\frac{p_\mu ^2}{2}+\frac{p_\nu ^2}{2}ϵ(\mu ^2+\nu ^2)+\frac{1}{8}\mu ^2\nu ^2(\mu ^2+\nu ^2)2,$$ (1) where $`ϵ=E\gamma ^{2/3}`$ is the scaled energy. $`h`$ has $`C_{4v}`$ and time-reversal symmetry. In view of the Birkhoff canonical coordinates, a Poincaré section has been chosen along a counterclockwise contour. We choose the $`\mu `$ and $`\nu `$ coordinate axes as a Poincaré section and determine a Poincaré map in terms of the entering direction of a orbit. In the first quadrant, the arc coordinate of the map is defined as $`s=\mu `$ along the $`\mu `$ axis and $`s=\nu `$ along the $`\nu `$ axis. Using the available transformation $`\frac{s}{1+\|s\|}s`$, we get $`s[1,1)`$. The second coordinate of the map is $`v=\frac{p_\mu }{\sqrt{p_\mu ^2+p_\nu ^2}}`$ along the $`\mu `$ axis and $`v=\frac{p_\nu }{\sqrt{p_\mu ^2+p_\nu ^2}}`$ along the $`\nu `$ axis. According to the $`C_4`$ symmetry, the Poincaré maps in the second, third and fourth quadrants can be obtained by rotating the map in the first quadrant with angles $`\pi /2`$, $`\pi `$ and $`3\pi /2`$, respectively, around the center anti-clockwisely and adding 2,4 and 6 to $`s`$. Thus, the dynamics on the Poincaré surface is represented by the map on the annulus $`s[1,7)`$ and $`v[1,1]`$. The reduced domain ($`s[0,2)`$, $`v[1,1]`$) is constructed by considering the rotational symmetry. In the same way, the minimal domain (MD) ($`s[0,1)`$, $`v[1,1]`$) is taken by discarding the $`\pi `$-rotation around the origin. To illustrate the method to find UPOs, we first discuss the ordering of stable and unstable manifolds in the Poincaré section. In Fig. 1, the stable and unstable manifolds in the MD display the local and global constricting, stretching and folding processes of the dynamics. When the tangencies associated with folding of one family of the manifolds occur, a partition line through them need to be introduced. Using the method given in Ref. , we roughly determine the partition line $`C_0`$ in Fig. 1. The demarcation line between two regions with the different rotation numbers in the lifted space is marked by $`B_0`$. In terms of the natural ordering in the lifted space and the tangencies of manifolds, we obtain the region partition in the MD with the symbols ($`L_0`$, $`R_0`$ and $`R_1`$) and the symbolic ordering of single letter for forward and backward sequences $$L_0<R_0<R_1,L_0<R_0<R_1.$$ (2) Except the ordering, the parity of a leading string in the sequences is defined by the oddness of the total number of $`R_0`$. Only odd leading string reverses the ordering. After the region partition and symbolic ordering are introduced, the families of the stable and unstable manifolds make a curvilinear coordinates system in the MD. All points on a stable (unstable) manifold converge to each other when iterated forward (backward). In terms of the language of symbolic dynamics, the points in each stable (unstable) manifold have the same forward (backward) sequence, i.e., the same ordering. A family of stable (unstable) manifolds is ordered according to their transverse intersections with an unstable (stable) manifold. The ordering is well-defined as long as there are no tangencies between the two manifolds. The occurrence of tangencies changes the ordering of folded stable or unstable manifolds, i.e., two families of sub-manifolds divided by the partition line have the reverse ordering. In Fig. 1, the ordering of stable (unstable) manifolds increases from the left-bottom (left-top) to the right-top (right-bottom) along each unstable (stable) manifold. For every point in the MD, there exists a stable or unstable manifold, which passes through the point and joins with the boundary line $`D_0`$. The ordering of stable and unstable manifolds in the Poincaré section can be displayed by using $`\alpha `$ and $`\beta `$ in the symbolic plane. Every forward sequence $`s_0s_1\mathrm{}s_n\mathrm{}`$ may correspond to a number $`\alpha `$, represented in base 3, between 0 and 1. The correspondence of symbol $`s_i`$ to number $`\xi _i\{0,1,2\}`$ is given as $`L_00`$, $`R_01`$, $`R_12`$ if the leading string $`s_0s_1\mathrm{}s_{i1}`$ is even, and as $`\{L_0,R_0,R_1\}\{2,1,0\}`$ for the odd case. Here, the number $`\alpha `$ is defined as $$\alpha =\underset{i=0}{\overset{\mathrm{}}{}}\xi _i3^{(i+1)}.$$ (3) Similarly, the number $`\beta `$ for a backward sequence $`\mathrm{}s_m`$ $`\mathrm{}s_2`$ $`s_1`$ is defined by $$\beta =\underset{j=1}{\overset{\mathrm{}}{}}\eta _j3^j,$$ (4) where $`\eta _j\{0,1,2\}`$ is determined by $`s_j`$ and the oddness of the leading string $`s_{j+1}\mathrm{}s_2s_1`$. In this way, forward (backward) sequences are ordered according to their $`\alpha `$ ($`\beta `$)-values. In particular, at the boundary line $`s=0^+`$, $`\alpha `$ ($`\beta `$) increases (decreases) monotonically from $`v=1`$ to $`v=1`$. In Fig. 2, the symbolic plane has a one to one correspondence with the MD. Many point sets are separated by the boundary line $`D_0`$, the partition line $`C_0`$ and their forward or backward maps. For a given symbolic sequence, its $`\alpha _{exact}`$ and $`\beta _{exact}`$ values may be calculated by using (3) and (4). In the moving process of a starting point, $`\alpha _{approx}`$ and $`\beta _{approx}`$ of the current point are respectively determined by using its forward and backward sequences formed in 50 forward and backward maps of the current point in the MD. Our method is described as follows: (i) We run near $`\alpha _{exact}`$ from $`v`$=-1 in the boundary line $`s=0^+`$ and then iterate $`v`$ to $`\alpha _{exact}`$ by using the half-division method. In the iteration process, if both $`\alpha _{approx}`$ and $`\beta _{approx}`$ for one of the two points are close to $`\alpha _{exact}`$ and $`\beta _{exact}`$ with a given error ($`10^6`$), we obtain an initial value of the UPO corresponding to the given symbolic sequence. The UPO is that passing through the origin. (ii) If $`\|\alpha _{approx}\alpha _{exact}\|<10^8`$, i.e., the difference between the values $`\alpha _{approx}`$ of two points surrounding $`\alpha _{exact}`$ is less than the error, we may make a stable manifold from the middle one between the two points. Along the stable manifold, $`\alpha _{approx}`$ remains approximately unchanged and $`\beta _{approx}`$ increases monotonically. When $`\beta _{approx}`$ is near $`\beta _{exact}`$, i.e., the current point leaps from one side of $`\beta _{exact}`$ to another, one should decrease the step size to achieve $`\|\beta _{approx}\beta _{exact}\|<10^8`$. However, in this case, the requirement $`\|\alpha _{approx}\alpha _{exact}\|<10^8`$ is not met. We may make an unstable manifold from the middle point and choose a direction according to $`\alpha _{approx}`$. Along the unstable manifold, $`\beta _{approx}`$ remains approximately unchanged and $`\alpha _{approx}`$ is monotonically changed. At this time, if the change of $`\beta _{approx}`$ or $`\alpha _{approx}`$ decreases monotonically, the above process will be continued until both the requirements $`\|\alpha _{approx}\alpha _{exact}\|<10^{15}`$ and $`\|\beta _{approx}\beta _{exact}\|<10^{15}`$ are met. Thus, in terms of the differences $`\|\alpha _{approx}\alpha _{exact}\|`$ and $`\|\beta _{approx}\beta _{exact}\|`$, we may extract an initial value of the UPO corresponding to the symbolic sequence, which is allowed. The UPO is that non-passing through the origin. If the change of $`\beta _{approx}`$ or $`\alpha _{approx}`$ does not decrease monotonically in the 10 circulations, i.e., the algorithm does not converge, the above process will be stopped. Thus, we give up the process for $`\alpha _{approx}`$ and repeat (ii) for $`\beta _{approx}`$. (iii) If $`\|\alpha _{approx}\alpha _{exact}\|>10^8`$ even when the step of the half-division method is enough small, i.e., $`\alpha _{approx}`$ jumps between two nearly invariant values around $`\alpha _{exact}`$, we give up the process (ii) for $`\alpha _{approx}`$ and repeat it for $`\beta _{approx}`$. (iv) If $`\|\beta _{approx}\beta _{exact}\|>10^8`$ even when the step of the half-division method is enough small, i.e., $`\beta _{approx}`$ also jumps between two nearly invariant values around $`\beta _{exact}`$, we may say the symbolic sequence is possibly forbidden. Using the above method for the given symbolic sequence, we can obtain an initial point of its corresponding UPO or a possible estimation of forbiddenness. For the latter case, we may further repeat the above steps to determine the admissibility of symbolic sequences produced by shifting the given sequence in its length. We stop the process as soon as the symbolic sequence is allowed. If all symbolic sequences generated by shifting the given sequence in its length are possibly forbidden, we may say the given symbolic sequence is forbidden. Finally, we can determine the allowance or forbiddenness of the given symbolic sequence and extract an initial point of the UPO for the allowed symbolic sequence. To illustrate the main running process of the above method, we take the symbolic sequences $`R_0^2R_1`$ and $`L_0R_0R_1`$ as two examples. For the sequence $`(R_0^2R_1)^{\mathrm{}}(R_0^2R_1)^{\mathrm{}}`$, using (3) and (4), we get $`\alpha _{exact}`$=0.5384615384615387 and $`\beta _{exact}`$=0.8461538461538460. The entire process to find the UPO is displayed in Fig. 1 and Fig. 2. First, we begin from the point ($`s_0`$,$`v_0`$)=($`10^{10}`$,-0.99999) in Fig. 1 and take its symbolic sequence to calculate ($`\alpha `$,$`\beta `$) in Fig. 2. We run along the line $`s=s_0`$ with a step $`\mathrm{\Delta }v_0`$ and determine the last point near $`\alpha _{exact}`$ by comparing $`\alpha _{approx}`$ with $`\alpha _{exact}`$. Then, we use the half-division method to decrease the difference between $`\alpha _{approx}`$ and $`\alpha _{exact}`$. When the difference is less than $`10^8`$, we draw the stable manifold from the point in Fig. 1. Along the stable manifold, we compare $`\beta _{approx}`$ with $`\beta _{exact}`$. When the difference between them is less than $`10^8`$, we compare $`\alpha _{approx}`$ with $`\alpha _{exact}`$ and get their difference 0.0006378404570228. The last point is marked by $`a`$ in Fig. 3(a). According to the comparison, we draw the unstable manifold from the point. When $`\|\alpha _{approx}\alpha _{exact}\|<10^{15}`$, we compare $`\beta _{approx}`$ with $`\beta _{exact}`$ and get their difference 0.0000029487355183. The last point is marked by $`b`$ in Fig. 3(b). Furthermore, we continue the above steps with the given error $`10^{15}`$. In the second circulation, when $`\|\beta _{approx}\beta _{exact}\|<10^{15}`$, the difference $`\|\alpha _{approx}\alpha _{exact}\|`$ is 0.0000000000441224. The last point is marked by $`c`$ in Fig. 3(c); when $`\|\alpha _{approx}\alpha _{exact}\|<10^{15}`$, the difference $`\|\beta _{approx}\beta _{exact}\|`$ is 0.0000000000000114. The last point is marked by $`d`$ in Fig. 3(d). In the third circulation, when $`\|\beta _{approx}\beta _{exact}\|<10^{15}`$, the difference $`\|\alpha _{approx}\alpha _{exact}\|`$ is also less than $`10^{15}`$. Thus, we get the point $`(s,v)`$=(0.6872027097581007,-0.3670108804483804) marked by $`e`$ in Fig. 3(d) corresponding to the sequence $`(R_0^2R_1)^{\mathrm{}}(R_0^2R_1)^{\mathrm{}}`$. Also, we may obtain its forward and backward sequences with 50 letters to conform the result. Finally, using Newton method, we exactly extract the UPO. For the sequence $`(L_0R_0R_1)^{\mathrm{}}(L_0R_0R_1)^{\mathrm{}}`$, we get $`\alpha _{exact}`$=0.1428571428571428 and $`\beta _{exact}`$=0.8571428571428572. First, beginning from the point ($`s_0`$,$`v_0`$)=($`10^{10}`$,-0.99999), we run along the line $`s=s_0`$ with a step $`\mathrm{\Delta }v_0`$ and determine the last point near $`\alpha _{exact}`$ by comparing $`\alpha _{approx}`$ with $`\alpha _{exact}`$. Then, using the half-division method, we decrease the difference between $`\alpha _{approx}`$ and $`\alpha _{exact}`$. When $`\mathrm{\Delta }v_0<10^{15}`$, the $`\alpha _{approx}`$ values of the last and current points are 0.1428571417083783 and 0.1428571738737871, respectively. Since $`\mathrm{\Delta }\alpha _{approx}>10^8`$, we give up the process for $`\alpha _{approx}`$ and repeat it for $`\beta _{approx}`$. When $`\mathrm{\Delta }v_0<10^{15}`$, the $`\beta _{approx}`$ values of the last and current points are 0.8553113553113367 and 0.9046940713607194, respectively. Since $`\mathrm{\Delta }\alpha _{approx}>10^8`$, the process for $`\beta _{approx}`$ is stopped. Thus, the sequence $`L_0R_0R_1`$ is possibly forbidden. Shifting the sequence $`L_0R_0R_1`$, we get $`R_0R_1L_0`$ and extract $`\alpha _{exact}`$=0.4285714285714285 and $`\beta _{exact}`$=0.2857142857142856. We repeat the above process. When $`\mathrm{\Delta }v_0<10^{15}`$ for $`\alpha _{approx}`$, $`\mathrm{\Delta }\alpha _{approx}=0.0000000964962270>10^8`$. When $`\mathrm{\Delta }v_0<10^{15}`$ for $`\beta _{approx}`$, $`\mathrm{\Delta }\beta _{approx}=0.0164609053497942>10^8`$. The sequence $`R_0R_1L_0`$ is possibly forbidden. Shifting the sequence $`R_0R_1L_0`$, we get $`R_1L_0R_0`$ and extract $`\alpha _{exact}`$=0.7142857142857142 and $`\beta _{exact}`$=0.5714285714285717. We repeat the above process. When $`\mathrm{\Delta }v_0<10^{15}`$ for $`\alpha _{approx}`$, $`\mathrm{\Delta }\alpha _{approx}=0.1481481481481481>10^8`$. When $`\mathrm{\Delta }v_0<10^{15}`$ for $`\beta _{approx}`$, $`\mathrm{\Delta }\beta _{approx}=0.0054869684499316>10^8`$. The sequence $`R_1L_0R_0`$ is possibly forbidden. Since the three sequences generated by shifting $`L_0R_0R_1`$ are possibly forbidden, the sequence $`L_0R_0R_1`$ is forbidden. ## III Correspondence between UPOs and Symbolic Sequences Before using the method to find UPOs, we discuss qualitatively the bottom-left region forbidden by the boundary line $`D_0`$ in Fig. 2. Let us focus on the orbits near the origin in the configuration space in Fig. 4. Since the orbits are nearly linear, we can obtain the relation between two continuous Poincaré maps. A current point near the origin is described as ($`s_0`$, $`v_0`$), its forward and backward mapping points can be written as ($`s_1`$,$`v_1`$) and ($`s_1`$,$`v_1`$), respectively. (i) If $`v_0(1,0)`$, an orbit between the points ($`s_0`$,$`v_0`$) and ($`s_1`$,$`v_1`$) is nearly linear and has not turning points. We have $$\begin{array}{c}s_1=\frac{s_0\sqrt{1v_0^2}}{v_0(s_01)+s_0\sqrt{1v_0^2}}s_0\sqrt{1v_0^2}/v_0>0,\hfill \\ v_1=\sqrt{1v_0^2},\hfill \end{array}$$ (5) where $`v_1(0,1)`$. We can obtain the symbolic description $`L_0`$ for the point ($`s_0`$,$`v_0`$) and $`R_0`$ or $`R_1`$ for the point ($`s_1`$,$`v_1`$) in Fig. 1. The symbolic string ($`L_0^2`$) is forbidden due to the geometry of the orbits. When the coordinate $`v_0`$ changes from -1 to 0, the symbolic sequence for the point ($`s_0,v_0`$) changes from $`R_1^{\mathrm{}}L_0R_0R_1^{\mathrm{}}`$ to $`R_1^{\mathrm{}}R_0L_0R_1^{\mathrm{}}`$ and $`\alpha `$ ($`\beta `$) increases (decreases) monotonically from $`\frac{1}{9}`$ (1) to $`\frac{1}{3}`$ ($`\frac{1}{3}`$). Thus, we get the left half part of the boundary line $`D_0`$ in Fig. 2. (ii) If $`v_0(0,1)`$, an orbit between the points ($`s_1`$,$`v_1`$) and ($`s_0`$,$`v_0`$) is nearly linear and has not turning points. We have $$\begin{array}{c}s_1=\frac{s_0\sqrt{1v_0^2}}{v_0(1s_0)+s_0\sqrt{1v_0^2}}s_0\sqrt{1v_0^2}/v_0>0,\hfill \\ v_1=\sqrt{1v_0^2},\hfill \end{array}$$ (6) where $`v_1(1,0)`$. We can obtain the symbolic description $`R_0`$ or $`R_1`$ for the point ($`s_0`$, $`v_0`$) and $`L_0`$ for the point ($`s_1`$, $`v_1`$) in Fig. 1. The symbolic string ($`L_0^2`$) is forbidden due to the geometry of orbits. When the coordinate $`v_0`$ changes from 0 to 1, the symbolic sequence changes from $`R_1^{\mathrm{}}L_0R_0R_1^{\mathrm{}}`$ to $`R_1^{\mathrm{}}R_0L_0R_1^{\mathrm{}}`$ and $`\alpha `$ ($`\beta `$) increases (decreases) monotonically from $`\frac{1}{3}`$ ($`\frac{1}{3}`$) to 1 ($`\frac{1}{9}`$). Thus, we get the right half part of the boundary line $`D_0`$ in Fig. 2. Therefore, combining (i) and (ii), the bottom-left region forbidden by the boundary line $`D_0`$ is obtained due to the linearity of the orbits between two continuous maps near the origin. In case of three symbols with the above exclusion rule, the associated connectivity matrix is $$\begin{array}{ccc}& & \begin{array}{ccc}L_0& R_0& R_1\end{array}\\ T=& \begin{array}{c}L_0\\ R_0\\ R_1\end{array}& \left(\begin{array}{ccc}0& 1& 1\\ 1& 1& 1\\ 1& 1& 1\end{array}\right)\end{array}.$$ The total number of symbolic sequences of the length $`n`$ is $`N_n=trT^n`$, which can be decomposed into the number $`M_d`$ of primary symbolic sequences of length $`d`$ dividing $`n`$: $`N_n=_{d\|n}dM_d.`$ By the Möbius inversion, we get $$M_n=\frac{1}{n}\underset{d\|n}{}\mu (\frac{n}{d})N_d,$$ where $`\mu (1)=1`$, $`\mu (n)=0`$ if $`n`$ contains the square of a prime and $`\mu (n)=(1)^k`$ if $`n`$ contains $`k`$ prime factors. All the calculated results are shown in Table I. In Ref. we have shown that the symbolic dynamics in the MD removes the rotational and reflectional symmetries and only preserves the time-reversal symmetry $`T`$, which just reverses the original sequence. We make a sequence set with different lengths by combining the three symbols. From the set, the symbolic sequences generated by shifting each primary sequence are removed. If a symbolic sequence generated by the time-reversal symmetry is only shifts of the primary symbolic sequence, it is also abandoned. Finally, all symbolic sequences including $`L_0^2`$ should be removed from the set. Thus, the number of the sequence set is equal to $`M_n`$ obtained from the connectivity matrix, which is shown in Table I. Under the system symmetry decomposition, only the time-reversal symmetry is preserved in the sequence set, which can be used in cycle expansion. Using the method mentioned above, we determine the admissibility of a sequence set up to length 8 and extract approximate initial points of the UPOs corresponding to the allowed symbolic sequences for $`ϵ=0`$. The number of allowed symbolic sequences is presented in Table I. In the primary sequences, their admissibility as physical UPOs can be further determined in terms of pruning fronts corresponding to the tangencies between stable and unstable manifolds and bound line. For example, in Fig. 5, the points ($`\alpha _{R_0^2R_1}`$, $`\beta _{R_0^2R_1}`$) and ($`\alpha _{R_0R_1R_0}`$, $`\beta _{R_0R_1R_0}`$) exist in the allowed $`R_0`$ region, and the point ($`\alpha _{R_1R_0^2}`$, $`\beta _{R_1R_0^2}`$) exists in the allowed $`R_1`$ region. So, the sequence $`(R_0^2R_1)^{\mathrm{}}(R_0^2R_1)^{\mathrm{}}`$ is allowed. The point ($`\alpha _{L_0R_0R_1}`$, $`\beta _{L_0R_0R_1}`$) does not exist in the allowed $`L_0`$ or $`R_0`$ or $`R_1`$ region, but in the forbidden region of boundary line $`D_0`$. So, the sequence $`(L_0R_0R_1)^{\mathrm{}}(L_0R_0R_1)^{\mathrm{}}`$ is forbidden. We use the method to determine the admissibility of primary sequences up to length 8. The result is consistent with the above one. For a real number, our Sun-workstation has significant digits with the 16’s space. Due to the restriction of precision, the exact UPO up to length 6 can be further exacted by using the Newton method. In order to check a one to one correspondence between the UPOs and their corresponding symbolic sequences, we display the allowed symbolic sequences up to period 6 in Table II and their some corresponding UPOs in Fig. 6. In terms of the rotational and reflectional symmetries, the Orbit Period is the same or twice or quadruple the Sequence Period, which can be determined in the configuration space. In the 38 UPOs of Table II, except the pairs of UPOs (5) and (6), (11) and (15), (14) and (16), (25) and (29), (26) and (32), (27) and (33) shown in Fig. 6, other UPOs have the different configurations and cannot be produced from one to another in terms of the time reversal. For each one of the 26 UPOs, the time-reversal orbit is the same as that produced by the rotational symmetry, or reflectional symmetry, or their combination. So, the time-reversal symmetry $`T`$ is degenerate, i.e., the UPO and its time reversal orbit correspond to the same symbolic sequences in Table II. For the pairs of UPOs (11) and (15), (14) and (16), (26) and (32), (27) and (33), they have also the different configurations, but the two UPOs of each pair have the same one. For each one of the two orbits, its time-reversal orbit is another one, which cannot be produced from the given UPO by the rotational symmetry, or reflectional symmetry, or their combination. So, the time-reversal symmetry $`T`$ is non-degenerate. The two orbits with the non-degenerate $`T`$ symmetry correspond to the different symbolic sequences in Table II. For the pairs of UPOs (5) and (6), (25) and (29) passing through the origin, they have also the different configurations, but the two UPOs of each pair have the same one even when the time reversal symmetry is considered. When the right limit $`s=0^+`$ of the UPOs is taken into account in the configuration space, the intersection points of the UPOs (5) and (29), (6) and (25) should be removed from the origin into the third and first quadrants, respectively. At this time, for the pair of UPOs (5) and (6) with the right limit, the two lines of the UPO in the center should be removed to pass through the axes $`\nu >0`$ and $`\nu <0`$, respectively; for the pairing UPOs (25) and (29) with the right limit, the two lines of the UPO in the center should be removed to pass through the axis $`\nu >0`$. The pairs of UPOs (5) and (6), (25) and (29) with the right limit have the similar configurations to the above UPOs no-passing through the origin and the same results can be concluded. Thus, a one to one correspondence between the UPOs and their corresponding symbolic sequences is shown under the system symmetry decomposition. ## IV Conclusion According to the ordering of stable and unstable manifolds, we propose a method to determine the admissibility of symbolic sequences and to find UPOs corresponding to allowed symbolic sequences for the diamagnetic Kepler problem. By searching the UPOs up to length 6, a one to one correspondence between the UPOs and their corresponding symbolic sequences has been shown under the system symmetry decomposition. ## ACKNOWLEDGMENTS This work was supported in part by Post-Doctoral Foundation and Nonlinear Science Project of China (ZBW) and National Natural Scientific Foundation of China (JYZ). FIGURE CAPTION
warning/0004/hep-th0004163.html	ar5iv	text	# 1 Introduction ## 1 Introduction The causal approach to renormalization theory of by Epstein and Glaser , leads to important simplification of the renormalization theory as well as of the computational aspects. This approach works for quantum electrodynamics , Yang-Mills theories , , -, , gravitation , , , etc. In this paper we investigate the rôle of dilation invariance in the causal approach. The pioneering works on scale covariance in perturbative field theory are , and . A mathematical refined analysis was developed in and , the main mathematical tool being the so-called quantum action principle (for a review see ). Our strategy will be based exclusively on the Epstein-Glaser construction of the chonological producs for the free fields. In the next Section we define the dilation invariance operator for various free fields. Next, we remind the basic facts about renormalization theory. We will emphasize the original Epstein-Glaser approach where one considers a set of (linearly independent) interaction Lagrangian and attaches to each of this Lagrangian a (space-time dependent) coupling constant. Then we are able to prove the basic theorem concerning the arbitrariness of the chronological products for the same set of interaction Lagrangian. This problem was already addressed in , , but we argue that the natural framework is the multi-valued coupling constant approach of . In Section 4 we obtain consequences about the scale behaviour of the chronological products. One expects that the action of the dilations operators on the chronological product should give the usual scaled chronological products, up to some scale anomalies. We will prove that one can reduce the analysis to a cohomological problem. An important difference appears in the study of this problem in the cases of a massless and of a massive field. If there exists at least one massive field in the theory then we can prove that one can choose the chronological product to be scale covariant. In the opposite case, one finds out that some anoumalous terms of logarithmic behaviour can appear. We emphasis again that these results hold for the chronological products of the free fields. We will comment on what one should expect for the case of the interacting fields in the last Section. We will apply these considerations for Yang-Mills models in Section 5 and obtain a restrictions on the possible form of the anomalies, namely the canonical dimension of such an anomalous expression must be $`5`$. ## 2 Dilation Invariance in Quantum Field Theory It is well known that the Fock space of the real scalar field of mass $`m`$ can be defined as: $$_m_{n=0}^{\mathrm{}}_m^{(n)}$$ (2.0.1) where $`_m^{(n)}`$ is the set of Borel function $`\mathrm{\Phi }^{(n)}:(X_m^+)^n`$ which are square integrable with respect to the Lorentz invariant measure: $`d\alpha _m^+(p)\frac{d𝐩}{\sqrt{𝐩^2+m^2}}`$ and completely symmetric in the all variables (see for notations). Then we have: ###### Proposition 2.1 Let us define for any $`\lambda _+`$ the operators $`𝒰_\lambda :_m_{\lambda ^1m}`$ as follows: $$\left(𝒰_\lambda \mathrm{\Phi }^{(n)}\right)(p_1,\mathrm{},p_n)=\lambda ^n\mathrm{\Phi }^{(n)}(\lambda p_1,\mathrm{},\lambda p_n).$$ (2.0.2) Then: (i) The operators $`𝒰_\lambda `$ are unitary; (ii) The following relations are verified for all $`\lambda ,\lambda ^{}_+`$: $$𝒰_\lambda 𝒰_\lambda ^{}=𝒰_{\lambda \lambda ^{}};$$ (2.0.3) (iii) If $`𝒰_{a,L}^{[m]}`$ is the representation of the Poincaré group in the Fock space $`_m^{(n)}`$, then: $$𝒰_{a,L}^{[m]}𝒰_\lambda =𝒰_\lambda 𝒰_{\lambda ^1a,L}^{[\lambda m]}$$ (2.0.4) for all translations $`a`$ and all Lorentz transformations $`L`$. Proof: The proof of the first assertion is based on the scaling properties of the measure $`d\alpha _m^+(p)`$. The next assertions follow from elementary computations. $`\mathrm{}`$ If we use the definition of the annihilation operators $$\left(a(q;m)\mathrm{\Phi }\right)^{(n)}(p_1,\mathrm{},p_n)=\sqrt{n+1}\mathrm{\Phi }^{(n+1)}(q,p_1,\mathrm{},p_n)$$ (2.0.5) then we immediately get the identity: $$𝒰_\lambda a(q;m)𝒰_\lambda ^1=\lambda a(\lambda ^1q;\lambda ^1m).$$ (2.0.6) By hermitian conjugation we get a similar identity for the creation operators $`a^{}(q)`$. The expression of the real scalar field of mass $`m`$ is: $$\varphi (x;m)\frac{1}{(2\pi )^{3/2}}𝑑\alpha _m^+(p)\left[e^{ixp}a(p;m)+e^{ixp}a^{}(p;m)\right]$$ (2.0.7) so we get from (2.0.6) the following relation: $$𝒰_\lambda \varphi (x;m)𝒰_\lambda ^1=\lambda \varphi (\lambda x;\lambda ^1m).$$ (2.0.8) ###### Remark 2.2 There is an alternative point of view. One can define the operators $`𝒰_\lambda ^{}:_m_m`$ according to $$\left(𝒰_\lambda ^{}\mathrm{\Phi }^{(n)}\right)(𝐩_1,\mathrm{},𝐩_n)=\underset{i=1}{\overset{n}{}}r_\lambda (𝐩_𝐢)\mathrm{\Phi }^{(n)}(\lambda 𝐩_1,\mathrm{},\lambda 𝐩_n)$$ (2.0.9) where $$r_\lambda (𝐩)\lambda ^{3/2}\sqrt{\frac{\omega _m(𝐩)}{\omega _m(\lambda 𝐩)}}.$$ (2.0.10) Because we have the cocycle identity $$r_\lambda (𝐩)r_\lambda ^{}(\lambda 𝐩)=r_{\lambda \lambda ^{}}(𝐩)$$ (2.0.11) the map $`\lambda 𝒰_\lambda ^{}`$ defined above is a representation of the multiplicative group $`_+`$ (the dilation) group in the Fock space of the scalar field. Moreover, the relations (2.0.4), (2.0.6) and (2.0.8) are valid only up terms of order $`O(m)`$ because we have $`r_\lambda (𝐩)=\lambda +O(m).`$ So, we see that some information is lost in this approach. It is easy to prove that relations of the same type as (2.0.8) are valid for other types of fields, namely fields of integer spin. This includes the electromagnetic potential, the Yang-Mills fields, the gravitational field and also the ghosts fields used in the process of quantization. For a Dirac field an important difference appears. Instead of (2.0.7) we have: $$\psi (x;M)\frac{1}{(2\pi )^{3/2}}𝑑\alpha _M^+(p)\left[e^{ixp}\underset{i=1}{\overset{2}{}}u_i(p;M)b_i(p;M)+e^{ixp}\underset{i=1}{\overset{2}{}}v_i(p;M)b_i^{}(p;M)\right]$$ (2.0.12) (see ) where $`b_i^\mathrm{\#}(p;M)`$ are the creation (annihilation) operators; the expressions $`u_i(p;M)`$ and $`v_i(p;M)`$ are solutions of the free Dirac equation of positive (negative) values. To have Poincaré covariance of the field operator $`\psi `$ one has to normalize in such a way these spinors such that we have: $$u_i(\lambda p;\lambda M)=\lambda ^{1/2}u_i(p;M),v_i(\lambda p;\lambda M)=\lambda ^{1/2}v_i(p;M).$$ (2.0.13) So we get instead of (2.0.8): $$𝒰_\lambda \psi (x;M)𝒰_\lambda ^1=\lambda ^{3/2}\psi (\lambda x;\lambda ^1M).$$ (2.0.14) We can obviously prove that the relations (2.0.4) are valid in the most general case, with fields of various spins. Let us note that if we apply to the relations (2.0.8) or (2.0.14) a derivation operator $`\frac{}{x^\mu }`$ we obtain a supplementary factor $`\lambda `$ in the right hand side. Finally, if $`W(x;𝐦)`$ is a Wick monomial in free fields of various masses $`𝐦=(m_1,\mathrm{},M_1,\mathrm{})`$ we obtain a generalization of the relations (2.0.8) and (2.0.14), namely: $$𝒰_\lambda W(x;𝐦)𝒰_\lambda ^1=\lambda ^{\omega (W)}W(\lambda x;\lambda ^1𝐦)$$ (2.0.15) where the number $`\omega (W)`$ is called the canonical dimension of the monomial $`W`$ and is computed according to the well known rule: one attributes to every integer (resp. half-integer) spin field the canonical dimension $`1`$ (resp. $`3/2`$) and to every derivative the canonical dimension $`1`$. Then one postulates that the canonical dimension is an additive function. One can extend these considerations to Wick monomials in many variables $`W(x_1,\mathrm{},x_n)`$. If the interaction Lagrangian of a model verifies a relation of the type (2.0.15) we say that the model is dilation (or scale)-covariant. It also well known that the canonical dimension of fields is an important property in renormalization theory. ## 3 Renormalization Theory ### 3.1 Bogoliubov Axioms We outline here the axioms of a multi-Lagrangian perturbation theory. Following Bogoliubov and Shirkov ideas, in one constructs the $`S`$-matrix as a formal series of operator valued distributions: $$S(𝐠)=1+\underset{n=1}{\overset{\mathrm{}}{}}\frac{i^n}{n!}_{^{4n}}𝑑x_1\mathrm{}𝑑x_nT_{j_1,\mathrm{},j_n}(x_1,\mathrm{},x_n)g_{j_1}(x_1)\mathrm{}g_{j_n}(x_n),$$ (3.1.1) where $`𝐠=\left(g_j(x)\right)_{j=1,\mathrm{}P}`$ is a multi-valued tempered test function in the Minkowski space $`^4`$ that switches the interaction and $`T_{j_1,\mathrm{},j_n}(x_1,\mathrm{},x_n)`$ are operator-valued distributions acting in the Fock space of some collection of free fields. These operator-valued distributions are called chronological products and verify some properties called in the following Bogoliubov axioms. It is necessary to note that there is a canonical projection $`pr`$ associating to the point $`x_i`$ the index $`j_i`$. One starts from a set of interaction Lagrangians $`T_j(x),j=1,\mathrm{},P`$ and tries to construct the whole series $`T_{j_1,\mathrm{},j_n},n2`$. The interaction Lagrangians must satisfy some requirements such like Poincaré invariance, hermiticity and causality. The natural candidates fulfilling these demands is a linearly independent set of Wick polynomials operating in the Fock space (describing a system of weakly interacting particles). The recursive process of constructing the chronological produces fixes the chronological products almost uniquely. We will study this arbitrariness in detail later. The physical $`S`$-matrix is obtained from $`S(𝐠)`$ taking the adiabatic limit which is , loosely speaking the limit $`g_j(x)1,j=1,\mathrm{},P.`$ We give here the set of axioms imposed on the chronological products $`T_{j_1,\mathrm{},j_n}`$ following the notations of . * First, it is clear that, without loosing generality, we can consider them completely symmetrical in all variables in the sense: $$T_{j_{\pi (1)},\mathrm{},j_{\pi (n)}}(x_{\pi (1)},\mathrm{}x_{\pi (p)})=T_p(x_1,\mathrm{}x_p),\pi 𝒫_p.$$ (3.1.2) * Next, we must have Poincaré invariance. Because we will also consider Dirac fields, we suppose that we have an unitary representation $`(a,A)𝒰_{a,A}`$ of the group $`inSL(2,)`$ (the universal covering group of the proper orthochronous Poincaré group $`𝒫_+^{}`$) and a finite dimensional representation $`AS(A)`$ of of the group $`SL(2,)`$ such that: $`𝒰_{a,A}T_{j_1,\mathrm{},j_n}(x_1,\mathrm{},x_p)𝒰_{a,A}^1=S(A^1)_{j_1k_1}\mathrm{}S(A^1)_{j_nk_n}\times `$ $`T_{k_1,\mathrm{},k_n}(\delta (A)x_1+a,\mathrm{},\delta (A)x_p+a),ASL(2,),a^4`$ (3.1.3) where $`SL(2,)A\delta (A)𝒫_+^{}`$ is the covering map. In particular, translation invariance is essential for implementing Epstein-Glaser scheme of renormalization. Sometimes it is possible to supplement this axiom by corresponding invariance properties with respect to inversions (spatial and temporal) and charge conjugation. For the standard model only the PCT invariance is available. * The central axiom seems to be the requirement of causality which can be written compactly as follows. Let us firstly introduce some standard notations. Denote by $`V^+\{x^4\|x^2>0,x_0>0\}`$ and $`V^{}\{x^4\|x^2>0,x_0<0\}`$ the upper (lower) lightcones and by $`\overline{V^\pm }`$ their closures. If $`X\{x_1,\mathrm{},x_m\}^{4m}`$ and $`Y\{y_1,\mathrm{},y_n\}^{4m}`$ are such that $`x_iy_j\overline{V^{}},i=1,\mathrm{},m,j=1,\mathrm{},n`$ we use the notation $`XY.`$ If $`x_iy_j\overline{V^+}\overline{V^{}},i=1,\mathrm{},m,j=1,\mathrm{},n`$ we use the notations: $`XY.`$ We use the compact notation $`T_J(X)T_{j_1,\mathrm{},j_n}(x_1,\mathrm{},x_n)`$ with the convention $$T_{\mathrm{}}(\mathrm{})\mathrm{𝟏}$$ (3.1.4) and by $`XY`$ we mean the juxtaposition of the elements of $`X`$ and $`Y`$. Then the causality axiom writes as follows: $$T_{J_1J_2}(X_1X_2)=T_{J_1}(X_1)T_{J_2}(X_2),X_1X_2;$$ (3.1.5) here $`J_i`$ are the indices corresponding to the the coordinates $`X_i`$ i.e $`J_ipr(X_i),i=1,2.`$ ¿From (3.1.5) one can derive easily: $$[T_{J_1}(X_1),T_{J_2}(X_2)]=0,\mathrm{if}X_1X_2.$$ (3.1.6) * The unitarity of the $`S`$-matrix can be expressed if one introduces, the formal series: $$\overline{S}(𝐠)=1+\underset{n=1}{\overset{\mathrm{}}{}}\frac{(i)^n}{n!}_{^{4n}}𝑑x_1\mathrm{}𝑑x_n\overline{T}_{j_1,\mathrm{},j_n}(x_1,\mathrm{},x_n)g_{j_1}(x_1)\mathrm{}g_{j_n}(x_n),$$ (3.1.7) where, by definition: $$(1)^{\|X\|}\overline{T}_J(X)\underset{r=1}{\overset{\|X\|}{}}(1)^r\underset{X_1,\mathrm{},X_rPart(X)}{}T_{J_1}(X_1)\mathrm{}T_{J_r}(X_r);$$ (3.1.8) here $`X_1,\mathrm{},X_r`$ is a partition of $`X`$, $`\|X\|`$ is the cardinal of the set $`X`$ and the sum runs over all partitions. In the lowest orders we have: $$\overline{T}_j(x)=T_j(x)$$ (3.1.9) and $$\overline{T}_{j_1j_2}(x_1,x_2)=T_{j_1j_2}(x_1,x_2)+T_{j_1}(x_1)T_{j_2}(x_2)+T_{j_2}(x_2)T_{j_1}(x_1).$$ (3.1.10) One calls the operator-valued distributions $`\overline{T}_{j_1,\mathrm{},j_n}(x_1,\mathrm{},x_n)`$ anti-chronological products. The series (3.1.7) is the inverse of the series (3.1.1) i.e. we have: $$\overline{S}(𝐠)=S(𝐠)^1$$ (3.1.11) in the sense of formal series. Then the unitarity axiom is then: $$\overline{T}_J(X)=T_J(X)^{},X.$$ (3.1.12) One can show that the following relations are identically verified: $$\underset{X_1,X_2Part(X)}{}(1)^{\|X_1\|}T_{J_1}(X_1)\overline{T}_{J_2}(X_2)=\underset{X_1,X_2Part(X)}{}(1)^{\|X_1\|}\overline{T}_{J_1}(X_1)T_{J_2}(X_2)=0.$$ (3.1.13) Also one has, similarly to (3.1.5): $$\overline{T}_{J_1J_2}(X_1X_2)=\overline{T}_{J_2}(X_2)\overline{T}_{J_1}(X_1),X_1X_2.$$ (3.1.14) A renormalization theory is the possibility to construct such a $`S`$-matrix starting from the first order terms: $`T_j(x),j=1,\mathrm{},P`$ which are linearly independent Wick polynomials called interaction Lagrangians which should verify the following axioms: $$𝒰_{a,A}T_j(x)𝒰_{a,A}^1=S(A^1)_{jk}T_k(\delta (A)x+a),ASL(2,),j=1,\mathrm{},P$$ (3.1.15) $$[T_j(x),T_k(y)]=0,x,y^4s.t.xy,j,k=1,\mathrm{},P$$ (3.1.16) and $$T_j(x)^{}=T_j(x),j=1,\mathrm{},P.$$ (3.1.17) Usually, these requirements are supplemented by covariance with respect to some discrete symmetries (like spatial and temporal inversions, or PCT), charge conjugations or global invariance with respect to some Lie group of symmetry. Some other restrictions follow from the requirement of the existence of the adiabatic limit, at least in the weak sense. The case of a single Lagrangian perturbation theory corresponds to $`P=1`$. In this case the expression $`T(x)=T_1(x)`$ is the interaction Lagrangian and the chronological products are $`T(X)T_{1\mathrm{}1}(X)`$. More generally, one can consider that the interaction Lagrangian is $$T(x)=c_jT_j(x)$$ (3.1.18) with $`c_j`$ some real constants. In this case, the chronological products of the theory are $$T(X)=c_{j_1}\mathrm{}c_{j_n}T_{j_1,\mathrm{},j_n}(X).$$ (3.1.19) ### 3.2 Epstein-Glaser Induction We summarize the steps of the inductive construction of Epstein and Glaser , . Let the interaction Lagrangians $`T_j(x),j=1,\mathrm{},P`$ be some linearly independent Wick monomials acting in a certain Fock space with $`\omega _j,j=1,\mathrm{},P`$ the corresponding canonical dimensions. We suppose that they generate the space of all Wick monomials of canonical dimension less that $`4`$. The causality property (3.1.16) is fulfilled, but we must make sure that we also have (3.1.15) and (3.1.17). Moreover, a certain generalization of the preceding formalism is needed in order express the operator-valued chronological product in terms of numerical distributions . It is convenient let the index $`j`$ to run from $`0`$ to $`P`$ and to give, by definition $$T_0\mathrm{𝟏}.$$ (3.2.1) Next, we define the summ $`j_1+j_2`$ of two indices $`j_1,j_2=0,\mathrm{},P`$ through the relation $$T_{j_1+j_2}(x)=:T_{j_1}(x)T_{j_2}(x):$$ (3.2.2) and then we extend the summation operation to $`n`$-uples of indices $`J=(j_1,\mathrm{},j_n)`$ componentwise. We will use the notation $$\omega _J\underset{jJ}{}\omega _j$$ (3.2.3) and we call it the canonical dimension of $`T_J(X)`$. We suppose that we have constructed the chronological products $`T_{j_1,\mathrm{},j_p}(x_1,\mathrm{},x_p)`$ (for all $`p=1,\mathrm{},n1`$) having the following properties: (3.1.2), (3.1.5) and (3.1.12) for $`pn1`$, (3.1.5) for $`\|X_1\|+\|X_2\|n1`$ and (3.1.6) for $`\|X_1\|,\|X_2\|n1`$. Moreover, we suppose that we have the following Wick expansion of the chronological products: $$T_J(X)=\underset{K+L=J}{}t_K(X):T_{l_1}(x_1)\mathrm{}T_{l_n}(x_n):$$ (3.2.4) for $`\|X\|n1`$; here $`t_K(X)`$ are numerical distributions (called renormalized Feynman amplitudes) with degree of singularity restricted by the following relation: $$\omega (t_K)\omega _K4(n1).$$ (3.2.5) Let us notice that from (3.2.4) we have: $$t_J(X)=(\mathrm{\Phi }_0,T_J(X)\mathrm{\Phi }_0).$$ (3.2.6) We want to construct the distribution-valued operators $`T_J(X),\|X\|=n`$ such that the the properties above go from $`1`$ to $`n`$. Here are the steps of the construction. 1. One constructs from $`T_J(X),\|X\|n1`$ the expressions $`\overline{T}_J(X),\|X\|n1`$ according to (3.1.8) and proves the properties (3.1.14) for $`\|X_1\|+\|X_2\|n1`$. 2. Next, we defines the expressions: $$A_{j_1,\mathrm{},j_n}^{}(x_1,\mathrm{},x_{n1};x_n)\underset{X_1,X_2Part(X)}{\overset{}{}}(1)^{\|X_2\|}T_{J_1}(X_1)\overline{T}_{J_2}(X_2),$$ (3.2.7) $$R_{j_1,\mathrm{},j_n}^{}(x_1,\mathrm{},x_{n1};x_n)\underset{X_1,X_2Part(X)}{\overset{}{}}(1)^{\|X_2\|}\overline{T}_{J_1}(X_1)T_{J_2}(X_2)$$ (3.2.8) where the sum $`^{}`$ goes over the partitions of $`X=\{x_1,\mathrm{},x_n\}`$ such that $`X_2\mathrm{},x_nX_1.`$ Next, we construct the expression $$D_{j_1,\mathrm{},j_n}(x_1,\mathrm{},x_{n1};x_n)A_{j_1,\mathrm{},j_n}^{}(x_1,\mathrm{},x_{n1};x_n)R_{j_1,\mathrm{},j_n}^{}(x_1,\mathrm{},x_{n1};x_n).$$ (3.2.9) and prove that it has causal support i.e. $`supp(D_{j_1,\mathrm{},j_n}(x_1,\mathrm{},x_{n1};x_n))\mathrm{\Gamma }^+(x_n)\mathrm{\Gamma }^{}(x_n)`$ where we use standard notations: $$\mathrm{\Gamma }^\pm (x_n)\{(x_1,\mathrm{},x_n)(^4)^n\|x_ix_nV^\pm ,i=1,\mathrm{},n1\}.$$ (3.2.10) 3. The distributions $`D_J(X)`$ can be written in a formula similar to (3.2.4): $$D_J(X)=\underset{K+L=J}{}d_K(X):T_{l_1}(x_1)\mathrm{}T_{l_n}(x_n):$$ (3.2.11) where $`d_K(X)`$ are numerical distributions; in analogy to (3.2.6) we have: $$d_J(X)=(\mathrm{\Phi }_0,D_J(X)\mathrm{\Phi }_0).$$ (3.2.12) It follows that the numerical distributions $`d_J(X)`$ have causal support i.e $`supp(d_J(X))\mathrm{\Gamma }^+(x_n)\mathrm{\Gamma }^{}(x_n)`$ and are $`SL(2,)`$-invariant. Moreover, their degree of singularity is restricted by $$\omega (d_K)\omega _K4(n1).$$ (3.2.13) 4. There exists a causal splitting $$d=ar,supp(a)\mathrm{\Gamma }^+(x_n),supp(r)\mathrm{\Gamma }^{}(x_n)$$ (3.2.14) which is also $`SL(2,)`$-invariant and such that the order of the singularity is preserved. So, there exists a $`SL(2,)`$-covariant causal splitting: $$D_{j_1,\mathrm{},j_n}(x_1,\mathrm{},x_{n1};x_n)=A_{j_1,\mathrm{},j_n}(x_1,\mathrm{},x_{n1};x_n)R_{j_1,\mathrm{},j_n}(x_1,\mathrm{},x_{n1};x_n)$$ (3.2.15) with $`supp(A_{j_1,\mathrm{},j_n}(x_1,\mathrm{},x_{n1};x_n))\mathrm{\Gamma }^+(x_n)`$ and $`supp(R_{j_1,\mathrm{},j_n}(x_1,\mathrm{},x_{n1};x_n))\mathrm{\Gamma }^{}(x_n)`$. The expressions $`A_n`$ and $`R_n`$ are the advanced (resp. retarded) products. 5. We have the relation $$D_{j_1,\mathrm{},j_n}(x_1,\mathrm{},x_{n1};x_n)^{}=(1)^{n1}D_{j_1,\mathrm{},j_n}(x_1,\mathrm{},x_{n1};x_n).$$ (3.2.16) The causal splitting obtained above can be chosen such that $$A_{j_1,\mathrm{},j_n}(x_1,\mathrm{},x_{n1};x_n)^{}=(1)^{n1}A_{j_1,\mathrm{},j_n}(x_1,\mathrm{},x_{n1};x_n).$$ (3.2.17) 6. Let us define $`T_{j_1,\mathrm{},j_n}(x_1,\mathrm{},x_n)A_{j_1,\mathrm{},j_n}(x_1,\mathrm{},x_{n1};x_n)A_{j_1,\mathrm{},j_n}^{}(x_1,\mathrm{},x_{n1};x_n)`$ $`R_{j_1,\mathrm{},j_n}(x_1,\mathrm{},x_{n1};x_n)R_{j_1,\mathrm{},j_n}^{}(x_1,\mathrm{},x_{n1};x_n).`$ (3.2.18) Then these expressions satisfy the $`SL(2,)`$-covariance, causality and unitarity conditions (3.1.3) (3.1.5) (3.1.6) and (3.1.12) for $`p=n`$. If we substitute $$T_{j_1,\mathrm{},j_n}(x_1,\mathrm{},x_n)\frac{1}{n!}\underset{\pi }{}T_{j_{\pi (1)},\mathrm{},j_{\pi (n)}}(x_{\pi (1)},\mathrm{},x_{\pi (n)})$$ (3.2.19) where the sum runs over all permutations of the numbers $`\{1,\mathrm{},n\}`$ then we also have the symmetry axiom (3.1.2). It is easy to see that the induction hypothesis is verified by the operators $`T_J(X),\|X\|=n`$ constructed in this way. ### 3.3 The Arbitrariness of the Chronological Products This problem was addressed in and , as we have already mention in the Introduction. We prefer to give an independent formulation based on the multi-Lagrangian Epstein-Glaser scheme presented above. We consider two solutions of the Bogoliubov axioms with the same “initial conditions” $`T_j,j=1,\mathrm{},P`$ chosen as a basis in the space of Wick monomials of degree $`4`$. We introduce the following notation: if $`X=\{x_1,\mathrm{},x_n\}`$ then $$\delta (X)\delta (x_1x_n)\mathrm{}\delta (x_{n1}x_n).$$ (3.3.1) We note the identity: $$\underset{x_lX}{}\frac{}{x_l^\mu }\delta (X)=0.$$ (3.3.2) Now, we have the following result: ###### Theorem 3.1 Let $`T_J(X)`$ and $`\stackrel{~}{T}_J(X)`$ be two solutions of the Bogoliubov axioms such that $`T_j(x)=\stackrel{~}{T}_j(x),j=1,\mathrm{},P`$ and both verify the restriction (3.2.5). Then we have the relation: $$\stackrel{~}{T}_J(X)=T_J(X)+\underset{r=1}{\overset{\|X\|1}{}}\frac{1}{r!}\underset{X_1,\mathrm{},X_rPart(X)}{}P_{J_1;k_1}(X_1)\mathrm{}P_{J_r;k_r}(X_r)T_{k_1,\mathrm{},k_r}(x_{i_1},\mathrm{},x_{i_r}),\|X\|2$$ (3.3.3) where summation over the indices $`k_1,\mathrm{},k_r=0,\mathrm{},P`$ is understood, $`P_{J;k}(X)`$ are distributions of the form $$P_{J;k}(X)=p_{J;k}()\delta (X)$$ (3.3.4) with $`p_{J;k}()`$ a Lorentz covariant polynomial with constant coefficients in the partial derivatives restricted by: $$deg(p_{J;k})+\omega _k\omega _J4(n1)$$ (3.3.5) and $`x_{i_p}X_p,p=1,\mathrm{},r`$. In the preceding equation, the convention $$P_{j;k}(X)\delta _{jk},\|X\|=1$$ (3.3.6) is understood. Proof: We use complete induction. For $`n=2`$ one obtains a possible expression $`T_{j_1j_2}(x_1,x_2)`$ by causally splitting the distribution $`D_{j_1j_2}(x_1,x_2)=[T_{j_1}(x_1),T_{j_2}(x_2)]`$. According to a general result in distribution splitting theory, two such splitting differ by a distribution with support in the set $`\{x_1=x_2\}`$ of the type $`_{k=0}^P[p_{j_1j_2;k}()\delta (x_1x_2)]T_k(x_2);`$ the limitation $`deg(p_{j_1j_2;k})+\omega _k\omega _{j_1}+\omega _{j_2}4`$ follows from the restrictions (3.2.5). The Lorentz covariance follows if we make the distribution splitting in a covariant way, which is known to be possible. We suppose that we have the expressions $`P_{J;k}(X),\|X\|n1`$ such that the formula from the statement is valid for $`\|X\|n1`$; we prove the formula for $`\|X\|=n`$. It is convenient to observe that one can write the formula (3.3.3) as follows: $$\stackrel{~}{T}_J(X)=\underset{r=1}{\overset{\|X\|}{}}\frac{1}{r!}\underset{X_1,\mathrm{},X_rPart(X)}{}P_{J_1;k_1}(X_1)\mathrm{}P_{J_r;k_r}(X_r)T_{k_1,\mathrm{},k_r}(x_{i_1},\mathrm{},x_{i_r}),\|X\|2$$ (3.3.7) where the convention (3.3.6) has been used. So, by the induction hypothesis, we have the previous relation for $`\|X\|n1`$. Let us consider in this case the expression $$\mathrm{\Delta }_J(X)\stackrel{~}{T}_J(X)\underset{r=2}{\overset{\|X\|}{}}\frac{1}{r!}\underset{X_1,\mathrm{},X_rPart(X)}{}P_{J_1;k_1}(X_1)\mathrm{}P_{J_r;k_r}(X_r)T_{k_1,\mathrm{},k_r}(x_{i_1},\mathrm{},x_{i_r})$$ (3.3.8) and show that it has the support in the set $`x_1=x_2=\mathrm{}=x_n`$. For this, let us suppose that the point $`(x_1,\mathrm{},x_n)`$ is outside this set. Then one can find a Cauchy surface separating this set in two non-void subsets $`Y`$ and $`Z`$ such that $`[Y][Z]`$. Because of the symmetry axiom (3.1.2) we can suppose, without loosing generality, that $`Y=\{x_1,\mathrm{},x_i\}`$ and $`Z=\{x_{i+1},\mathrm{},x_n\}`$. In that case, let us notice that in the sum appearing in the preceding formula we can have non-zero contributions only from those partitions $`X_1,\mathrm{},X_r`$ such that for every $`p=1,\mathrm{},r`$ we have either $`X_pY`$ or $`X_pZ`$. This means that for such a choice of $`(x_1,\mathrm{},x_n)`$ we have: $`\mathrm{\Delta }_J(X)\stackrel{~}{T}_J(X){\displaystyle \underset{2s+t\|X\|}{}}{\displaystyle \underset{Y_1,\mathrm{},Y_sPart(Y)}{}}{\displaystyle \frac{1}{s!}}{\displaystyle \underset{Z_1,\mathrm{},Z_tPart(Z)}{}}{\displaystyle \frac{1}{t!}}`$ $`P_{J_1;k_1}(Y_1)\mathrm{}P_{J_s;k_s}(Y_s)P_{J_{s+1};k_{s+1}}(Z_1)\mathrm{}P_{J_{s+t};k_{s+t}}(Z_t)T_{k_1,\mathrm{},k_{s+t}}(x_{i_1},\mathrm{},x_{i_{s+t}})`$ (3.3.9) with $`x_{i_p}Y_p,p=1,\mathrm{},s`$ and $`x_{i_{s+p}}Z_p,p=1,\mathrm{},t`$. Now, we can use in the right hand side the causality property (3.1.5) for the chronological products $`T_J(X)`$ and $`\stackrel{~}{T}_J(X)`$.. We have easily get $`\mathrm{\Delta }_J(X)=0`$. The support property of the distribution $`\mathrm{\Delta }_J(X)`$ is proved. Using Wick theorem and well known facts about the structure of numerical distribution with support included in the set $`x_1=x_2=\mathrm{}=x_n`$ we get the formula (3.3.3) for $`\|X\|=n`$. The Lorentz covariance follows like in the case $`n=2`$. This finished the induction. $`\mathrm{}`$ It is clear now why do we need the multi-Lagrangian generalisation of Epstein-Glaser formalism. Even if we work in a theory with a single Lagrangian, the best we can do is to choose it among the set of linearly independent Wick polynomials $`T_j`$ say, $`T(x)=T_1(x)`$ and the usual chronological products of a single Lagrangian theory are $`T(X)=T_{1\mathrm{}1}(X)`$ (see the end of Subsection 3.1). To sets of chronological products $`T(X)`$ and $`\stackrel{~}{T}(X)`$ with the same “initial condition” $`T(x)=\stackrel{~}{T}(x)`$ will be connected by a formula of the following type: $$\stackrel{~}{T}(X)=T(X)+\underset{r=1}{\overset{\|X\|1}{}}\frac{1}{r!}\underset{X_1,\mathrm{},X_rPart(X)}{}P_{k_1}(X_1)\mathrm{}P_{k_r}(X_r)T_{k_1,\mathrm{},k_r}(x_{i_1},\mathrm{},x_{i_r}),\|X\|2$$ (3.3.10) where we have denoted $`P_k(X)P_{\{1\mathrm{}1\};k}(X)`$ with $`\|X\|`$ entries of the figure $`1`$. So, in the difference between two solutions of the problem will certainly appear other chronological products that $`T(X)`$. ## 4 Dilation Covariance of the Chronological Products ### 4.1 A General Characterization of Dilation Properties We will use the result from the preceding Subsection to study the generic behaviour of the chronological products with respect to the dilation invariance operators which was defined in Section 2. More explicitly, we consider a certain choice for the chronological products and we will emphasize the mass dependence in an obvious way: $`T_J(X;m)`$; these expressions are not completely fixed for $`\|X\|2`$ because there is the possibility of finite renormalizations. Nevertheless, we make a concrete choice in accordance with Bogoliubov axioms and we have: ###### Proposition 4.1 We suppose that the framework from the preceding Section is valid. Then the following relations are valid for all $`\|X\|2`$: $`𝒰_\lambda T_J(X;m)𝒰_\lambda ^1=\lambda ^{\omega _J}T_J(\lambda X;\lambda ^1m)`$ $`+{\displaystyle \underset{r=1}{\overset{\|X\|1}{}}}{\displaystyle \frac{1}{r!}}{\displaystyle \underset{X_1,\mathrm{},X_rPart(X)}{}}P_{J_1;k_1;m;\lambda }(\lambda X_1)\mathrm{}P_{J_r;k_r;m;\lambda }(\lambda X_r)T_{k_1,\mathrm{},k_r}(\lambda x_{i_1},\mathrm{},\lambda x_{i_r};\lambda ^1m)`$ (4.1.1) where the distributions $`P_{J;k;m;\lambda }(X)`$ are of the form $$P_{J;k;m;\lambda }(X)=\underset{\alpha }{}c_{J;k;\alpha }(\lambda ,m)^\alpha \delta (X);$$ (4.1.2) here $`\alpha `$ are multi-indices and $`\|\alpha \|`$ is the corresponding length. Moreover, the following relation is verified: $$P_{J;k;m;1}=0c_{J;k;\alpha }(1,m)=0.$$ (4.1.3) Proof: Let us consider the following expressions $$T_J^\lambda (X)\lambda ^{\omega _J}T_J(X;\lambda ^1m),\stackrel{~}{T}_J^\lambda (X)𝒰_\lambda T_J(\lambda ^1X;m)𝒰_\lambda ^1,\|X\|2$$ (4.1.4) both acting in the same Fock space: $`_{\lambda ^1m}`$ and having the same “initial conditions” $$T_j^\lambda (x)\lambda ^{\omega _j}T_j(x),j=1,\mathrm{},P$$ (4.1.5) due to (2.0.15). Also, these expressions verify the Bogoliubov axioms: the unitarity and the causality are obvious, but for the Poincaré covariance one had to use the relation (2.0.4). We can apply theorem 3.1 and obtain that the difference between the two expressions $`\stackrel{~}{T}_J^\lambda (X)`$ and $`T_J^\lambda (X)`$ is a sum of the type appearing in the right hand side of the relation (3.3.3) but with the polynomials depending on the parameter $`\lambda `$. If we make the substitution $`X\lambda X`$ we get the relation from the statement. If we take $`\lambda =1`$ then we get (4.1.3). $`\mathrm{}`$ ###### Remark 4.2 If we change the choice of the chronological products $`T_J(X;m),\|X\|2`$ adding some finite renormalizations, then the distributions $`P_{J;k;m;\lambda }(X)`$ will change also by some “coboundary” contribution. We will use this freedom later to simplify their form. The central result of this paper describes the explicit $`\lambda `$-dependence of the polynomials appearing in the preceding proposition. We study separately the cases $`m0`$ and $`m=0`$. ### 4.2 The Case $`m0`$ In this case the following result hold: ###### Theorem 4.3 In the case $`m0`$ one can choose the chronological products such that the following relations are valid for any $`\|X\|2`$: $$𝒰_\lambda T_J(X;m)𝒰_\lambda ^1=\lambda ^{\omega _J}T_J(\lambda X;\lambda ^1m).$$ (4.2.1) Proof: Is done by induction. (i) First, we consider the case $`\|X\|=2`$. We start from the relation (4.1.1) from the preceding proposition and apply $`𝒰_\lambda ^{}\mathrm{}𝒰_\lambda ^{}^1`$. We easily obtain the cocycle identity $$P_{J;k;m;\lambda \lambda ^{}}(X)=\lambda ^{\omega _J}P_{J;k;\lambda ^1m;\lambda ^{}}(X)+(\lambda ^{})^{\omega _k}P_{J;k;m;\lambda }(\lambda ^1X).$$ (4.2.2) If we substitute here the expression (4.1.2) for $`P_{J;k;\lambda }(X)`$ one finds out immediately from the preceding cocycle identity that we have $$c(\lambda \lambda ^{};m)=\lambda ^{\omega _J}c(\lambda ^{};\lambda ^1m)+(\lambda ^{})^{4+\omega _k+\|\alpha \|}c(\lambda ;m).$$ (4.2.3) where we omit for simplicity the dependence on $`J,k`$ and $`\alpha `$. More conveniently, one defines $$d(\lambda ,m)\lambda ^{\omega _J}c(\lambda ,m)$$ (4.2.4) and has the cocycle identity: $$d(\lambda \lambda ^{},m)=d(\lambda ^{},\lambda ^1m)+(\lambda ^{})^sd(\lambda ,m)$$ (4.2.5) where we have denoted $$s4+\omega _k+\|\alpha \|\omega _J.$$ (4.2.6) ¿From (4.1.3) we have the “initial condition”: $$c(1,m)=0d(1,m)=0.$$ (4.2.7) The equation (4.2.5) can be analysed elementary: if we take $`m=\lambda `$ the following equation emerges: $$d(\lambda \lambda ^{},\lambda )=d(\lambda ^{},1)+(\lambda ^{})^sd(\lambda ,\lambda )$$ (4.2.8) or, if we make $`\lambda =M`$ and $`\lambda \lambda ^{}=\mathrm{\Lambda }`$ we have $$d(\mathrm{\Lambda },M)=d(\frac{\mathrm{\Lambda }}{M},1)+\left(\frac{\mathrm{\Lambda }}{M}\right)^sd(M,M).$$ (4.2.9) Now, we substitute this expression into the initial relation (4.2.5) we immediately get $$d(\frac{\lambda }{m},1)=\left(\frac{\lambda }{m}\right)^sd(\lambda ^1m,\lambda ^1m)$$ (4.2.10) which makes sense because $`m0`$. Finally, we substitute this expression into (4.2.9) and obtain the most general solution of the cohomological equation (4.2.5) $$d(\lambda ,m)=\lambda ^sp(m)p(\lambda ^1m)$$ (4.2.11) where $`p(m)=\frac{1}{m^s}d(m,m)`$. This proves, in particular, that the general solution of the cohomological equation (4.2.5) is in this case a coboundary. In the end we have $$c_{J;k;\alpha }(\lambda ,m)=\lambda ^{4+\|\alpha \|+\omega _k}p_{J;k;\alpha }(m)\lambda ^{\omega _J}p_{J;k;\alpha }(\lambda ^1m)$$ (4.2.12) where $`p_{J;k;\alpha }`$ are arbitrary functions on $`m`$. It follows that the polynomial $`P_{J;k;m;\lambda }`$ appearing in the relation (4.1.1) for $`\|X\|=2`$ has the generic form $$P_{J;k;m\lambda }(X)=\underset{\alpha }{}\left[\lambda ^{4+\|\alpha \|+\omega _k}p_{J;k;\alpha }(m)\lambda ^{\omega _J}p_{J;k;\alpha }(\lambda ^1m)\right]^\alpha \delta (X).$$ (4.2.13) Now, it is an easy exercise to prove that this expression is a “coboundary” in the sense that if we redefine the chronological products $`T_J(X;m)`$ in order $`2`$ according to $$\stackrel{~}{T}_J(X,m)T_J(X,m)\underset{\alpha }{}p_{J;k;\alpha }(m)\left[^\alpha \delta (X)\right]T_k(x_2,m)$$ (4.2.14) then we will have for $`\|X\|=2`$ a relation of the type (4.1.1) with $`P_{J;k;m;\lambda }\stackrel{~}{P}_{J;k;m;\lambda }=0`$. This proves the assertion from the statement in the case $`\|X\|=2`$ that is, we can redefine the chronological products in the second order such that we have the relation from the satement for $`\|X\|=2`$. (ii) We suppose that the formula from the statement is valid for $`2\|X\|n1`$ and we prove it for $`\|X\|=n`$. The induction hypothesis amounts to $$P_{J;k;m;\lambda }(X)=0,\|X\|n1.$$ (4.2.15) Then, we get from (4.1.1) for $`\|X\|=n`$ the following relation $$𝒰_\lambda T_J(X;m)𝒰_\lambda ^1=\lambda ^{\omega _J}T_J(\lambda X;\lambda ^1m)+P_{J;k;m;\lambda }(\lambda X)T_k(\lambda x_n;\lambda ^1m)\|X\|=n$$ (4.2.16) which is a relation of the same type as the relation for $`T_J(X,m),\|X\|=2`$ obtained above. One can obtain a cohomological relation for the polynomials $`P_{J;k;m;\lambda }`$ which coincides in fact with (4.2.2). For the coefficients $`c_{J;k;\alpha }(\lambda ,m)`$ we will obtain again an equation of the type (4.2.3): $$c(\lambda \lambda ^{};m)=\lambda ^{\omega _J}c(\lambda ^{};\lambda ^1m)+(\lambda ^{})^{4(\|X\|1)+\omega _k+\|\alpha \|}c(\lambda ;m).$$ (4.2.17) Then we find out that the relation (4.2.5) is valid with (4.2.6) modified to $$s=4(\|X\|1)+\|\alpha \|+\omega _k\omega _J.$$ (4.2.18) In the end, the most general expression of the polynomial $`P_{J;k;m;\lambda }`$ is $$P_{J;k;m\lambda }(X)=\underset{\alpha }{}\left[\lambda ^{4(\|X\|1)+\|\alpha \|+\omega _k}p_{J;k;\alpha }(m)\lambda ^{\omega _J}p_{J;k;\alpha }(\lambda ^1m)\right]^\alpha \delta (X).$$ (4.2.19) As before, one can make $`P_{J;k;\lambda }(X)=0`$ by a suitable redefinition of the chronological products $`T_J(X,m),\|X\|=n`$. The induction is finished. $`\mathrm{}`$ ### 4.3 The Case $`m=0`$ We remark that in the relation (4.1.1) the distributions $`P_{J;k;0;\lambda }`$ appear and for simplicity we can skip the entry $`0`$ from the index. We adopt the same convetion for the coefficients $`c_{J;k;\alpha }(\lambda ,0)`$ appearing in the generic expression (4.1.2). In this case $`m=0`$ we have a much more interesting result: ###### Theorem 4.4 In the case $`m=0`$ one can redefine the chronological products in such a way that the distributions $`P_{J;k;\lambda }(X)`$ are of the following form: $$P_{J;k;\lambda }(X)=\lambda ^{\omega _J}ln(\lambda )\underset{\|\alpha \|=\omega _J4(\|X\|1)\omega _k}{}c_{J;k;\alpha ;\lambda }^\alpha \delta (X).$$ (4.3.1) Proof: As before, is done by induction. (i) First, we consider the case $`\|X\|=2`$. We start from the relation (4.1.1) from the proposition 4.1.3 and apply $`𝒰_\lambda ^{}\mathrm{}𝒰_\lambda ^{}^1`$. We obtain the cocycle identity of the same type as (4.2.2): $$P_{J;k;\lambda \lambda ^{}}(X)=\lambda ^{\omega _J}P_{J;k;\lambda ^{}}(X)+(\lambda ^{})^{\omega _k}P_{J;k;\lambda }(\lambda ^1X)$$ (4.3.2) where now we have no mass dependence. Insetead of the relation (4.2.3) we get: $$c(\lambda \lambda ^{})=\lambda ^{\omega _J}c(\lambda ^{})+(\lambda ^{})^{4+\|\alpha \|+\omega _k}c(\lambda ).$$ (4.3.3) As before, one defines $`d(\lambda )`$ according to (4.2.4) and has the cocycle identity: $$d(\lambda \lambda ^{})=d(\lambda ^{})+(\lambda ^{})^sd(\lambda ),s4+\omega _k+\|\alpha \|\omega _J.$$ (4.3.4) Again we have from (4.1.3) the “initial condition”: $$c(1)=0d(1)=0.$$ (4.3.5) The equation (4.3.4) can be analysed elementary if we differentiate with respect to $`\lambda ^{}`$ and put $`\lambda ^{}=1.`$ The following differential equation emerges: $$\lambda d^{}(\lambda )=d_0+sd(\lambda )$$ (4.3.6) where $`d_0d^{}(1).`$ We have two cases: (a) $`s0`$ The homogeneous equation $`\lambda D^{}(\lambda )=sD(\lambda )`$ has the solution $`D(\lambda )=A\lambda ^s`$. With the methods of variation of constants, we look for a solution of the preceding equation of the form $`d(\lambda )=A(\lambda )\lambda ^s`$ with the initial condition $`A(1)=0.`$ The function $`A(\lambda )`$ will verify the equation: $$A^{}=d_0\lambda ^{s1}$$ (4.3.7) with the solution $$A(\lambda )=\frac{d_0}{s}(1\lambda ^s)$$ (4.3.8) ¿From here we get the solution $$d(\lambda )=\frac{d_0}{s}(\lambda ^s1)$$ (4.3.9) which verifies identically the initial equation (4.3.4). (b) $`s=0`$ The equation (4.3.6) becomes: $$\lambda d^{}(\lambda )=d_0$$ (4.3.10) with the solution $`d(\lambda )=d_0\mathrm{ln}(\lambda )`$ which, again, identically verifies the initial equation (4.3.4). We get in this case the solution: $$c(\lambda )=d_0\lambda ^{\omega _J}ln(\lambda ).$$ (4.3.11) In the end, we get, instead of the formula (4.2.13) $`P_{J;k;m\lambda }(X)={\displaystyle \underset{\|\alpha \|\omega _J4\omega _k}{}}\left[\lambda ^{4+\|\alpha \|+\omega _k}c_{J;k;\alpha }\lambda ^{\omega _J}c_{J;k;\alpha }\right]^\alpha \delta (X)`$ $`+{\displaystyle \underset{\|\alpha \|=\omega _J4\omega _k}{}}\lambda ^{\omega _J}ln(\lambda )c_{J;k;\alpha }^\alpha \delta (X).`$ (4.3.12) If we make a redefinition of the chronological products $`T_J(X,0),\|X\|=2`$ we can get rid of the first sum in the preceding expression as in the case $`m0`$. This means that one can take $$P_{J;k;m\lambda }(X)=\lambda ^{\omega _J}ln(\lambda )\underset{\alpha =\omega _J4\omega _k}{}c_{J;k;\alpha }^\alpha \delta (X)$$ (4.3.13) which proves the assertion from the statement in the case $`\|X\|=2.`$ (ii) We suppose that the formula from the statement is valid for $`2\|X\|n1`$ and we prove it for $`\|X\|=n.`$ We establish a cocycle identity for $`P_{J;k;\lambda }(X),\|X\|=n.`$ Instead of (4.3.2) we obtain in the same way: $`P_{J;k;\lambda \lambda ^{}}(X)=\lambda ^{\omega _J}P_{J;k;\lambda ^{}}(X)+(\lambda ^{})^{\omega _k}P_{J;k;\lambda }(\lambda ^1X)`$ $`+{\displaystyle \underset{r=2}{\overset{\|X\|1}{}}}{\displaystyle \frac{1}{r!}}{\displaystyle \underset{X_1,\mathrm{},X_rPart(X)}{}}P_{J_1;m_1;\lambda }(\lambda ^1X_1)\mathrm{}P_{J_r;m_r;\lambda }(\lambda ^1X_r)P_{m_1,\mathrm{},m_r;k;\lambda ^{}}(x_{i_1},\mathrm{},x_{i_r})`$ (4.3.14) This relation goes into (4.2.2) for $`n=2`$ because the sum disappears. The preceding relation gives, instead of (4.3.3) the following: $$c(\lambda \lambda ^{})=\lambda ^{\omega _J}c(\lambda ^{})+(\lambda ^{})^{4(\|X\|1)+\omega _k+\|\alpha \|}c(\lambda )+(\lambda \lambda ^{})^{\omega _J}ln(\lambda ^{})\underset{r=2}{\overset{\|X\|1}{}}c_rln^r(\lambda )$$ (4.3.15) where, again, the multi-index $`\alpha `$ was omitted and $`c_r`$ are some constants; their value will not be needed. If we define the function $`d(\lambda )`$ by (4.2.4) we get: $`d(\lambda \lambda ^{})=d(\lambda ^{})+(\lambda ^{})^sd(\lambda )+ln(\lambda ^{}){\displaystyle \underset{r=2}{\overset{\|X\|1}{}}}c_rln^r(\lambda ),`$ $`s4(\|X\|1)+\omega _k+\|\alpha \|\omega _J.`$ (4.3.16) We know certainly from the general theorem 4.1.3 that this equations must have solutions. The only problem is to determine the $`\lambda `$ dependence from the preceding equation. As before we get from this relation the differential equation: $$\lambda d^{}(\lambda )=d_0+sd(\lambda )+\underset{r=2}{\overset{\|X\|1}{}}c_rln^r(\lambda )$$ (4.3.17) We have the same cases as before. (a) $`s0`$ The homogeneous equation is again $`\lambda D^{}(\lambda )=sD(\lambda )`$ with the the solution $`D(\lambda )=A\lambda ^s`$. If we look for a solution of the equation (4.3.17) of the form $`d(\lambda )=A(\lambda )\lambda ^s`$ with the initial condition $`A(1)=0`$ we get for $`A(\lambda )`$ the equation: $$A^{}=\lambda ^{s1}\left[d_0+\underset{r=2}{\overset{\|X\|1}{}}c_rln^r(\lambda )\right].$$ (4.3.18) In the end, we get, instead of (4.3.9) the following expression for the functions $`d`$: $$d(\lambda )=\frac{d_0}{s}(\lambda ^s1)+\underset{r=1}{\overset{\|X\|1}{}}a_rln^r(\lambda )$$ (4.3.19) with $`a_r`$ some constants. We substitute in the original equation (4.3.16) for the function $`d`$ and obtain that the only possibility is to have $`a_r=0`$ (and $`c_r=0`$) so we have in fact the solution: $$d(\lambda )=\frac{d_0}{s}(\lambda ^s1).$$ (4.3.20) (b) $`s=0`$ The equation (4.3.6) becomes: $$\lambda d^{}(\lambda )=d_0+\underset{r=2}{\overset{\|X\|1}{}}c_rln^r(\lambda ).$$ (4.3.21) with the solution $$d(\lambda )=d_0\mathrm{ln}(\lambda )+\underset{r=2}{\overset{\|X\|1}{}}\frac{c_r}{r+1}ln^{r+1}(\lambda ).$$ (4.3.22) We substitute in the initial equation (4.3.16) and obtain that $`c_r=0`$ so $$d(\lambda )=d_0\mathrm{ln}(\lambda )c(\lambda )=d_0\lambda ^{\omega _J}\mathrm{ln}(\lambda ).$$ (4.3.23) It follows that the most general expression of the polynomials $`P_{J;k;\lambda }(X),\|X\|=n`$ is $`P_{J;k;m\lambda }(X)={\displaystyle \underset{\|\alpha \|\omega _J4(\|X\|1)\omega _k}{}}\left[\lambda ^{4(\|X\|1)+\|\alpha \|+\omega _k}c_{J;k;\alpha }\lambda ^{\omega _J}c_{J;k;\alpha }\right]^\alpha \delta (X)`$ $`+{\displaystyle \underset{\|\alpha \|=\omega _J4(\|X\|1)\omega _k}{}}\lambda ^{\omega _J}ln(\lambda )c_{J;k;\alpha }^\alpha \delta (X).`$ (4.3.24) If we make a redefinition of the chronological products $`T_J(X,0),\|X\|=n`$ we can get rid of the first sum in the preceding expression as in the case $`m0`$. This means that one can take $$P_{J;k;m\lambda }(X)=\lambda ^{\omega _J}ln(\lambda )\underset{\|\alpha \|=\omega _J4(\|X\|1)\omega _k}{}c_{J;k;\alpha }^\alpha \delta (X)$$ (4.3.25) which proves the assertion from the statement in the case $`\|X\|=n.`$ The induction is finished. $`\mathrm{}`$ If we substitute the preceding result into the proposition 4.1.3 we get the following result: ###### Theorem 4.5 The following relations are valid for any $`\|X\|2`$: $`𝒰_\lambda T_J(X;0)𝒰_\lambda ^1=\lambda ^{\omega _J}[T_J(\lambda X;0)+`$ $`{\displaystyle \underset{r=1}{\overset{\|X\|1}{}}}{\displaystyle \frac{ln^r(\lambda )}{r!}}{\displaystyle \underset{X_1,\mathrm{},X_rPart(X)}{}}\lambda ^{(\omega _{k_1}+\mathrm{}+\omega _{k_r})}P_{J_1;k_1}(X_1)\mathrm{}P_{J_r;k_r}(X_r)\times `$ $`T_{k_1,\mathrm{},k_r}(\lambda x_{i_1},\mathrm{},\lambda x_{i_r};0)]`$ (4.3.26) where the distributions $`P_{J;k}(X)`$ are of the form $$P_{J;k}(X)=\underset{\|\alpha \|=\omega _J4(\|X\|1)\omega _k}{}c_{J;k;\alpha }^\alpha \delta (X).$$ (4.3.27) We also have: $`𝒰_\lambda \overline{T}_J(X;0)𝒰_\lambda ^1=\lambda ^{\omega _J}[\overline{T}_J(\lambda X;0)+`$ $`{\displaystyle \underset{r=1}{\overset{\|X\|1}{}}}{\displaystyle \frac{ln^r(\lambda )}{r!}}{\displaystyle \underset{X_1,\mathrm{},X_rPart(X)}{}}\lambda ^{(\omega _{k_1}+\mathrm{}+\omega _{k_r})}\overline{P}_{J_1;k_1}(X_1)\mathrm{}\overline{P}_{J_r;k_r}(X_r)\times `$ $`\overline{T}_{k_1,\mathrm{},k_r}(\lambda x_{i_1},\mathrm{},\lambda x_{i_r};0)].`$ (4.3.28) Let us also remark that in the case $`m=0`$ the operators $`𝒰_\lambda `$ act in the same Fock space $`_0^+`$ and so, they form a unitary representation of the dilation group $`_+`$. This means that we can define the infinitesimal generatios of the dilations: let us consider the continuous unitary representation of the additive group $``$ given by $$V_\chi 𝒰_{exp(\chi )}$$ (4.3.29) and denote by $`D`$ its infinitesimal generator obtained via Stone-von-Neumann theorem: $$V_\chi =e^{i\chi D}.$$ (4.3.30) Then we have from (2.0.8) the following commutation relation: $$[D,\varphi (x)]i\left(1+x^\mu \frac{}{x^\mu }\right)\varphi (x).$$ (4.3.31) The infinitesimal form of the relations (2.0.15) and (4.3.26) are: $$[D,W(x)]i\left[\omega (W)+x^\mu \frac{}{x^\mu }\right]W(x).$$ (4.3.32) and respectively: $$[D,T_J(X)]i\left(\omega _J+\underset{lX}{}x^\mu \frac{}{x_l^\mu }\right)T_J(X)iP_{J;k}(X)T_k(x_n).$$ (4.3.33) ### 4.4 Scaling Properties of the Renormalized Feynman Amplitudes We translate the preceding results for the renormalized Feynman amplitudes. From this analysis one can obtain the asymptotic behaviour of these amplitudes as it is done in the classic paper of Weinberg . We use the expression (3.2.4) for $`T_J(X;m)`$ emphasizing the mass dependence: $$T_J(X;m)=\underset{K+L=J}{}t_K(X;m):T_{l_1}(x_1;m)\mathrm{}T_{l_n}(x_n;m):$$ (4.4.1) and we have: ###### Theorem 4.6 The following relations are verified: (1) in the case $`m0`$: $$t_J(X;m)=\lambda ^{\omega _J}t_J(\lambda X;\lambda ^1m);$$ (4.4.2) (2) in the case $`m=0`$: $`t_J(X;0)=\lambda ^{\omega _J}[t_J(\lambda X;0)+`$ $`{\displaystyle \underset{r=1}{\overset{\|X\|1}{}}}{\displaystyle \frac{ln^r(\lambda )}{r!}}{\displaystyle \underset{X_1,\mathrm{},X_rPart(X)}{}}\lambda ^{(\omega _{k_1}+\mathrm{}+\omega _{k_r})}P_{J_1;k_1}(X_1)\mathrm{}P_{J_r;k_r}(X_r)\times `$ $`t_{k_1,\mathrm{},k_r}(\lambda x_{i_1},\mathrm{},\lambda x_{i_r};0)].`$ (4.4.3) The proof is done using the formula (3.2.6) into the relations (4.2.1) and (4.3.26). The preceding theorem elucidates the logarithmic behaviour of the renormalized Feynman amplitudes in the case $`m=0`$. Presumably, the terms proportional with $`ln^r`$ correspond to graphs with $`r`$ loops. One can obtain the infinitesimal form of the preceding relation: we make $`\lambda =e^\chi `$, differentiate with respect to the variable $`\chi `$ and put $`\chi =0`$. If we take into account that $$t_j(x)=\delta _{j,0}$$ (4.4.4) the following relations emerges: (1) for $`m0`$: $$\left(\underset{l=1}{\overset{n}{}}x_l^\mu \frac{}{x_l^\mu }m\frac{}{m}+\omega _J\right)t_{J;K}(X;m)=0.$$ (4.4.5) If we take into account translation invariance, we can express the Feynman amplitudes $`t_J(X;m)`$ in terms of translation-invariant variables: $`\xi _ix_ix_n,i=1,\mathrm{},n1`$ and we have: $$\left(\underset{l=1}{\overset{n1}{}}\xi _l^\mu \frac{}{\xi _l^\mu }m\frac{}{m}+\omega _J\right)t_J(\mathrm{\Xi };m)=0$$ (4.4.6) or, for the Fourier transform: $$\left(\underset{l=1}{\overset{n1}{}}p_l^\mu \frac{}{p_l^\mu }+m\frac{}{m}\omega _J\right)\stackrel{~}{t}_{J;K}(P;m)=0.$$ (4.4.7) (2) for $`m=0`$: $$\left(\underset{l=1}{\overset{n}{}}x_l^\mu \frac{}{x_l^\mu }+\omega _J\right)t_{J;K}(X;m)+P_{J;0}(X)=0$$ (4.4.8) or, in translationally invariant variables: $$\left(\underset{l=1}{\overset{n1}{}}\xi _l^\mu \frac{}{\xi _l^\mu }+\omega _J\right)t_J(\mathrm{\Xi };m)+P_{J;0}(\mathrm{\Xi })=0$$ (4.4.9) or, for the Fourier transforms: $$\left(\underset{l=1}{\overset{n1}{}}p_l^\mu \frac{}{p_l^\mu }\omega _J\right)\stackrel{~}{t}_{J;K}(P;m)+\stackrel{~}{P}_{J;0}(P)=0.$$ (4.4.10) ### 4.5 The General Case Suppose that we have a theory with a finite number of fields, some of them of zero-mass and some of non-zero mass. We suppose that there exists at least one field of non-zero-mass. Then one can implement the analysis from the case $`m0`$ in the following way. Let us denote the non-zero masses of the theory as $`𝐦(m_1,\mathrm{},m_t)`$; then we have instead of (4.2.3) the relation: $$c(\lambda \lambda ^{};𝐦)=\lambda ^{\omega _J}c(\lambda ^{};\lambda ^1𝐦)+(\lambda ^{})^{4+\omega _k+\|\alpha \|}c(\lambda ;𝐦).$$ (4.5.1) It is covenient to work in “polar” coordinates $`(m,\mu )^{}\times S^{t1}`$ where $`m\|𝐦\|`$ and $`\mu _i\frac{m_i}{m},i=1,\mathrm{},t`$. Then, the previous relation wites as follows: $$c(\lambda \lambda ^{};m,\mu )=\lambda ^{\omega _J}c(\lambda ^{};\lambda ^1m,\mu )+(\lambda ^{})^{4+\omega _k+\|\alpha \|}c(\lambda ;m,\mu )$$ (4.5.2) and we see that the variables $`\mu S^{t1}`$ play no rôle. Then one can implement the proof of the theorem 4.2.1 without any change. So, it follows that if at least a non-zero mass particle is present in the theory, then the conclusions of the theorem 4.2.1 and of the case (1) considered above for the numerical distributions are true in this case also. More precisely, we have: $$𝒰_\lambda T_J(X;𝐦)𝒰_\lambda ^1=\lambda ^{\omega _J}T_J(\lambda X;\lambda ^1𝐦),$$ (4.5.3) $$\left(\underset{l=1}{\overset{n}{}}x_l^\mu \frac{}{x_l^\mu }\underset{i=1}{\overset{t}{}}m_i\frac{}{m_i}+\omega _J\right)t_{J;K}(X;𝐦)=0,$$ (4.5.4) $$\left(\underset{l=1}{\overset{n1}{}}\xi _l^\mu \frac{}{\xi _l^\mu }\underset{i=1}{\overset{t}{}}m_i\frac{}{m_i}+\omega _J\right)t_J(\mathrm{\Xi };𝐦)=0$$ (4.5.5) and $$\left(\underset{l=1}{\overset{n1}{}}p_l^\mu \frac{}{p_l^\mu }+\underset{i=1}{\overset{t}{}}m_i\frac{}{m_i}\omega _J\right)\stackrel{~}{t}_{J;K}(P;𝐦)=0.$$ (4.5.6) Because these relations follow from scale invariance, they can be called the Callan-Symanzik equation in the framework of Epstein-Glaser perturbative scheme. However, we do not obtain the anomalous dimension in this way. The usual Callan-Symanzik equation , expresses the action of the (infinitesimal) dilation operator on the generating functional of the Green function of the interacting field, but it is natural to suppose that the two formalisms to be, in some way, equivalent. We will comment more in the last Section about this point where we will indicate the way to connect our result to the standard arguments based on the action principle. We close this Section with an important remark. The results contained in the Subsections 4.2 and in the general case here are rather natural. Indeed, in the case fron Subsection 4.2 of a fields of mass $`m>0`$ one could preceed more directly as follows. One can make a choice for the chronological products $`T_J(X;m_0)`$ for a certain fixed mass $`m_0>0`$ and define the chronological products $`T_J(X;m)`$ for any other mass $`m`$ such that one has the simple behaviour described by the equation (4.2.1). The argument can be immediatedly adapted to the general case of more mass, but with at least one non-zero. It is nevertheless, interesting to work out these cases from pure cohomological considerations. So, it follows that the really non-trivial case is the one when all particles have null masses and when, in principle, one cannot avoid the non-trivial cocycles of logarithmic type. Let us also remark that another way of proving the result from Subsection 4.2 is by observing that in this case one can apply the central solution for the distribution splitting and in this way the scaling properties of the numerical distributions are preserved. ## 5 Yang-Mills Theories In this Section we analyse the scale covariance of the Standard Model (SM) and the consequences of this property for the structure of possible anomalies. ### 5.1 The Fock Space of the Bosons We give some basic facts about the quantization of a spin $`1`$ Boson of mass $`m>0`$. One can proceed in a rather close analogy to the case of the photon; for more details see and references quoted there. Let us denote the Hilbert space of the Boson by $`\mathrm{H}_m`$; it carries the unitary representation of the orthochronous Poincaré group $`𝖧^{[m,1]}`$. The Hilbert space of the multi-Boson system should be, as before, the associated symmetric Fock space $`_m^+(\mathrm{H}_m)`$. We construct this Fock space as before in the spirit of algebraic quantum field theory. One considers the Hilbert space $`^{gh}`$ generated by applying on the vacuum $`\mathrm{\Phi }_0`$ the free fields $`A^\mu (x),u(x),\stackrel{~}{u}(x),\mathrm{\Phi }(x)`$ which are completely characterize by the following properties: * Equation of motion: $$(\mathrm{}+m^2)A^\mu (x),(\mathrm{}+m^2)u(x)=0,(\mathrm{}+m^2)\stackrel{~}{u}(x)=0,(\mathrm{}+m^2)\mathrm{\Phi }(x)=0.$$ (5.1.1) * Canonical (anti)commutation relations: $`[A^\mu (x),A^\rho (y)]=g^{\mu \rho }D_m(xy)\times \mathrm{𝟏},`$ $`[A^\mu (x),u(y)]=0,[A^\mu (x),\stackrel{~}{u}(y)]=0,[A^\mu (x),\mathrm{\Phi }(y)]=0,`$ $`\{u(x),u(y)\}=0,\{\stackrel{~}{u}(x),\stackrel{~}{u}(y)\}=0,\{u(x),\stackrel{~}{u}(y)\}=D_m(xy)\times \mathrm{𝟏},`$ $`[\mathrm{\Phi }(x),\mathrm{\Phi }(y)]=D_m(xy)\times \mathrm{𝟏},[\mathrm{\Phi }(x),u(y)]=0.`$ (5.1.2) * $`SL(2,)`$-covariance: $`U_{a,A}A^\mu (x)U_{a,A}^1=\delta (A^1)_{}^{\mu }{}_{\nu }{}^{}A^\nu (\delta (A)x+a),`$ $`U_{a,A}u(x)U_{a,A}^1=u(\delta (A)x+a),U_{a,A}\stackrel{~}{u}(x)U_{a,A}^1=\stackrel{~}{u}(\delta (A)x+a)`$ $`U_{a,A}\mathrm{\Phi }(x)U_{a,A}^1=\mathrm{\Phi }(\delta (A)x+a)`$ (5.1.3) * PCT covariance. $`U_{PCT}A_\mu (x)U_{PCT}^1=A_\mu (x),U_{PCT}\mathrm{\Phi }(x)U_{PCT}^1=\mathrm{\Phi }(x)`$ $`U_{PCT}u(x)U_{PCT}^1=u(x),U_{PCT}\stackrel{~}{u}(x)U_{PCT}^1=\stackrel{~}{u}(x),`$ $`U_{PCT}\mathrm{\Phi }_0=\mathrm{\Phi }_0.`$ (5.1.4) ###### Remark 5.1 Although we could give the expressions for $`U_{I_s},U_{I_t}`$ and $`U_C`$ separately, we prefer to give only the expression of the PCT transform because the interaction Lagrangian of the standard model is not invariant with respect to these three operations but it is PCT-covariant. We give as before in $`^{gh}`$ the sesqui-linear form $`<,>`$ which is completely characterize by requiring: $$A_\mu (x)^{}=A_\mu (x),u(x)^{}=u(x),\stackrel{~}{u}(x)^{}=\stackrel{~}{u}(x),\mathrm{\Phi }(x)^{}=\mathrm{\Phi }(x).$$ (5.1.5) Now, the expression of the supercharge gets an extra term: $$Q=_^3d^3x\left[^\mu A_\mu (x)+m\mathrm{\Phi }(x)\right]\stackrel{}{_0}u(x)$$ (5.1.6) and one can see that we have $`[Q,A_\mu ]=i_\mu u,\{Q,u\}=0,\{Q,\stackrel{~}{u}\}=i(_\mu A^\mu +m\mathrm{\Phi }),[Q,\mathrm{\Phi }]=imu`$ (5.1.7) We still have $$Q^2=0Im(Q)Ker(Q)$$ (5.1.8) and also $$U_{a,A}Q=QU_{a,A},U_{PCT}Q=QU_{PCT}.$$ (5.1.9) Finally: ###### Theorem 5.2 The sesqui-linear form $`<,>`$ factorizes to a well-defined scalar product on the completion of the factor space $`Ker(Q)/Im(Q)`$. Then there exists the following Hilbert spaces isomorphism: $$\overline{Ker(Q)/Im(Q)}_m;$$ (5.1.10) The representation of the Poincaré group and the PCT operator are factorizing to $`Ker(Q)/Im(Q)`$ and are producing unitary operators (resp. an anti-unitary operator). If $`𝒲`$ the linear space of all Wick monomials in the fields $`A_\mu ,u,\stackrel{~}{u}`$ and $`\mathrm{\Phi }`$ acting in the Fock space $`^{gh}`$ then the expression of the BRST operator is determined by $`d_Qu=0,d_Q\stackrel{~}{u}=i(^\mu A_\mu +m\mathrm{\Phi }),d_QA_\mu =i_\mu u,d_Q\mathrm{\Phi }=imu.`$ (5.1.11) and, as a consequence we have $$d_Q^2=0.$$ (5.1.12) If one adds matter fields we proceed as before. In particular, this will mean that the BRST operator acts trivially on the matter fields. Now we can define the Yang-Mills field. We must consider the case when we have $`r`$ fields of spin $`1`$ and some of them will have zero mass and the others will be considered of non-zero mass. Apparently, we need the scalar ghosts only in the last case. However it can be shown that with this assumption, there are no non-trivial models. To avoid this situation, we make the following amendment. All the fields considered above will carry an additional index $`a=1,\mathrm{},r`$ i.e. we have the following set of fields: $`A_{a\mu },u_a,\stackrel{~}{u}_a,\mathrm{\Phi }_aa=1,\mathrm{},r.`$ If one of the fields $`A_{a\mu }`$ has zero mass we postulate that the corresponding scalar fields $`\mathrm{\Phi }_a`$ are physical fields and they will be called Higgs fields. Moreover, we do not have to assume that they are massless i.e. if some Boson field $`A_a^\mu `$ has zero mass $`m_a=0`$, we can suppose that the corresponding Higgs field $`\mathrm{\Phi }_a`$ has a non-zero mass: $`m_a^H`$. It is convenient to use the compact notation $$m_a^{}\{\begin{array}{ccc}\hfill m_a& \text{for}& m_a0\hfill \\ \hfill m_a^H& \text{for}& m_a=0\hfill \end{array}$$ (5.1.13) These fields verify the following equations of motion: $`(\mathrm{}+m_a^2)A_a^\mu (x)=0,(\mathrm{}+m_a^2)u_a(x)=0,(\mathrm{}+m_a^2)\stackrel{~}{u}_a(x)=0,(\mathrm{}+(m_a^{})^2)\mathrm{\Phi }_a(x)=0`$ (5.1.14) The rest of the formalism stays unchanged. The canonical (anti)commutation relations are: $`[A_{a\mu }(x),A_{b\nu }(y)]=\delta _{ab}g_{\mu \nu }D_{m_a}(xy)\times \mathrm{𝟏},`$ $`\{u_a(x),\stackrel{~}{u}_b(y)\}=\delta _{ab}D_{m_a}(xy)\times \mathrm{𝟏},[\mathrm{\Phi }_a(x),\mathrm{\Phi }_b(y)]=\delta _{ab}D_{m_a^{}}(xy)\times \mathrm{𝟏};`$ (5.1.15) and all other (anti)commutators are null. The supercharge is given by $$Q=\underset{a=1}{\overset{r}{}}_^3d^3x\left[^\mu A_{a\mu }(x)+m_a\mathrm{\Phi }_a(x)\right]\stackrel{}{_0}u_a(x)$$ (5.1.16) and verifies all the expected properties. The Krein operator is determined by: $$A_{a\mu }(x)^{}=A_{a\mu }(x),u_a(x)^{}=u_a(x),\stackrel{~}{u}_a(x)^{}=\stackrel{~}{u}_a(x),\mathrm{\Phi }_a(x)^{}=\mathrm{\Phi }_a(x).$$ (5.1.17) The ghost degree is defined in an obvious way and the expression of the BRST operator is similar to the previous one. In particular we have: $$d_Qu_a=0,d_Q\stackrel{~}{u}_a=i(_\mu A_a^\mu +m_a\mathrm{\Phi }_a),d_QA_a^\mu =i^\mu u_a,d_Q\mathrm{\Phi }_a=im_au_a,a=1,\mathrm{},r.$$ (5.1.18) Finally, the condition of gauge invariance is (see ): $$d_QT(X)=i\underset{x_lX}{}\frac{}{x_l^\mu }T_l^\mu (X)$$ (5.1.19) for some Wick polynomials $`T_l^\mu (X),l=1,\mathrm{},\|X\|`$. ### 5.2 Matter Fields and the Interaction Lagrangian of the SM In this case the matter field is a set of Dirac fields of mass $`M_A,A=1,\mathrm{},N`$ denoted by $`\psi _A(x)`$. These fields are characterized by the following relations ; here $`A,B=1,\mathrm{},N`$: * Equation of motion: $$(i\gamma +M_A)\psi _A(x)=0.$$ (5.2.1) * Canonical (anti)commutation relations: $`[\psi _A(x),A_a^\mu (y)]=0,[\psi _A(x),u_a(y)]=0,[\psi _A(x),\stackrel{~}{u}_a(y)]=0,[\psi _A(x),\mathrm{\Phi }_a(y)]=0`$ $`\{\psi _A(x),\psi _B(y)\}=0,\{\psi _A(x),\overline{\psi _B}(y)\}=\delta _{AB}S_{M_A}(xy)\times \mathrm{𝟏}.`$ (5.2.2) * Covariance properties with respect to the Poincaré group: $$U_{a,A}\psi _A(x)U_{a,A}^1=S(A^1)\psi _A(\delta (A)x+a).$$ (5.2.3) * PCT-covariance: $$U_{PCT}\psi _A(x)U_{I_s}^1=\gamma _1\gamma _2\gamma _3\overline{\psi _A}(x)^t.$$ (5.2.4) The condition of gauge invariance remains the same (5.1.19) and one can prove that this condition for $`n=1,2`$ determines quite drastically the interaction Lagrangian of canonical dimension $`\omega (T(x))=4`$: $`T(x)f_{abc}[:A_{a\mu }(x)A_{b\nu }(x)^\nu A_a^\mu (x)::A_a^\mu (x)u_b(x)_\mu \stackrel{~}{u}_c(x):],`$ $`+f_{abc}^{}[:\mathrm{\Phi }_a(x)_\mu \mathrm{\Phi }_b(x)A_c^\mu (x):m_b:\mathrm{\Phi }_a(x)A_{b\mu }(x)A_c^\mu (x):m_b:\mathrm{\Phi }_a(x)\stackrel{~}{u}_b(x)u_c(x):]`$ $`+f_{abc}^\mathrm{"}:\mathrm{\Phi }_a(x)\mathrm{\Phi }_b(x)\mathrm{\Phi }_c(x):+j_a^\mu (x)A_{a\mu }(x)+j_a(x)\mathrm{\Phi }_a(x)`$ (5.2.5) where: $$j_a^\mu (x)=:\overline{\psi _A}(x)(t_a)_{AB}\gamma ^\mu \psi _B(x):+:\overline{\psi _A}(x)(t_a^{})_{AB}\gamma ^\mu \gamma _5\psi _B(x):$$ (5.2.6) and $$j_a(x)=:\overline{\psi _A}(x)(s_a)_{AB}\psi _B(x):+:\overline{\psi _A}(x)(s_a^{})_{AB}\gamma _5\psi _B(x):$$ (5.2.7) are the so-called currents. The conditions of $`SL(2,)`$ and PCT-covariance of the interaction Lagrangian are easy to prove as well as the causality condition. The hermiticity conditions are equivalent to the fact that the complex $`N\times N`$ matrices $`t_a,t_a^{},s_a,a=1,\mathrm{}r`$ are hermitian and $`s_a^{},a=1,\mathrm{},r`$ is anti-hermitian. The constants $`f_{abc}`$ are completely anti-symmetric and verify Jacobi identity so they generate a compact semi-simple Lie group quite naturally. There are other conditions on the rest of the constants as well, but because we do not need these properties in the subsequent analysis, we refer to the literature , and references quoted there. Moreover, it can be proved that the condition of gauge invariance (5.1.19) is valid for $`n=1,2`$ and we can take $`T^\mu (x)`$ to be of canonical dimension $`\omega (T^\mu (x))=4`$ with the explicit form: $`T^\mu (x)=f_{abc}(:u_a(x)A_{b\nu }(x)F_c^{\nu \mu }(x):{\displaystyle \frac{1}{2}}:u_a(x)u_b(x)^\mu (x)\stackrel{~}{u}_c(x):)`$ $`+f_{abc}^{}(m_a:A_a^\mu (x)\mathrm{\Phi }_b(x)u_c(x):+:\mathrm{\Phi }_a(x)^\mu \mathrm{\Phi }_b(x)u_c(x):).+u_a(x)j_a^\mu (x).`$ (5.2.8) The following relations are verified: * $`SL(2,)`$-covariance: for any $`ASL(2,)`$ we have $$U_{a,A}T(x)U_{a,A}^1=T(\delta (A)x+a),U_{a,A}T^\mu (x)U_{a,A}^1=\delta (A^1)_{}^{\mu }{}_{\rho }{}^{}T^\rho (\delta (A)x+a).$$ (5.2.9) * PCT-covariance: $$U_{PCT}T(x)U_{PCT}^1=T(x),U_{PCT}T^\mu (x)U_{PCT}^1=T^\mu (x).$$ (5.2.10) * Causality: $$[T(x),T(y)]=0,[T^\mu (x),T^\rho (y)]=0,[T^\mu (x),T(y)]=0,x,y^4s.t.xy.$$ (5.2.11) * Unitarity: $$T(x)^{}=T(x),T^\mu (x)^{}=T^\mu (x).$$ (5.2.12) * Ghost content: $$gh(T(x))=0,gh(T^\mu (x))=0.$$ (5.2.13) We mention that in -, the condition of gauge invariance is analysed up to the order $`3`$. ### 5.3 Dilation Covariance of the Standard Model In this Subsection we generalize the arguments from the Sections 2.1 for the standard model. We denote the set of all masses by $`𝐦(m_a,m_a^{},M_A)_{a=1.\mathrm{},r;A=1,\mathrm{},N}`$ and the Fock space of all particles (physical or ghosts) by $`_𝐦^{gh}`$. This Hilbert space is generated from the vacuum by applying the operators: $`A_a^\mu (x;m_a),u_a(x;m_a),\stackrel{~}{u}_a(x;m_a),\mathrm{\Phi }_a(x;m_a^{})`$ and $`\psi _A(x;M_A)`$. We define the dilation operators in the total Hilbert space in analogy to (2.0.2) and the result from the proposition 2.1 stays true; we also have the commutations relations with the Poincaré (2.0.4). Finally, we have from (2.0.8) and (2.0.14): $`𝒰_\lambda A_a^\mu (x;m_a)𝒰_\lambda ^1=\lambda A_a^\mu (\lambda x;\lambda ^1m_a),𝒰_\lambda \mathrm{\Phi }_a(x;m_a)𝒰_\lambda ^1=\lambda \mathrm{\Phi }_a(\lambda x;\lambda ^1m_a),`$ $`𝒰_\lambda u_a(x;m_a)𝒰_\lambda ^1=\lambda u_a(\lambda x;\lambda ^1m_a),𝒰_\lambda \stackrel{~}{u}(x;m_a)𝒰_\lambda ^1=\lambda \stackrel{~}{u}(\lambda x;\lambda ^1m_a),`$ $`𝒰_\lambda \psi _A(x;M_A)𝒰_\lambda ^1=\lambda ^{3/2}\psi _A(\lambda x;\lambda ^1M_A),a=1,\mathrm{},r,A=1,\mathrm{},N.`$ (5.3.1) ¿From these relations and from the expressions (5.2.5) and (5.2.8) we obtain particular cases of the relation (2.0.15): $$𝒰_\lambda T_1(x;𝐦)𝒰_\lambda ^1=\lambda ^4T_1(\lambda x;\lambda ^1𝐦),𝒰_\lambda T_1^\mu (x;𝐦)𝒰_\lambda ^1=\lambda ^4T_1^\mu (\lambda x;\lambda ^1𝐦);$$ (5.3.2) this means that both expressions have canonical dimension equal to $`4`$ which is also the dimension of the Minkowski space-time. Let us suppose from now on that there are non-zero masses into the theory. Then we can apply the argument presented at the end of the previous Section and obtain that for all $`\|X\|1`$ we have the following formulæ: $`𝒰_\lambda T(X;𝐦)𝒰_\lambda ^1=\lambda ^{4\|X\|}T(\lambda X;\lambda ^1𝐦)`$ $`𝒰_\lambda T_l^\mu (X;𝐦)𝒰_\lambda ^1=\lambda ^{4\|X\|}T_l^\mu (\lambda X;\lambda ^1𝐦).`$ (5.3.3) We also mention the following result which easily follows from the definitions: ###### Lemma 5.3 The following relations is valid for every Wick monomial: $$𝒰_\lambda \left[d_QW(X;𝐦)\right]𝒰_\lambda ^1=\lambda ^{\omega (W)+1}W(\lambda X;\lambda ^1𝐦).$$ (5.3.4) Proof: If the expression $`W`$ is one of the fields $`A_a^\mu (x;m_a),u_a(x;m_a),\stackrel{~}{u}_a(x;m_a),\mathrm{\Phi }_a(x;m_a^{})`$ or $`\psi _A(x;M_A),\overline{\psi }_A(x;M_A)`$ the formula from the statement follows elementary; then we extend to any Wick monomial by induction, using the derivative properties of the BRST operator. $`\mathrm{}`$ ### 5.4 The Structure of the Anomalies in the Standard Model We consider the standard model as defined by the Lagrangian (5.2.5). and suppose that there are no anomalies up to the order $`n1`$ i.e. we have (5.1.19) up to this order. The purpose of this Subsection is to find if possible anomalous terms can appear in this relation in order $`n`$ and what limitation are imposed by scale covariance. The analysis is similar to the case of the quantum electrodynamics . However, we prefer to use the formalism developped in Subsections 3.1 and 3.2. (i) Suppose that we have constructed the chronological products $`T_J(X),\|X\|n1`$ verifying all the induction hypothesis from Subsection 3.2. Then we will be able to use the formulæ of the type (3.2.6). We must have, in analogy to (3.1.18), a developpment of the type: $$T^\mu (x)=c_j^\mu T_j(x)$$ (5.4.1) with $`c_j^\mu `$ some real constants; then we will have in analogy to (3.1.19): $$T_l^\mu (X)=c_{j_1}\mathrm{}c_{j_l}^\mu \mathrm{}c_{j_n}T_{j_1,\mathrm{},j_n}(X)$$ (5.4.2) for all $`\|X\|n1`$. In particular, the following conventions hold: $$T(\mathrm{})\mathrm{𝟏},T_l^\mu (\mathrm{})0,T_l^\mu (X)0,\mathrm{for}x_lX.$$ (5.4.3) We supplement the induction hypothesis adding: * ghost number content: $$gh(T(X))=0,gh(T_l^\mu (X))=1,\|X\|n1;$$ (5.4.4) * gauge invariance: $$d_QT(X)=i\underset{l=1}{\overset{n}{}}\frac{}{x_l^\mu }T_l^\mu (X),\|X\|n1.$$ (5.4.5) * scale covariance: $$𝒰_\lambda T_J(X;𝐦)𝒰_\lambda ^1=\lambda ^{\omega _J}T_J(\lambda X;\lambda ^1𝐦),\|X\|n1.$$ (5.4.6) (ii) Now we can construct the expressions $`D_J(X;𝐦)`$ according to the formula (3.2.9) from Subsection 3.2 such that we have the well known properties of causality, Poincaré covariance and unitarity. We consider now a causal splitting of the type (3.2.15) such that we preserve Poincaré covariance and the order of singularity. The chronological products can be obtained from the formula (3.2.18), but we still have some freedom in the choice of the splitting which will shall use in the following. It can be proved as in that we have, instead of the relation (5.4.5) a somewhat weaker form, namely: $$d_QA(X;𝐦)=i\underset{l=1}{\overset{n}{}}\frac{}{x_l^\mu }A_l^\mu (X;𝐦)+P(X;𝐦),\|X\|=n$$ (5.4.7) where the expressions $`A(X;𝐦)`$ and $`A_l^\mu (X;𝐦)`$ are constructed from the expressions $`A_J(X)`$ according to the prescriptions (3.1.19) and (5.4.2). In the right hand side $`P(X;𝐦)`$ is a Wick polynomial (called anomaly) of the following structure: $$P(X)=\underset{J}{}\left[p_J()\delta (X)\right]:T_{j_1}(x_1)\mathrm{}T_{j_n}(x_n):$$ (5.4.8) with $`p_J`$ polynomials in the derivatives with the maximal degree restricted by $$deg(p_J)+\omega _J5,J.$$ (5.4.9) If we argue like in theorem 4.2.1 then we can see that from the induction hypothesis we have the following scaling behaviour of the chronological products $`T_J(X),\|X\|=n`$: $$𝒰_\lambda T_J(X;𝐦)𝒰_\lambda ^1=\lambda ^{\omega _J}T_J(\lambda X;\lambda ^1𝐦)+P_{J;k;𝐦;\lambda }(X)T_k(x_n;𝐦)$$ (5.4.10) for some quasi-local distribution $`P_{J;k;𝐦;\lambda }(X)`$ having an expression of the form (3.3.4). Moreover, these distributions have a coboundary structure (because we are in the case when we have massive fields in the theory) so one can redefine the chronological products $`T_J(X)`$ such that one gets rid of the expression $`P_{J;k;𝐦;\lambda }(X)`$. This implies, redefinitions for the causal splitting (3.2.15), in particular redefinitions for the distributions $`A(X;𝐦)`$ and $`A_l^\mu (X;𝐦)`$. It is clear that in this way we will not affect the general structure of the equation (5.4.7), that is we eventually modify the anomaly $`P`$ without spoiling the Poincaré covariance and the order of singularity of the splitting. Moreover, we will have in this way, instead of (5.4.10), the relation (5.4.6) for $`\|X\|=n`$ also. Because we can prove from the induction hypothesis that we have $$𝒰_\lambda A_J^{}(X;𝐦)𝒰_\lambda ^1=\lambda ^{\omega _J}A_J^{}(\lambda X;\lambda ^1𝐦)$$ (5.4.11) for $`\|X\|=n`$, we obtain that in this case we have also $$𝒰_\lambda A_J(X;𝐦)𝒰_\lambda ^1=\lambda ^{\omega _J}A_J(\lambda X;\lambda ^1𝐦)$$ (5.4.12) and in particular: $$𝒰_\lambda A(X;𝐦)𝒰_\lambda ^1=\lambda ^{4n}A(\lambda X;\lambda ^1𝐦),𝒰_\lambda A_l^\mu (X;𝐦)𝒰_\lambda ^1=\lambda ^{4n}A_l^\mu (\lambda X;\lambda ^1𝐦).$$ (5.4.13) This result can be obtained combining with the gauge invariance condition given by the equation (5.4.7) in the following way: we apply to $`𝒰_\lambda \mathrm{}𝒰_\lambda ^1`$ to (5.4.7), we use the lemma 5.3.4 and then the previous relations. The result is the following identity verified by the anomaly: $$𝒰_\lambda P(X;𝐦)𝒰_\lambda ^1=\lambda ^{4n+1}P(\lambda X;\lambda ^1𝐦).$$ (5.4.14) In other words, the anomaly must have the canonical dimension equal to $`4n+1`$. By “integrations by parts” (see ) we can exhibit the anomaly as follows: $$P(X)=i\underset{l=1}{\overset{n}{}}\frac{}{x_l^\rho }N_l^\rho (X)+P^{}(X)$$ (5.4.15) where $`P^{}(X)`$ is of the following form: $$P^{}(X)=\delta (X)𝒫(x_n)$$ (5.4.16) with $`𝒫(x)`$ a Wick polynomial in one variable. So, by redefining the expressions $`A_l^\mu (X)`$ we can take the anomaly of the form $$P(X)=\delta (X)𝒫(x_n).$$ (5.4.17) It is obvious that the “integration by parts” process will not affect the properties of the anomaly that we have already obtained. In consequence, the Wick polynomial $`𝒫(x)`$ will verify the following restrictions: * $`SL(2,)`$-covariance: $$U_{a,A}𝒫(x)U_{a,A}^1=𝒫(\delta (A)x+a),(a,A)inSL(2,).$$ (5.4.18) * Ghost numbers restrictions: $$gh(𝒫(x))=1.$$ (5.4.19) * Scale covariance: $$𝒰_\lambda 𝒫(x;𝐦)𝒰_\lambda ^1=\lambda ^5𝒫(\lambda X;\lambda ^1𝐦).$$ (5.4.20) If we use the generic structure $$𝒫(x;𝐦)=\underset{j}{}c_j(𝐦)T_j(x;𝐦)$$ (5.4.21) with $`c_j(𝐦)`$ some mass-dependent constants, in the last equation we immediately get: $$c_j(\lambda 𝐦)=\lambda ^{5\omega _j}c_j(𝐦).$$ (5.4.22) If we now construct the chronological products $`T(X),T_l^\mu (X)`$ one can fix in a standard way (see ) the properties of symmetrization and unitarity without spoiling the other relations we have already obtained. In the same one can fix PCT-covariance. In the end, we will have a relation of the type: $$d_QT(X;𝐦)=i\underset{l=1}{\overset{n}{}}\frac{}{x_l^\mu }T_l^\mu (X;𝐦)+P(X;𝐦),\|X\|=n$$ (5.4.23) where the anomaly $`P(X,𝐦)`$ has the form (5.4.17) and verifies all the preceding restrictions: Poincaré covariance, ghost number restriction, scale covariance. Moreover, it obviously verifies gauge invariance: $$d_QP(X;𝐦)=i\underset{l=1}{\overset{n}{}}\frac{}{x_l^\mu }P_l^\mu (X;𝐦)$$ (5.4.24) for some Wick polynomials $`P_l^\mu (X;𝐦)`$ and also PCT-covariance: $$U_{PCT}𝒫(x)U_{PCT}^1=(1)^n𝒫(x)$$ (5.4.25) and unitarity: $$𝒫(x)^{}(1)^n𝒫(x).$$ (5.4.26) (iii) The list of possible anomalies can be written now as in . We only remark that the restrictions imposed above do not lead to the conclusion that there are no anomalies in order $`n`$. In fact, a number of hard anomalies remain such as: $$𝒫_1=c_{abcde}^1\underset{m_a=m_b=m_c=m_d=m_e=0}{}u_a:\mathrm{\Phi }_b\mathrm{\Phi }_c\mathrm{\Phi }_d\mathrm{\Phi }_e:$$ (5.4.27) $$𝒫_2=c_{abc}^2\underset{m_a=m_b=m_c=0}{}u_a:^\mu \mathrm{\Phi }_b_\mu \mathrm{\Phi }_c:$$ (5.4.28) $$𝒫_3=c_{abc}^3\epsilon _{\mu \nu \rho \sigma }u_a:F_b^{\mu \nu }F_c^{\sigma \rho }:$$ (5.4.29) where $$F_a^{\mu \nu }^\mu A_a^\nu ^\nu A_a^\mu .$$ (5.4.30) $$𝒫_4=\underset{m_a=m_b=0}{}[:\overline{\psi }_A(K_{ab})_{AB}\psi _B::\overline{\psi }_A(K_{ab}^{})_{AB}\gamma _5\psi _B:]u_a\mathrm{\Phi }_b.$$ (5.4.31) One can show that from unitarity (or PCT-covariance) that we have $$c_{\mathrm{}}=(1)^nc_{\mathrm{}},K_{ab}^{}=(1)^nK_{ab},(K_{ab}^{})^{}=(1)^nK_{ab}^{}.$$ (5.4.32) The list of hard anomalies is larger: all the anomalies appearing in the second and in the third order of perturbation theory (see and ) should appear. ## 6 Conclusions The expression (5.4.29) is the famous Adler-Bardeen-Bell-Jackiw anomaly (ABBJ). So, we see that the various symmetries of the standard model (including scale covariance) are not sufficient to prove the anomalies are absent in higher orders of the perturbation theory if they are absent in orders $`n=1,2,3`$ (at least in Epstein-Glaser approach). In fact, if a certain type of anomaly is present in low orders of perturbation theory, this means that the corresponding expression is not in conflict with the various symmetries of the model. Then it is hard to imagine why such a conflict would appear in higher orders of perturbation theory. Such a result would be possible in our formalism only if in the equation (5.4.22) the number $`n`$ (the order of the perturbation theory) would survive. To obtain the cancelation of anomalies in all orders in our formalism a more refined formula for the distribution splitting seems to be needed. There appears to be a contradiction between our result and the analysis from (see also and ) where it is showed that the ABBJ anomaly can appear only in the order $`n=3`$. The discrepancy can be explained if one admits that in these references one works with the interaction fields. One know that one can construct such fields from the chronological products $`T_J(X)`$ as formal series ( see formula (76) of ) of the type $$\mathrm{\Phi }(x)=\frac{i^n}{n!}𝑑y_1\mathrm{}𝑑y_nR(y_1,\mathrm{},y_n;x)g(y_1)\mathrm{}g(y_n)$$ (6.0.1) where $`R`$ are the retarded products. If we perform formally the adiabatic limit we have $$\mathrm{\Phi }(x)=\underset{n=0}{\overset{\mathrm{}}{}}g^n\varphi _n(x)$$ (6.0.2) where $`\varphi _0`$ is the free field and $`g`$ is the coupling constant. Now supposes that a formula of the following type is valid: $$𝒰_\lambda \mathrm{\Phi }(x)𝒰_\lambda ^1=\lambda ^{\omega \gamma (g)}\mathrm{\Phi }(\lambda x),$$ (6.0.3) where $`\omega `$ is the canonical dimension of the field $`\varphi `$ and $$\gamma (g)=\underset{n=1}{\overset{\mathrm{}}{}}\gamma _ng^n$$ (6.0.4) is a formal series in the coupling constant called anomalous dimension. One can write $$\lambda ^{\gamma (g)}=e^{\gamma (g)ln(\lambda )}$$ (6.0.5) perform the Taylor expansion in $`g`$ and substitute in the preceding formula. Then one finds out by regrouping the terms that we have $$𝒰_\lambda \varphi _n(x)𝒰_\lambda ^1=\lambda ^\omega \underset{p+q=n}{}\underset{m_1,\mathrm{},m_p}{}\frac{(ln\lambda )^{m_1+\mathrm{}m_p}}{p!}\gamma _{m_1}\mathrm{}\gamma _{m_p}\mathrm{\Phi }_q(\lambda x).$$ (6.0.6) This relation should be interpreted as a relation on the retarded products $`R`$. In fact such a relation is compatible with our analysis in the scaling limit, where all momenta are very large. In this region it is plausible to assume that all masses of the theory are zero, so we can apply the result of Subsection 4.3 which leads to a formula having the structure (6.0.6). It is an interesting problem to make this analysis completely rigourous in the framework of Epstein-Glaser formalism. Acknowledgement: The author wishes to thank prof. K. Fredenhagen and dr. M. Dütsch for many discussions and critical remarks.
warning/0004/quant-ph0004052.html	ar5iv	text	# Quantum to Classical Transition from the Cosmic Background Radiation ## I Introduction In the last decade several proposals to modify the standard Hamiltonian dynamics, ranging from master equations to stochastic quantum mechanics, have been advanced to try to set up an unified description for microscopic and macroscopic physical phenomena. In the pioneer work by Ghirardi, Rimini, and Weber , quantum mechanics with spontaneous localization (QMSL), the state vector collapse, leading from quantum to classical dynamics results from the instantaneous action of a spontaneous random hitting process. Such a Poisson process is described by a “localization” operator, a gaussian function acting on each microscopic constituent of any system. The localization operator carries two free parameters; a mean frequency $`\lambda `$ and a localization width $`\alpha ^{1/2}`$, understood as new constants of nature (the spontaneous localization is argued to be a fundamental physical process). Through these basic assumptions the QMSL consists in an explicit model allowing an unified description for microscopic and macroscopic systems. It forbids the occurrence of linear superposition of states localized in far away spatial regions and induces a dynamics that agree with the predictions of classical mechanics. Pursuing the program of the QMSL model, Diosi presented an interesting connection between the original GRW hitting process and a modified Schrödinger equation. Another significant achievement concerning a dynamical reduction model, a stochastic equation for physical ensemble, was reported by Gisin . Next, Pearle described the QMSL model through an Itô stochastic differential equation. Basically, Pearle replaced the Poisson process of instantaneous hits in GRW model by a Markov process described as a stochastic modification of the Schrödinger equation, so that a continuous evolution of the state vector was accomplished. By considering a specific choice of the operators defining the Markov process (expressed in terms of creation and annihilation operators), Ghirardi, Pearle, and Rimini have described the mechanism known as continuous spontaneous localization (CSL) of systems of identical particles (the QMSL model has consistency only in the case of systems of distinguishable particles). Other investigations dealing with dynamical reduction models have recently been considered , among them it is worth to mention the model for intrinsic decoherence proposed by Milburn . While in the GRWP model the addition of stochastic terms in the Schrödinger evolution automatically destroys the quantum coherence of the physical properties of the system that attain a macroscopic level, the modification of the Liouville equation proposed by Milburn destroys the coherence even at microscopic level. In the CSL model the Itô stochastic equation for the evolution of the state vector reads $$d\|\psi =\left(\frac{i}{\mathrm{}}Hdt+dh\frac{1}{2}\overline{(dh)^2}\right)\|\psi ,$$ (1) where $`dh`$ is a linear self-adjoint operator, whose random fluctuation may increase or decrease the norm of the state vector. Using Itô formula (with the notation $`\|d\psi d\|\psi `$), $$d\psi ^2=\psi \|d\psi +d\psi \|\psi +\overline{d\psi \|d\psi },$$ (2) it is easy to see that Eq. (1) does not conserve the norm of $`\|\psi `$. Thus, the introduction of a norm conserving nonlinear process is mandatory. This process, whose random operator depends on the state vector, reads $$d\|\varphi =\left(\frac{i}{\mathrm{}}Hdt+dh_\varphi \frac{1}{2}\overline{(dh_\varphi )^2}\right)\|\varphi .$$ (3) Now, it is necessary to distinguish between raw ( Eq. (1)) and physical ( Eq. (3)) ensembles of state vectors to correctly understand the effect of the non-Hamiltonian terms. To this end a precept is adopted, namely, that the square norm of each (unnormalized) state vector represents the weight associated with that (normalized) state vector in the ensemble coming from the Itô stochastic equation . This precept is a generalization of the GRW assumption that the frequency of hits is proportional to the squared norm of the state vector. Therefore, in the GRW prescription the quantum theory prediction about the associated probabilities in a measurement process is recovered. By considering such a precept for the physical ensemble, the linearity of the raw equation and the Markov nature of the Itô stochastic process leads to the physical stochastic differential equation for the $`N`$-particle state vector $$d\|\mathrm{\Psi }_N=(\frac{i}{\mathrm{}}Hdt+𝐙.d𝐁\frac{1}{2}\gamma 𝐙^{}.𝐙dt)\|\mathrm{\Psi }_N,$$ (4) where $`𝐙\left\{Z_i\right\}`$ are operators on the Hilbert space of the system and the set of random operators $`𝐁\left\{B_i\right\}`$ is characterized through a real Wiener process, satisfying the following means and correlations over ensemble $$\overline{dB_i}=0,\overline{dB_idB_j}=\gamma \delta _{ij}dt.$$ (5) The statistical operator $`\rho _N=\overline{\|\mathrm{\Psi }_N\mathrm{\Psi }_N\|}`$ of the physical ensemble and its evolution equation are directly obtained from Eq. (4); using the Itô calculus in evaluating $`d\rho _N/dt`$ one gets $$\frac{d\rho _N}{dt}=\frac{i}{\mathrm{}}[H,\rho _N]+\gamma 𝐙\rho _N.𝐙^{}\frac{\gamma }{2}\{𝐙^{}.𝐙,\rho _N\},$$ (6) which is exactly the Lindblad form for the generator of a quantum dynamical semigroup. In the present work our main concern is to achieve the decoupling between the state vector dynamics of the center-of-mass (CM) and internal motion of a system of particles. In the GRWP model this decoupling results from a hypothesis of spontaneous localization of the system’s wave function due to a fundamental stochastic hitting process on the particles, which induces an increase of total mean energy of the Universe claimed to be the origin of the Cosmic Background Radiation (CBR). Contrarily to this argument, in the present work we assume the point of view of standard cosmology: the nowadays CBR is a clue that the Universe began its expansion from a Big Bang . This assumption is introduced with the purpose to avoid the unconventional increase of the total mean energy of the Universe. Formally, we hypothesize that the state vector, the Hamiltonian $`H`$ and operators $`𝐙`$, $`𝐙^{}`$ in Eq.(4) represent both, the system of particle and CBR radiation; the set of random functions $`\{B_i\}`$ describes the intervention of the CBR on the system and substitute the spontaneous localization process. Instead of elaborating on the formal microscopic problem of the interaction of a system with a reservoir , we assume ad-hoc that the evolution of the system of particles, under the influence of the CBR, is described by an Itô equation having stochastic coupling parameters. Therefore, in the present conservative continuous reduction model (the total energy of system plus CBR is conserved) we argue that: 1) the increase or decrease of the system’s mean energy is attributed to the CBR; 2) the positional space is not privileged with respect to the momentum space, as required when the localization operator is involved; 3) we do not claim for an additional assumption to decouple the collective and internal motion, namely the width parameter $`\alpha ^{1/2}10^5cm`$ in the CSL model; 4) as above-mentioned, more admissible results are obtained for decoherence times, while in the CSL model the value $`10^7s`$ obtained for a system of particles to undergo from quantum to classical dynamics seems to be too large (as well as the localization width $`\alpha ^{1/2}10^5cm`$ also seems too large when considering typical atomic distances about $`10^8cm`$, or even superposition of the center-of-mass coordinate different by more than about $`\alpha ^{1/2}`$ ), and finally, 5) instead of the two free parameters required in the GRWP model ($`\alpha ^{1/2}`$ and the mean frequency $`\lambda `$), the random function describing the interaction between the system and the CBR carries just a single strength parameter with dimension of inverse of time. In fact, the coupling constant of the CBR photons to the $`N`$-particle system, as the strength parameter in the GRWP model, defines the inverse of a characteristic time which is associated to the net effect of the random pseudo-“potential” $`dh`$ . Also, as in the GRWP model, our strength parameter is such small that nothing changes in the Hamiltonian dynamics of a single particle even in the case in which it has an extended wave function . Finally, we mention that Joos and Zeh , have previously argued that scattering of photons even at a relatively low temperature can induce the localization of the wave packet of a macroscopic system. So, their treatment, based on a master equation proposed by Wigner , suggests that the intergalactic cold CBR cannot simply be neglected . The model here presented goes exactly on this point, i.e., we consider the process of random scattering of the CBR photons by a system of particles as responsible for the superselection rules and the micro to macro transition of its dynamical description. In this way, despite inducing the superselection rules the CBR also induces the mechanism of separating the center-of-mass (CM) coordinate from the internal motion. Besides, we present a brief cosmological analysis of our results, discussing the roles played by both the CBR temperature and the number of particles of the system, in its way from quantum to classical dynamics, as the universe evolved from a hot to a cold state. In Section II we briefly review the GRWP model presenting its main achievements. In Section III we construct our model: beginning from an Itô stochastic equation we derive a pre-master equation for a system of $`N`$ particles and the CBR; tracing over the CBR degrees of freedom we obtain a master equation for the system of particles only and In Section IV we show that structurally it shows exactly the Lindblad form. In Section V we estimate the coupling parameter and in Section VI we estimate the decoherence time for the system of particles. In Section VII we show that at low temperature limit our master equation and the GRWP Itô equation are equivalent, thus this last one is a particular situation of the former; these equations allow the decoupling of the state vector dynamics into two separate equations, one for the CM and the other for the internal motion. In section VIII we calculate the entropy and analyze the problem of selection of a preferred basis. Finally, in Section IX we present a summary and conclusions. ## II The Ghirardi-Rimini-Weber-Pearle Model As explained in the introduction, in CSL model the random operator $`dh`$ contains in its definition the length parameter $`\alpha ^{1/2}`$ and a strength parameter $`\zeta `$ which is related to the mean hitting frequency $`\lambda `$. In this section we present a brief review of the CSL model as a class of Markov processes in Hilbert space . We will consider a system of $`N`$ identical particles so that the localization operator must involve globally the whole set of particles in order to preserve the symmetry properties of the wave function . For this purpose let us consider the creation and annihilation field operators $`a^{}(𝐪,s)`$, $`a(𝐪,s)`$ of a particle at the point $`𝐪`$ in some reference frame with spin component $`s`$, satisfying the canonical commutation or anticommutation relations. From these operators a locally averaged number density operator is defined as $$N(𝐱)=(\frac{\alpha }{2\pi })^{3/2}\underset{s}{}d^3𝐪\mathrm{exp}\left[\frac{1}{2}\alpha (𝐪𝐱)^2\right]a^{}(𝐪,s)a(𝐪,s).$$ (7) The operator $`N(𝐱)`$ is self adjoint and its commutator for different values of $`𝐱`$ vanishes. The total number operator is defined as $`N=d^3𝐱N(𝐱)`$, and the symmetrized (antisymmetrized) states containing $`n`$ particles at the indicated positions, $$𝐪,s=𝒩a^{}(𝐪_1,s_1)\mathrm{}a^{}(𝐪_n,s_n)0,$$ (8) constitutes the normalized common eigenvectors related to the eigenvalue equation $`N(𝐱)𝐪,s=n_𝐱𝐪,s`$, with $$n_𝐱=(\frac{\alpha }{2\pi })^{3/2}\underset{i=1}{\overset{N}{}}\mathrm{exp}\left[\frac{1}{2}\alpha (𝐱𝐪_i)^2\right].$$ (9) Applying the stochastic process established by Eq. (4) to a system of identical particles and considering the locally averaged density operator defined by Eq. (7), one gets the physical stochastic nonlinear differential equation for the state vector as $$d\psi _N=\left[iHdt+d^3𝐱N(𝐱)𝑑B(𝐱)\frac{1}{2}\zeta d^3𝐱N^2(𝐱)𝑑t\right]\psi _N.$$ (10) where the Wiener process $`B(𝐱)`$ satisfies $`\overline{dB(𝐱)}`$ $`=`$ $`0,`$ (12) $`\overline{dB(𝐱)dB(𝐲)}`$ $`=`$ $`\zeta \delta ^3(𝐱𝐲)dt.`$ (13) From Eq. (10) the evolution equation of the $`N`$-particle statistical operator obtained from Itô calculus reads $$\frac{\rho _N}{t}=i[H,\rho _N]+\zeta d^3𝐱N(𝐱)\rho _NN(𝐱)\frac{1}{2}\zeta \{d^3𝐱N^2(𝐱),\rho _N\}.$$ (14) and it can be checked that taking $`\lambda =\zeta (\alpha /4\pi )^{3/2}`$, Eq. (14) reduces to the correspondent equation for a single particle considered in the QMSL model. To discuss the physical implications of the modified dynamical equation (10), the separation of the CM motion will be made. If $`𝐐`$ is the CM coordinate of the system and $`\stackrel{~}{𝐪}_i`$ its internal coordinates (measured from the CM of the particles), one can define the particle coordinates as $$𝐪_i=𝐐+\stackrel{~}{𝐪}_i(\left\{𝐫_𝐢\right\}),$$ (15) where $`\left\{𝐫_𝐢\right\}`$ represents a set of $`3N3`$ independent variables. In the GRWP model the set $`\left\{𝐫_𝐢\right\}`$ does not contain macroscopic variables. As a consequence, assuming that the Hamiltonian can be written as $`H=H_Q+H_{r_i}`$, we consider the wave function $`\varphi (𝐪,s)`$ $`=`$ $`\mathrm{\Psi }(𝐐)\chi (𝐫_𝐢,s),`$ (17) $`\chi (𝐫_𝐢,s)`$ $`=`$ $`\left({\displaystyle \genfrac{}{}{0pt}{}{A}{B}}\right)\mathrm{\Delta }(𝐫_𝐢,s),`$ (18) where the symbol $`\left(\genfrac{}{}{0pt}{}{A}{B}\right)`$ specify the symmetrization or antisymmetrization of the internal coordinate wave function. Under the assumption that the length parameter $`\alpha ^{1/2}`$ is such that the internal wave function $`\mathrm{\Delta }(𝐫_𝐢,s)`$ is sharply peaked around the value $`𝐫_{𝐢}^{}{}_{0}{}^{}`$ of $`𝐫`$ (with respect to $`\alpha ^{1/2}`$ ), the action of the operator $`N(𝐱)`$ on the wave function (17) turns out to be $$N(𝐱)\mathrm{\Psi }(𝐐)\chi (𝐫_𝐢,s)=F(𝐐𝐱)\mathrm{\Psi }(𝐐)\chi (𝐫_𝐢,s),$$ (19) with $$F(𝐐𝐱)=\underset{i}{}(\frac{\alpha }{2\pi })^{3/2}\mathrm{exp}\left\{\frac{1}{2}\alpha \left[𝐐+\stackrel{~}{𝐪}_i(𝐫_0)𝐱\right]^2\right\}.$$ (20) Therefore, the operator $`N(𝐱)`$ acts only on $`\mathrm{\Psi }`$ and the separately normalized wave functions $`\mathrm{\Psi }`$ and $`\chi `$ satisfy the equations $`d\mathrm{\Psi }`$ $`=`$ $`\left[iH_Qdt+{\displaystyle d^3𝐱F(𝐐𝐱)𝑑B(𝐱)}{\displaystyle \frac{1}{2}}\zeta {\displaystyle d^3𝐱F^2(𝐐𝐱)𝑑t}\right]\mathrm{\Psi },`$ (22) $`d\chi `$ $`=`$ $`iH_{r_i}\chi dt.`$ (23) By assuming a large enough length parameter and an internal wave function which is independent of the macroscopic variables, the internal motion decouples as in the absence of the stochastic terms in Eq. (10). From this fact, the reduction rates which are characteristic of the GRWP theory together with the position and momentum spreading can be obtained. In particular, in the positional representation of Eq. (14), it is possible to verify with the help of the macroscopic density approximation and the sharp scanning approximation , that the macroscopic frequency associated to the system of identical particles is $$\mathrm{\Gamma }=\zeta D_0n_{out}.$$ (24) Here a homogeneous macroscopic body of density $`D_0`$ was considered and $`n_{out}`$ is the number of particles of the body at position $`𝐐^{}`$ that do not lie in the volume occupied by the body at position $`𝐐^{\prime \prime }`$. In the case of distinguishable particles, one gets the direct result $$\lambda _{CM}=n\lambda ,$$ (25) $`n`$ being the total number of particles, so that for a typical macroscopic number $`n10^{23}`$, one obtains $`\lambda _{CM}10^7s`$, as mentioned above. The position and momentum spreading obtained from the approximations leading to Eq. (24), are written as $`Q_i^2`$ $`=`$ $`Q_i^2_s+\zeta \delta _i{\displaystyle \frac{\mathrm{}^2}{6M^2}}t^3,`$ (27) $`P_i^2`$ $`=`$ $`P_i^2_s+{\displaystyle \frac{1}{2}}\zeta \delta _i\mathrm{}^2t,`$ (28) where the suffix $`s`$ indicates the Schrödinger evolution, and $$\delta _i=d^3𝐲\left(\frac{F(𝐲)}{y_i}\right)^2.$$ (29) Now, using the macroscopic density approximation applied to the identical particles system, Eq. (20) is modified to $$F(𝐐𝐱)=d^3\stackrel{~}{𝐲}D(\stackrel{~}{𝐲})(\frac{\alpha }{2\pi })^{3/2}\mathrm{exp}\left[\frac{1}{2}\alpha (𝐐+\stackrel{~}{𝐲}𝐱)^2\right],$$ (30) where $`D(𝐲)`$ is the number of particles per unit volume in the neighborhood of the point $`𝐲=𝐐+\stackrel{~}{𝐲}`$. The evaluation of the factor $`\delta _i`$ for the case of a homogeneous macroscopic box containing the $`N`$ particles, through the Eq. (30) gives the result $$\delta _i=(\alpha /\pi )^{1/2}D_0^2S_i,$$ (31) where $`S_i`$ is the transversal section of the macroscopic box. From Eq. (28) it is evident that the momentum variance implies that the CM energy increases per unit time as $$\frac{\mathrm{\Delta }E}{t}=\frac{\zeta \delta _i\mathrm{}^2}{M}10^{32}(gcms^1)S_icm^2,$$ (32) with the GRWP choice $`\alpha ^{1/2}10^5cm`$ together with $`D_010^{24}cm^3`$. From the requirement that the macroscopic frequency associated to the system of identical particles Eq. (24) is exactly the same as for distinguishable particles Eq. (25), GRWP have chosen $`\zeta 10^{30}cm^3s^1`$. ## III Decoherence from the Cosmic Background Radiation Our approach uses the stochastic dynamical equation (4), where we identify the continuous component (in frequency space) of the operator responsible for the interaction of the $`N`$-particle system to the CBR as $$𝐙(\mathrm{\Omega })\underset{k=1}{\overset{N}{}}\left(A(\mathrm{\Omega })𝐚_k^{}+A^{}(\mathrm{\Omega })𝐚_k\right),𝐚_k=(a_{k,x},a_{k,y},a_{k,z}).$$ (33) where $$𝐚_k=\frac{1}{\sqrt{2\mathrm{}m\omega }}\left(m\omega 𝐪_k+i𝐩_k\right),$$ (34) and $`𝐚_k^{}`$ is its hermitian conjugate $`\left([a_{k,i},a_{k^{},j}^{}]=\delta _{k,k^{}}\delta _{i,j},i=x,y,z\right)`$, $`𝐪_k`$ and $`𝐩_k`$ are respectively position and momentum operators of the $`k^{th}`$ particle of mass $`m`$. $`\mathrm{}\omega `$ is a characteristic energy of the system of particles associated to the quantum fluctuation of the CM. The operators $`A^{}(\mathrm{\Omega }),A(\mathrm{\Omega })`$ stand for the creation and annihilation of a quantum of energy $`\mathrm{}\mathrm{\Omega }`$ from the environment. The coupling parameter is defined by the continuous stochastic Wiener process $`𝐁(\mathrm{\Omega })`$ satisfying $`\overline{d𝐁(\mathrm{\Omega })}`$ $`=`$ $`0,`$ (36) $`\overline{dB_i(\mathrm{\Omega })dB_j(\mathrm{\Omega }^{^{}})}`$ $`=`$ $`\gamma (\mathrm{\Omega })\delta _{i,j}\delta (\mathrm{\Omega }\mathrm{\Omega }^{^{}})dt,`$ (37) with $`\gamma (\mathrm{\Omega })=\mathrm{\Lambda }\mathrm{\Gamma }(\mathrm{\Omega })`$ accounting for a strength parameter $`\mathrm{\Lambda }`$ and a frequency distribution function $`\mathrm{\Gamma }(\mathrm{\Omega })`$. Note that $`\mathrm{\Gamma }(\mathrm{\Omega })`$ refers to the effective frequency distribution of the CBR photons which interact with the system of particles at energy around $`\mathrm{}\omega `$. We next consider that the system of particles and CBR interacting almost resonantly with Lorentzian spectrum $$\mathrm{\Gamma }(\mathrm{\Omega })=\frac{1}{\pi }\frac{\tau _c}{\tau _c^2(\mathrm{\Omega }\omega )^2+1}.$$ (38) In view of Eq. (38) it follows from the Fourier transform of Eq. (37) that $$\overline{dB_i(t)dB_j(t^{})}=\frac{\mathrm{\Lambda }}{2\pi }e^{i\omega (tt^{^{}})}e^{(tt^{^{}})/\tau _c}dt,$$ (39) where the correlation time $`\tau _c`$ defines the memory time over which the stochastic function changes appreciably. From Eq. (39) we conclude that when considering $`\tau _c`$ extremely short, i.e., much less than all other times of interest (evolution of the particle system) so that in a good approximation $`\overline{dB_i(t)dB_j(t^{})}\delta (tt^{^{}})dt`$, the system is Markovian. Through Eqs. (36) and (37) the physical stochastic differential equation (4) reads $`d\|\mathrm{\Psi }_{N+CBR}`$ $`=`$ $`\{{\displaystyle \frac{i}{\mathrm{}}}H_{N+CBR}dt+{\displaystyle }d\mathrm{\Omega }{\displaystyle \underset{k=1}{\overset{N}{}}}(A(\mathrm{\Omega })𝐚_k^{}+A^{}(\mathrm{\Omega })𝐚_k)d𝐁(\mathrm{\Omega })`$ (41) $`{\displaystyle \frac{\mathrm{\Lambda }}{2}}{\displaystyle }d\mathrm{\Omega }\mathrm{\Gamma }(\mathrm{\Omega })\left[{\displaystyle \underset{k=1}{\overset{N}{}}}(A(\mathrm{\Omega })𝐚_k^{}+A^{}(\mathrm{\Omega })𝐚_k)\right]^2dt\}\|\mathrm{\Psi }_{N+CBR}.`$ It must be emphasized that Eq. (41) describes the evolution of the state vector of system of $`N`$ particles and CBR differently from the stochastic differential equation in the CSL model. The Hamiltonian $`H_{N+CBR}`$ in this equation describes the free evolution of the system of particles and CBR, while the two remaining terms account for the stochastic interaction between the CBR and partices. By defining both, the Wiener process $`d𝐁`$ and the operator $`𝐙`$ depending on the CBR frequency space, the positional space will not be anymore privileged with respect to the momentum space, as occurs in the CSL model. We now proceed to the separation of the CM motion of the modified dynamical equation (41). The substitution of the operators $`𝐚_k^{},𝐚_k`$ as position and momentum operators $`𝐩_k,𝐪_k`$, permits us to express Eq. (41) in terms of the CM coordinates $`𝐐=\frac{1}{N}_k𝐪_k`$ and $`𝐏=_k𝐩_k`$ as $`d\|\mathrm{\Psi }_{N+CBR}`$ $`=`$ $`\{{\displaystyle \frac{i}{\mathrm{}}}H_{N+CBR}dt+{\displaystyle }d\mathrm{\Omega }(A(\mathrm{\Omega })𝐗^{}+A^{}(\mathrm{\Omega })𝐗)d𝐁(\mathrm{\Omega })`$ (42) $``$ $`{\displaystyle \frac{\mathrm{\Lambda }}{2}}{\displaystyle }d\mathrm{\Omega }\mathrm{\Gamma }(\mathrm{\Omega })(A(\mathrm{\Omega })𝐗^{}+A^{}(\mathrm{\Omega })𝐗)^2dt\}\|\mathrm{\Psi }_{N+CBR}.`$ (43) where the operator $`𝐗`$ accounting for the macroscopic object reads $$𝐗=\frac{1}{\sqrt{2\mathrm{}m\omega }}\left(Nm\omega 𝐐+i𝐏\right),$$ (44) while $`𝐗^{}`$ is its hermitian conjugate. These operators satisfy the commutation relation $`[X_i,X_j^{}]=N\delta _{i,j}\widehat{1}`$. As mentioned earlier the coupling constant of the interaction between the CBR and the system of particles defines a characteristic time $`\mathrm{\Lambda }^1`$ which is associated to the net effect of the random pseudo-“potential” described by the last two terms on the right-hand side of Eq. (43). As the stochastic operator in Eq. (43) automatically acts only on the joint wave vector of the CM degree of freedom and the CBR $`\|\mathrm{\Psi }_{CM+CBR}`$, the separately normalized state vectors $`\|\mathrm{\Psi }_{CM+CBR}`$ and $`\|\varphi _{\{𝐫_i\}}`$, the latter for the internal degrees of freedom, satisfy the equations $`d\|\mathrm{\Psi }_{CM+CBR}`$ $`=`$ $`[{\displaystyle \frac{i}{\mathrm{}}}H_{CM+CBR}dt+{\displaystyle }d\mathrm{\Omega }(A(\mathrm{\Omega })𝐗^{}+A^{}(\mathrm{\Omega })𝐗)d𝐁(\mathrm{\Omega })`$ (47) $`{\displaystyle \frac{\mathrm{\Lambda }}{2}}{\displaystyle }d\mathrm{\Omega }\mathrm{\Gamma }(\mathrm{\Omega })(A(\mathrm{\Omega })𝐗^{}+A^{}(\mathrm{\Omega })𝐗)^2dt]\|\mathrm{\Psi }_{CM+CBR},`$ $`d\|\varphi _{\{𝐫_i\}}`$ $`=`$ $`iH_{_{\{𝐫_i\}}}\|\varphi _{\{𝐫_i\}}dt.`$ (48) It should be noted that the above Eqs. (47) and (48), differently from those in the CSL model (Eqs. (22) and (23)), involve also the CBR degrees of freedom. As will be shown later, the present approach in the low temperature limit allows obtaining separately the normalized wave functions for the system of particles, $`\|\mathrm{\Psi }_{CM}`$ and $`\|\varphi _{\{𝐫_i\}}`$, satisfying equations similar to those in the CSL. Next, from Eqs. (6) and (43) the statistical operator $`\rho _{N+CBR}=\overline{\|\mathrm{\Psi }_{N+CBR}\mathrm{\Psi }_{N+CBR}\|}`$ reads $`{\displaystyle \frac{d\rho _{N+CBR}}{dt}}`$ $`=`$ $`{\displaystyle \frac{i}{\mathrm{}}}[H_{N+CBR},\rho _{N+CBR}]{\displaystyle \frac{\mathrm{\Lambda }}{2}}{\displaystyle }d\mathrm{\Omega }\mathrm{\Gamma }(\mathrm{\Omega })[\{(A(\mathrm{\Omega })𝐗^{}+A^{}(\mathrm{\Omega })𝐗)^2,\rho _{N+CBR}\}`$ (50) $`2(A(\mathrm{\Omega })𝐗^{}+A^{}(\mathrm{\Omega })𝐗)\rho _{N+CBR}(A(\mathrm{\Omega })𝐗^{}+A^{}(\mathrm{\Omega })𝐗)],`$ which is a pre-master equation in that it contains operators from both, the $`N`$ particles system and the CBR, allowing to calculate correlations between operators of system of particles and CBR. However, since we only have at our disposal the statistical properties of the CBR field, the obvious procedure is to trace over the CBR degrees of freedom, considered thermalized at temperature $`T`$, which leads to the reduced density operator of the system of particles only, containing the average number of photons of the CBR as a parameter. Back to Eqs. (39), when the following assumptions are met: i) a short correlation time $`\tau _c`$ ($`\mathrm{\Lambda }^1`$), leading to the Markovian approximation; ii) the interaction between the system of particles and CBR is sufficiently small (exactly the purpose at hand), the density operator of the global system can be written as $`\rho _{N+CBR}(t)=\rho _N(t)\rho _{CBR}(t)+\rho _{correl}(t)`$, where the correlation term $`\rho _{correl}`$ can be neglected . By considering the thermalized CBR density operator $`\rho _{CBR}=\mathrm{exp}\left(\beta H_{CBR}(A^{},A)\right)/\mathrm{Tr}\left[\mathrm{exp}\left(\beta H_{CBR}(A^{},A)\right)\right]`$, with $`\beta =k_BT`$, $`k_B`$ being the Boltzmann’s constant and $`T`$ the CBR temperature, we find the master equation for the N-particle system $`{\displaystyle \frac{d\rho _N}{dt}}`$ $`=`$ $`{\displaystyle \frac{i}{\mathrm{}}}[H_N,\rho _N]{\displaystyle \frac{\mathrm{\Lambda }}{2}}{\displaystyle }d(\mathrm{\Omega })\mathrm{\Gamma }(\mathrm{\Omega })\{[𝐗^{}.,𝐗\rho _N]+[\rho _N𝐗^{}.,𝐗]`$ (52) $`+n_\mathrm{\Omega }([𝐗^{},[𝐗,\rho _N]]+[𝐗,[𝐗^{},\rho _N]])\},`$ where $`\rho _N`$ is the reduced density operator of the N-particle system and $`n_\mathrm{\Omega }=1/\left(\mathrm{exp}(\beta \mathrm{}\mathrm{\Omega })1\right)`$ is the thermal averaged photon number. As time goes on, it is expected that the stochastic coupling induces the $`N`$-particle system to a thermal equilibrium with the CBR. By evaluating the rate of energy change between the system and the CBR we shall estimate the strength parameter $`\mathrm{\Lambda }`$ and improve our understanding about the nature of this stochastic coupling. In order to estimate the energy mean-value let us consider the mean value of a generic dynamical variable $`𝒱`$ whose equation of motion is obtained through Eq.(52) as $`{\displaystyle \frac{d𝒱}{dt}}`$ $`=`$ $`{\displaystyle \frac{i}{\mathrm{}}}\mathrm{tr}\left([𝒱,H_N]\rho _N\right){\displaystyle \frac{\mathrm{\Lambda }}{2}}{\displaystyle }d(\mathrm{\Omega })\mathrm{\Gamma }(\mathrm{\Omega })\mathrm{Tr}\{[[𝒱,𝐗^{}]𝐗+𝐗^{}[𝐗,𝒱]`$ (54) $`+n_\mathrm{\Omega }([[𝒱,𝐗^{}],𝐗]+[[𝒱,𝐗],𝐗^{}])]\rho _N\},`$ By applying Eq.(54) to the position and momentum variables consecutively, we observe that not only the pure Schrödinger evolution is modified but also the results from the CSL model, such that the equations of motion become $`{\displaystyle \frac{d𝐐}{dt}}`$ $`=`$ $`{\displaystyle \frac{1}{M}}𝐏{\displaystyle \frac{1}{2}}N\mathrm{\Lambda }𝐐,`$ (56) $`{\displaystyle \frac{d𝐏}{dt}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}N\mathrm{\Lambda }𝐏.`$ (57) These equations lead to the results $`𝐏_t=\mathrm{exp}\left(\frac{1}{2}N\mathrm{\Lambda }t\right)𝐏_s`$ and $`𝐐_t=\mathrm{exp}\left(\frac{1}{2}N\mathrm{\Lambda }t\right)𝐐_s`$, where the subscript $`s`$ indicates the pure Schrödinger evolution: $`𝐏_s=𝐏_{t=0}`$ and $`𝐐_s=𝐐_0+𝐏_{t=0}/Mt`$. For $`𝒱=𝐐^2,𝐐𝐏+𝐏𝐐`$ and $`𝐏^2`$ successively, the equations of motion for the mean values become, respectively, $`{\displaystyle \frac{d𝐐^2}{dt}}`$ $`=`$ $`{\displaystyle \frac{1}{M}}𝐐𝐏+𝐏𝐐N\mathrm{\Lambda }𝐐^2+{\displaystyle \frac{3\mathrm{}\mathrm{\Lambda }}{2m\omega }}{\displaystyle 𝑑\mathrm{\Omega }\mathrm{\Gamma }(\mathrm{\Omega })\left(1+2n_\mathrm{\Omega }\right)},`$ (59) $`{\displaystyle \frac{d𝐐𝐏+𝐏𝐐}{dt}}`$ $`=`$ $`{\displaystyle \frac{2}{M}}𝐏^2N\mathrm{\Lambda }𝐐𝐏+𝐏𝐐,`$ (60) $`{\displaystyle \frac{d𝐏^2}{dt}}`$ $`=`$ $`N\mathrm{\Lambda }𝐏^2+{\displaystyle \frac{3N^2\mathrm{\Lambda }m\mathrm{}\omega }{2}}{\displaystyle 𝑑\mathrm{\Omega }\mathrm{\Gamma }(\mathrm{\Omega })\left(1+2n_\mathrm{\Omega }\right)},`$ (61) which differ from the pure Schrödinger evolution since $`\mathrm{\Lambda }0`$. ## IV The Master Equation and Itô Dynamics It will be useful to remind the conventional treatment of the problem of interaction of a N-particle system with the reservoir (R). Under the Hamiltonian $`H=H_N+H_R+V`$, $`V`$ being the interaction between both systems, the reduced density operator of the N-particle system, $`\rho _N(t)=\mathrm{Tr}_R\left[\rho _N(t)\right]`$, evolves, up to the second order in the interaction, according to the generalized master equation $$\frac{d\rho _N(t)}{dt}=\frac{i}{\mathrm{}}[H_N,\rho _N(t)]\frac{1}{\mathrm{}^2}\mathrm{Tr}_R_0^t[V,e^{iL_0(tt\stackrel{´}{})}[V,\rho _N(t\stackrel{´}{})\rho _R]]𝑑t\stackrel{´}{},$$ (62) where $`L_0()[H_N+H_R,]`$ is the Liouvillian operator of the free Hamiltonian. The second term in Eq.(62), acting as a source of noise for the system and also as a sink (or source) of energy, is responsible for the irreversibility of the process and the loss of coherence in $`\rho _N(t)`$. As such, the Itô calculus is justified when the stochastic terms are introduced into the Schrödinger equation. So, the CBR is responsible for the variation of the mean energy of the system and the increase of entropy. As shown by Isar et al. , choosing conveniently the interaction term $`V`$ it is possible to obtain Eq. (52) (the Lindblad form) from Eq. (62). It is worth noting that the master equation (52) can be written as $$\frac{d\rho _N(t)}{dt}=\frac{i}{\mathrm{}}[H_N,\rho _S(t)]+\underset{n=1}{\overset{2}{}}𝒮[c_n]\rho _N(t),$$ (63) where the superoperator $`𝒮[c_n]`$ is defined as $$𝒮[c_n]\rho _N=c_n.\rho _Nc_n^{}\frac{1}{2}\{c_n^{}.c_n,\rho _N\},$$ (64) with $`c_1=\left[\mathrm{\Lambda }𝑑\mathrm{\Omega }\mathrm{\Gamma }(\mathrm{\Omega })n_\mathrm{\Omega }\right]^{1/2}𝐗^{}`$ and $`c_2=\left[\mathrm{\Lambda }𝑑\mathrm{\Omega }\mathrm{\Gamma }(\mathrm{\Omega })\left(1+n_\mathrm{\Omega }\right)\right]^{1/2}𝐗`$. Written as in Eq. (63) our master equation resembles the Lindblad form for the decay of a mode of the eletromagnetic field inside a cavity . In summary, we have assumed ad-hoc that the evolution of the system of particles in its way from quantum to classical dynamics, under the influence of the CBR, is described by an Itô stochastic equation. However, here we showed that the usual master equation formalism can be viewed as a subdynamics of the Itô dynamics, without any need to use perturbation methods as is done in the conventional derivation. ## V The strength parameter Back to the equations of motion (III), their solutions are $`𝐐^2`$ $`=`$ $`𝐐^2_se^{N\mathrm{\Lambda }t}{\displaystyle \frac{3\mathrm{}\omega }{M}}\left[{\displaystyle \frac{t}{N\mathrm{\Lambda }}}\left(1{\displaystyle \frac{N\mathrm{\Lambda }t}{2}}\right)e^{N\mathrm{\Lambda }t}\left({\displaystyle \frac{1}{N^2\mathrm{\Lambda }^2}}+{\displaystyle \frac{1}{2\omega ^2}}\right)\left(1e^{N\mathrm{\Lambda }t}\right)\right],`$ (66) $`\{𝐐,𝐏\}`$ $`=`$ $`\{𝐐,𝐏\}_se^{N\mathrm{\Lambda }t}3\mathrm{}\omega \left[te^{N\mathrm{\Lambda }t}{\displaystyle \frac{1}{N\mathrm{\Lambda }}}\left(1e^{N\mathrm{\Lambda }t}\right)\right],`$ (67) $`𝐏^2`$ $`=`$ $`𝐏^2_se^{N\mathrm{\Lambda }t}+{\displaystyle \frac{3Nm\mathrm{}\omega }{2}}\left(1e^{N\mathrm{\Lambda }t}\right),`$ (68) where $`𝐐^2_s`$ $`=`$ $`𝐐^2_0+{\displaystyle \frac{1}{M}}\left(\{𝐐,𝐏\}_0t+{\displaystyle \frac{1}{M}}𝐏^2_0t^2\right),`$ (70) $`\{𝐐,𝐏\}_s`$ $`=`$ $`\{𝐐,𝐏\}_0+{\displaystyle \frac{2}{M}}𝐏^2_0t,`$ (71) $`𝐏^2_s`$ $`=`$ $`𝐏^2_0.`$ (72) The effect of the CBR temperature is present in the integral $`=𝑑\mathrm{\Omega }\mathrm{\Gamma }(\mathrm{\Omega })\left(1+2n_\mathrm{\Omega }\right)`$. It is worth noting that the time evolution of the operators in Eqs. (V) does not show the additive property with respect to the Schrödinger terms as obtained in the CSL model. As a consequence, Eq. (68) differs from the corresponding one in the CSL model, Eq. (28), because instead of the diffusion, inducing a steady increase of the mean value of the kinetic energy, the present model exhibits, asymptotically, thermalization due to the CBR, $$K=\left(K_sK_{eq}\right)e^{N\mathrm{\Lambda }t}+K_{eq},$$ (73) where the equilibrium kinetic energy reads $`K_{eq}=3\mathrm{}\omega /4.`$ So, $`\omega `$ is a characteristic frequency proportional to the thermalized mean kinetic energy of the CM. As mentioned above, in the CSL model the localization of a single particle of the system is sufficient to localize the whole system; as a consequence, the CM energy increases linearly with the “interaction” parameter $`N\mathrm{\Lambda }t`$. However, from Eq. (73) we conclude that the stochastic coupling accounts for a CM energy which grows or decays exponentially with $`N\mathrm{\Lambda }t`$, depending on the negative or positive value for $`K_sK_{eq}`$, respectively. In order to estimate the strength parameter $`\mathrm{\Lambda }`$, from Eq. (73) we assume that the relaxation time follows from relation $`(K_sK_{eq})e^{N\mathrm{\Lambda }\tau _R}K_{eq}`$, so that $$\mathrm{\Lambda }\frac{1}{N\tau _R}\mathrm{ln}\left(\frac{K_sK_{eq}}{K_{eq}}\right).$$ (74) For a system of $`N10^{23}`$ particles initially at room temperature the equipartition energy theorem gives a mean kinetic energy $`K_s10^9`$ergs. The integral $``$ accounting for the effect of the temperature of the CBR has been estimated in the Appendix for $`\beta \tau _c\mathrm{}`$, with $`\omega \tau _c1`$. The result $``$ $`1+2n_\omega `$, holds for both, low- and high-frequency regimes. So, we find for the equilibrium energy at low-frequency regime ($`\mathrm{}\omega k_BT`$, so that $`n_\omega k_BT/\mathrm{}\omega `$), $`K_{eq}k_BT10^{16}ergs`$. At high-frequency regime ($`\mathrm{}\omega k_BT`$), the equilibrium energy obeys $`K_{eq}k_BT`$. (We are referring to low- and high-frequency regimes since the nowadays CBR temperature, $`T3`$K is assumed). Taking $`K_{eq}`$ at low-frequency regime (in fact, due to the $`\mathrm{ln}`$ function, choosing $`K_{eq}`$ in low- or high- frequency will not change appreciably the value of $`\mathrm{\Lambda }`$), and the relaxation time $`\tau _R`$ of the order of the age of the Universe, about $`10^{16}s`$ (what seems to be reasonable when considering, as obtained below, such a small coupling of the system to the CBR), we get $$\mathrm{\Lambda }10^{38}s^1,$$ (75) a value to be compared with the above-mentioned coupling in the CSL model $`\zeta 10^{30}cm^3s^1`$. Thus, the parameter $`\mathrm{\Lambda }`$ is of the order of the upper limit of the excitation rate for nucleons estimated by Pearle and Squires , by comparison with neutrino-induced process. As already pointed out, such a value hardly affects the dynamics of a microscopic particle. ## VI Wave-packets reduction rates Back to Eq. (52), in the CM positional representation, the density matrix $`\rho _N(𝐐,𝐐^{})`$ evolves according to the differential equation $`{\displaystyle \frac{\rho _N(𝐐,𝐐^{},𝐭)}{t}}`$ $`=`$ $`\{{\displaystyle \frac{\mathrm{}}{2iM}}({\displaystyle \frac{^2}{𝐐^2}}{\displaystyle \frac{^2}{𝐐_{}^{}{}_{}{}^{2}}})𝒟[(𝐐𝐐^{})^2{\displaystyle \frac{\mathrm{}^2}{\left(M\omega \right)^2}}({\displaystyle \frac{}{𝐐}}+{\displaystyle \frac{}{𝐐^{}}})^2]`$ (76) $``$ $`{\displaystyle \frac{1}{2}}N\mathrm{\Lambda }[(𝐐{\displaystyle \frac{}{𝐐^{}}}+𝐐^{}{\displaystyle \frac{}{𝐐}})1]\}\rho _N(𝐐,𝐐^{},𝐭).`$ (77) The first term on the right-hand side comes from the commutator in Eq. (52), the terms multiplied by the diffusion constant $`𝒟=NM\mathrm{\Lambda }\omega \left(1+2n_\omega \right)/4\mathrm{}`$ (as well as the remaining terms, which are independent of temperature) account for the fluctuations (or random kicks) and for the energy changes due to the stochastic coupling, respectively. To analyze the wave-packet reduction rates we will not consider Eq. (77) in detail, since the effect of the second term on quantum superposition will be of much greater interest . For a brief estimation of the off-diagonal matrix elements of Eq. (77) will decay exponentially as $$𝐐\|\rho _S(t)\|𝐐^{}=e^{\zeta t}𝐐\|\rho _S(0)\|𝐐^{},$$ (78) where $`\zeta =𝒟(\mathrm{\Delta }𝐐)^2`$ and $`(\mathrm{\Delta }𝐐)^2=\left(𝐐𝐐^{}\right)^2`$. It follows from Eq. (78) that the quantum coherence of a macroscopic system will disappear on a decoherence time scale $$\tau _D\frac{1}{𝒟(\mathrm{\Delta }𝐐)^2}=\frac{1}{\left(1+2n_\omega \right)}\frac{\mathrm{}}{NM\mathrm{\Lambda }\omega (\mathrm{\Delta }𝐐)^2}.$$ (79) Analyzing Eq. (79) in terms of the CBR temperature, it is interesting to note that in the low-temperature limit (nowadays universe, $`T3K`$), i.e., $`n_\omega 0`$, the number of particles $`N`$ plays a crucial role in the decoherence process induced by the CBR. In the high-temperature limit, i.e, $`n_\omega \mathrm{}`$ (the early universe in the present model), we conclude that Eq. (79) leads from quantum to classical physics even a system composed by a small number of particles. This is a key result which help supporting the assumptions considered in the present model. Let us now estimate the decoherence time for both, a macroscopic and a microscopic object in nowadays Universe, i.e, $`T3K`$. In order to compare our results with that presented in literature, we consider the low-frequency regime, such that Eq. (79) reduces to $$\tau _D\frac{1}{𝒟(\mathrm{\Delta }𝐐)^2}=\frac{\mathrm{}^2}{2N\mathrm{\Lambda }Mk_BT(\mathrm{\Delta }𝐐)^2}.$$ (80) By considering a system of $`N`$ ($`10^{23})`$ hydrogen atoms with mass $`M1`$g and separation $`\mathrm{\Delta }𝐐1`$cm, quantum coherence would be destroyed in $`\tau _D10^{24}s`$. Such a value turns to be significantly smaller than the one obtained by GRWP, $`\lambda _{CM}10^7s`$ , Eq. (25), and comparable with that obtained through the linear response model of the Caldeira and Leggett (CL) , where, also at low-frequency regime, $`\tau _D/\tau _R\mathrm{}^2/2mk_BT(\mathrm{\Delta }𝐐)^2`$, $`\tau _R`$ being a relaxation time. For the above-mentioned system of $`N`$ atoms, and assuming $`\tau _R10^{16}s`$, as we have done to obtain $`\mathrm{\Lambda }`$, Eq. (74), we get from CL model $`\tau _D10^{23}s`$. So, Eq. (79), and consequently Eq. (80), arise from a theory that, despite assuring the essential character of the GRWP model, gives a more realistic value for the decoherence time of a macroscopic object. As far as a microscopic object is concerned, for example a single atom, $`m10^{24}g`$ on atomic scale $`\mathrm{\Delta }𝐐10^8cm`$, we observe the persistence of quantum coherence since $`\tau _D10^{41}s`$. Finally, we note that when considering a tiny Weber bar , $`\mathrm{\Delta }𝐐10^{19}m`$, at cryogenic temperatures, $`T10^3K`$, we also observe the persistence of quantum coherence from Eq. (79), as should be expected. Back to Eq. (78), when interpreting the exponential damping factor $`\zeta `$ by the light of the CSL model (Eqs. (24) and (25)), we conclude that the strength $`\mathrm{\Lambda }`$ plays the role of a microscopic frequency hitting parameter. ## VII The CM and Internal Motion By construction we assumed that the CBR acts only on the CM coordinates of the system of particles. Such assumption automatically decouples the dynamics of the collective and internal motions in the master equation (52). Next, we show that even the vector state dynamics for the CM and the internal motion decouple, as in the CSL model. Of course, our analysis will be restricted to the low temperature limit where, as obtained in Eq. (79), the macroscopic character of the system becomes really important due to the number of particles $`N`$. In this limit Eq. (52) simplifies to $$\frac{d\rho _N}{dt}=\frac{i}{\mathrm{}}[H_N,\rho _N]+\mathrm{\Lambda }𝐗\rho _N𝐗^{}\frac{\mathrm{\Lambda }}{2}\{𝐗^{}𝐗,\rho _N\}.$$ (81) The stochastic differential equation for the state vector of the system of particles which leads to Eq. (81) can be written as $$d\|\mathrm{\Psi }_N=(\frac{i}{\mathrm{}}H_Ndt+𝐗d𝐖\frac{\mathrm{\Lambda }}{2}𝐖𝐖^dt)\|\mathrm{\Psi }_S,$$ (82) now with the Wiener process $`\overline{dW_i}=0,\overline{dW_idW_j}=\mathrm{\Lambda }\delta _{ij}dt.`$ The assumption made in the CSL model, that the set $`\{𝐫_i\}`$ in Eq. (15) does not contain macroscopic variables, implies that the state vector for the macroscopic object factorizes as $`\mathrm{\Psi }_N(\{𝐪_k\})=\psi _{CM}(𝐐)\varphi _{int}(\{𝐫_i\})`$. The additional assumption that the CM motion is decoupled from the internal degrees of freedom means that the Hamiltonian must be written as a sum of two terms, $`H_N=H_{CM}+H_{int}`$ . Under these assumptions the Itô calculus, $`d\mathrm{\Psi }_N=d(\psi _{CM}\varphi _{int})=(d\psi _{CM})\varphi _{int}+\psi _{CM}(d\varphi _{int})+\overline{(d\psi _{CM}(d\varphi _{int})}`$, shows that the wave functions $`\psi _{CM}(𝐐)`$ and $`\varphi _{int}(𝐫_i)`$, similarly to Eqs. (22) and (23), satisfy equations $$d\|\psi _CM=\left(\frac{i}{\mathrm{}}H_{CM}dt+𝐗d𝐖\frac{\mathrm{\Lambda }}{2}𝐖^{}𝐖dt\right)\|\mathrm{\Psi }_{CM},$$ (84) $$d\|\varphi _{int}=\frac{i}{\mathrm{}}H_{int}\|\varphi _{int}dt.$$ (85) The stochastic terms do not affect the internal structure of the system of particles, i.e., nothing changes in the Schrödinger dynamics of microscopic particles. It is worth noting that in the CSL model the additional assumption of a large enough localization width parameter (besides of an internal wave function independent of macroscopic variables) is necessary to decouple the dynamics of $`\psi _{CM}`$ from $`\varphi _{int}`$. In fact, as shown in Ref. , a width parameter of order of atomic size leads to the breakdown of the translational symmetry of the system and the interaction between the CM and the relative coordinates (i.e., $`H=H_{CM}+H_{int}+V`$), has to be taken into account. However, in the present model, since we have assumed that the CBR acts only on the CM coordinates of the system of particles, no additional conjectures was requested about the random operator $`𝐙(\mathrm{\Omega })`$, Eq. (33), to achieve the remarkable result of the CM decoupling from the internal motion, as if the stochastic terms in Eq. (52) were absent. The operator $`𝐙(\mathrm{\Omega })`$ has thus the advantage of not needing additional conjectures about the width parameter of the localization process. ## VIII Decoherence and Entropy The decoherence process resulting from the interaction of the state vector for a macroscopic object with the CBR can be quantified by the rate of increase of either the linear or the statistical entropy. In terms of the density matrix, the statistical entropy, a measure of our ignorance, is defined as $`𝒮_s=\mathrm{Tr}\left(\rho \mathrm{ln}\left(\rho \right)\right)`$ (the subscript $`s`$ refers to statistical). This definition does not require that the system be in a thermal equilibrium state. Alternatively, a good measure of the loss of purity for states of an evolving open system is based on the increase of the linear entropy (subscript $`l`$) $$𝒮_l=\mathrm{Tr}\left(\rho \rho ^2\right).$$ (86) Next, we estimate the rate of increase of the linear entropy through the evolution of the density matrix given in the operator form by Eq. (77). Considering a weak strength parameter $`\left(\mathrm{\Lambda }0\right)`$ and the state vector remaining approximately pure $`\left(\mathrm{Tr}\rho ^21\right)`$, up to first order in $`\mathrm{\Lambda }`$ we obtain $$\dot{𝒮}_l=4𝒟\left(\left(\mathrm{\Delta }𝐐\right)^2+\frac{1}{\left(Nm\omega \right)^2}\left(\mathrm{\Delta }𝐏\right)^2\right),$$ (87) where $`\left(\mathrm{\Delta }𝐐\right)^2`$ and $`\left(\mathrm{\Delta }𝐏\right)^2,`$ obtained from Eqs. (66) - (72), stand for the variances of the position and momentum operators and can be rewritten as function of their initial values $`𝐐_0`$ and $`𝐏_0`$. In order to better understand the rate of increase of the linear entropy in Eq. (87), it is worth to compare it with that obtained by Zurek who used the linear response model of Caldeira and Leggett (in the high temperature limit). With the above approximations Zurek obtained $`\dot{𝒮}_l=4𝒟\left(\mathrm{\Delta }𝐐\right)^2`$ (for a single oscillator), so that the rate of increase of linear entropy (in quantum Brownian motion) is proportional to the dispersion in position coordinate only - the preferred observable singled out by the interaction Hamiltonian. In our approach, from Eq. (87) we observe that no preferred observable emerge from the dynamic equation (52) (the dispersion in momentum is also present), contrarily even to the CSL model where the position representation is taken from the outset as privileged. However, for a large number of particles ($`N1`$), Eq. (87) indicates that the dispersion in momentum is considerably smaller when compared with that in position which, in this situation, emerges as the preferred observable. In the weak-coupling limit we integrate Eq. (87) replacing the general evolution in Eq. (52) by the free von Neumann equation to obtain $`𝒮_l`$ $`=`$ $`4𝒟[(\left(\mathrm{\Delta }𝐐\right)^2_0+{\displaystyle \frac{1}{\left(Nm\omega \right)^2}}\left(\mathrm{\Delta }𝐏\right)^2_0)t+{\displaystyle \frac{1}{2M}}\mathrm{\Delta }\{𝐐,𝐏\}_0t^2`$ (88) $`+`$ $`{\displaystyle \frac{1}{3M^2}}\left(\mathrm{\Delta }𝐏\right)^2_0t^3],`$ (89) with $`\mathrm{\Delta }\{𝐐,𝐏\}\{𝐐,𝐏\}2𝐐𝐏`$. The dispersions appearing in the Eq.(89) are computed for the pure initial state. Back to the preferred basis problem, we remind that Zurek considered the free Heisenberg equations for the oscillator operators ($`P,Q`$) and obtained the linear entropy $`2𝒟\left(\left(\mathrm{\Delta }Q\right)^2_0+\frac{1}{\left(Nm\omega \right)^2}\left(\mathrm{\Delta }P\right)^2_0\right)`$ ($`N=1`$), averaged over one oscillator period. So, this result corresponds only to the coefficient for the linear time-dependence in Eq.(89), where additional terms as square and cubic time-dependent behavior also take place. Such a behavior indicates that, in spite of the large number of particles, for large times the momentum plays an important role in the problem of the preferred observable because we have considered the free motion of a $`N`$-particle system instead of a single harmonic oscillator. ## IX Summary and Conclusions In the GRWP model of continuous dynamical reduction of the state vector it is assumed that each microscopic constituent of a system of $`N`$ particles is subject to a sudden collapse due to a spontaneous random hitting process consisting in a localization of the wave function of the particle within an appropriate range . In what turns to be a remarkable result the localization of a single constituent of the system of particles is sufficient to localize the whole system. Such a spontaneous localization, considered as a fundamental physical process, induces a steady increase of the mean energy value of the physical system and so the increase in temperature per unit time of the universe. When taking into account that the age of the universe is about $`10^{16}s`$, the GRWP model leads to a total temperature increase from the beginning of the universe of $`10^3K`$, a value claimed to be comparable with the cosmic background radiation (CBR) of $`3`$K. In the present model for continuous dynamical reduction, also based in a stochastic differential equation describing a Markovian evolution of state vectors, the random hitting process in GRWP model is substituted by the intervention of the CBR. Such a strategy is intended to maintain (i) the principle of conservation of energy, and (ii) the claim that the Universe originated from the Big Bang leaving the CBR as a signature. In (i) the increase or decrease of the CM mean energy of the system of $`N`$ particles subject to a stochastic interaction with the CBR, which acts as a reservoir. In (ii), taking the opposite direction to the GRWP argument (which claim that the present temperature of the universe comes from the increase of the total energy arising from the random hitting process), we propose that the CBR temperature plays an important role in the reduction of the $`N`$-particle wave packet. So, we assumed, in agreement with the standard cosmology, that the Universe has originated from a hot state and cooling during its expansion, with decreasing mean photon energy. The Planck law for the thermal average boson number in CBR, indeed the best blackbody known, has recently been tested by the COBE satellite . The temperature of the CBR, decreasing as the mean photon energy decrease due to the cosmic expansion makes the mass of the system increasingly more important for the transition from quantum to classical description. On this basis one can argue that the quantum nature of the Universe becomes increasingly important as it is cooling. In fact, for the early Universe, the number of particles does not play a fundamental role in estimating the decoherence time, where higher temperatures (by itself) turn the system from micro to macro dynamics. However, as the Universe becomes cooler the number of particles becomes increasingly important. Moreover, the present model leads to realistic results for decoherence times. While in the GRWP model the value $`10^7s`$ obtained for a system of particles to go from micro to macro dynamics seems to be too large, the value $`10^{24}s`$ here obtained for a system of $`N`$ atoms in the low-frequency regime is comparable to the decoherence time obtained from the Caldeira-Leggett model. As mentioned, whereas the GRWP model requires a privileged positional space, in the present model, by construction, the stochastic operator acts on the CBR spectrum, carrying the same status for both, the position and the momentum space. The GRWP’s result - the wave function collapse of a single particles induces the collapse of the wave function of the whole system - was obtained exactly from the choice of the position as a preferred basis. The same result follows from our model without the choice of the position as a preferred basis. However, It has to be mentioned that in spite of attributing the same status for the position and the momentum space, when analyzing the entropy under the process of decoherence, the position coordinate still emerges as a preferred basis when considering a system with a large number of particles $`N`$. So, the preferred basis is directly related to the number of particles in the system. Another interesting feature is that we do not claim for an additional assumption to decouple collective from internal motion as the required large width parameter $`\alpha ^{1/2}10^5cm`$ in the GRWP model. The random operator $`𝐙(\mathrm{\Omega })`$ here assumed, besides being a more conventional choice since it is associated to a reservoir (CBR), leads to the advantage of decoupling the CM and internal motion without additional assumption beyond that usually assumed for a reservoir. The random operator describing the interaction between the system and the CBR carries only one parameter, the strength $`\mathrm{\Lambda }`$, instead of the two free parameters, as required in the GRWP model ($`\alpha ^{1/2}`$ and the mean frequency $`\lambda `$). In our model, the coupling of the CBR to the system, proportional to $`\mathrm{\Lambda }`$, corresponds to the random pseudo-“potential” $`dh`$ of the GRWP model. As well as the parameter $`\lambda `$ in GRWP model, our $`\mathrm{\Lambda }`$ is weak enough in the sense that it does not affect the dynamics of a unique particle, even in the case in which its wave function is spatially spread . Finally, we point out that the Itô equation is not derived from a physical picture of the background and associated scattering processes of the CBR by the system of particles. Instead of considering a particular interaction and choose some specific particle property sensible to the electric and magnetic field of the CBR, we approached the problem by modeling the interaction by a stochastic coupling, such that the dynamics could be described by an Itô equation. We have considered an effective strenght parameter $`\mathrm{\Lambda }`$ accounting for all kind of light-particle scattering processes. We also stress that our pre-master equation (34) (with respect to the particles) has still information on both, the system of particles and the CBR, since it contains operators of both subsystems. This approach is different from the usual one where for getting a master equation it is necessary to trace over the environment degrees of freedom, as is done in the theories of Joos and Zeh and Caldeira-Leggett or even in quantum optics. In our model it is possible to calculate correlations between observables of both subsystems. However, we have get rid of CBR degrees of freedom, Eq. (35), just because the available information on the CBR subsystem is sparse, consisting of the blackbody radiation distribution function at 3K. Thus the master equation (35) expressed in the CM positional representation, Eq. (49), incorporates the similar equations obtained in both theories, Joos and Zeh and Caldeira-Leggett. The main difference between the three approaches stem in the nature of the diffusion constant (DC): In Joos and Zeh the DC originates from the scattering of electromagnetic waves by small objects; in Caldeira-Leggett it comes from the fluctuations arising from energy dissipation of the system of interest to a thermal reservoir. In our model the DC originates from the stochastic interaction between $`N`$ particles of mass $`m`$ and the CBR at temperature $`T`$. ###### Acknowledgements. MCO and NGA thank FAPESP, São Paulo, Brazil, for total financial support. MHYM and SSM thank CNPq, Brazil, for partial financial support. The authors wish to thank Prof. R. J. Napolitano for helpful discussions. ## A Calculation of integral $``$ Due to the normalized Lorentzian spectrum (Eq. (38)), the integral accounting for the temperature of the CBR reads $`=1+2𝑑\mathrm{\Omega }\mathrm{\Gamma }(\mathrm{\Omega })n_\mathrm{\Omega }`$. Now, since the Planck’s distribution $`n_\mathrm{\Omega }`$ diverges when $`\mathrm{\Omega }`$ goes to zero, the same occurs to the remaining integral $`𝑑\mathrm{\Omega }\mathrm{\Gamma }(\mathrm{\Omega })n_\mathrm{\Omega }`$. However, as usual, we assume that the spectrum $`\mathrm{\Gamma }(\mathrm{\Omega })`$ has its maximum far away from zero in order to cancel the divergence coming from $`n_\mathrm{\Omega }`$. In what follows we are going to estimate under which conditions this approximation is valid. After the transformations $`\mathrm{\Omega }\tau _c=x`$ and $`\gamma =\mathrm{}/k_BT\tau _c`$, the remaining integral reads $$\frac{1}{\pi }_{\mathrm{}}^+\mathrm{}𝑑x\frac{1}{[x(\omega \tau +i)][x(\omega \tau _ci)]}\frac{1}{e^{\gamma x}1},$$ (A1) which can be solved in the complex space through Jordan’s lema, leading to the result $$2i\left\{\frac{1}{2i}\frac{1}{e^{\gamma (\omega \tau _c+i)}1}+\frac{1}{\gamma }\underset{n=0}{\overset{\mathrm{}}{}}\left(1\frac{1}{2}\delta _{n,0}\right)\frac{1}{[\omega \tau +i(1\frac{2\pi n}{\gamma })][\omega \tau _ci(1+\frac{2\pi n}{\gamma })]}\right\}.$$ (A2) It can be shown that the imaginary term coming from the above result is zero. Now, denoting $`\gamma =p/\xi `$, where the parameter $`p`$ is equal to $`\mathrm{}\omega /k_BT`$ whereas $`\xi =\omega \tau _c`$, the real term coming from (A2), reads $$\frac{\mathrm{cos}(\xi /p)e^\xi 1}{e^\xi [e^\xi 2\mathrm{cos}(\xi /p)]+1}8\pi \frac{p^3}{\xi ^2}\underset{n=1}{\overset{\mathrm{}}{}}\frac{n}{[1+p^2(2\pi np/\xi )^2]+(4\pi np^2/\xi )^2}.$$ (A3) For large $`n`$ the second term of (A3) reduces to $$\frac{\xi ^2}{p}\underset{n=1}{\overset{\mathrm{}}{}}\frac{1}{n^3}.$$ (A4) The analysis of the above result will be restricted to the condition $`\xi /p1`$, with $`\xi 1`$, under which the sum in (A4) can be disregarded (since even $`\xi ^2/p1)`$, and the first term in (A3) gives us $`1e^{\mathrm{}\omega /k_BT}1)`$, in a way that the Lorentzian distribution $`\mathrm{\Gamma }(\mathrm{\Omega })`$ acts practically as a delta function ($`\delta (\mathrm{\Omega }\omega )`$). In fact, the limit$`\xi 1`$, leads to the condition $`\omega \tau _c^1`$, so that the frequency can be taken far away from zero since, as above-discussed, we are considering an extremely short correlation time (Markovian approximation). Under such condition it is expected that the lorentzian function $`\mathrm{\Gamma }`$ acts indeed as a delta function, what means that the action of the reservoir over the system of particles is restricted to the oscillators whose frequencies is closely related to $`\omega `$. So, the problem of how far $`\omega `$ has to be from zero, in order to eliminate the divergence coming from Planck’s distribution when $`\omega 0,`$ depends exactly on the lorentzian height in its maximum. Moreover, the condition $`\xi /p1`$, with $`\xi 1`$, holds for both, the low- and high-frequency regimes. When $`\xi 1`$ (so that $`\omega \tau _c^1`$), we get the high-frequency regime $`\mathrm{}\omega k_BT`$, whereas for $`\xi 1`$ even the low-frequency regime is allowed. For the latter case we have to assure that $`0\omega \tau _c^1`$, not only to get rid of the divergence arising from $`n_{\omega \text{ }}`$, but also to hold the assumption of highly excited oscillations of the CBR leading to the Markovian approximation. Summarizing, under the conditions established above we get the result $`1+2n_\omega `$, which holds for high- and low-frequency regime.
warning/0004/hep-th0004177.html	ar5iv	text	# Regular (2+1)-dimensional black holes within non-linear Electrodynamics. ## Abstract Abstract:(2+1)-regular static black hole solutions with a nonlinear electric field are derived. The source to the Einstein equations is an energy momentum tensor of nonlinear electrodynamics, which satisfies the weak energy conditions and in the weak field limit becomes the (2+1)-Maxwell field tensor. The derived class of solutions is regular; the metric, curvature invariants and electric field are regular everywhere. The metric becomes, for a vanishing parameter, the (2+1)-static charged BTZ solution. A general procedure to derive solutions for the static BTZ (2+1)-spacetime, for any nonlinear Lagrangian depending on the electric field is formulated; for relevant electric fields one requires the fulfillment of the weak energy conditions. Keywords: 2+1 dimensions, Non-Linear black hole PACS numbers: 04.20.Jb, 97.60.Lf In general relativity the literature on regular black hole solutions is rather scarce . In (3+1)-gravity it is well known that electrovacuum asymptotically flat metrics endowed with timelike and spacelike symmetries do not allow for the existence of regular black hole solutions. Nevertheless, in the vacuum plus cosmological constant $`\mathrm{\Lambda }`$ case, the de-Sitter solution with positive cosmological constant is known to be a regular non-asymptotically flat solution (the scalar curvature is equal to $`4\mathrm{\Lambda }`$ and all the invariants of the conformal Weyl tensor are zero.) In order to be able to derive regular (black hole) gravitational–nonlinear electromagnetic fields one has to enlarge the class of electrodynamics to nonlinear ones . On the other hand in (2+1)-gravity, which is being intensively studied in these last years , in the vacuum case all solutions are locally Minkowski (the Riemann tensor is zero); the extension to the vacuum plus cosmological constant allows for the existence of the rotating anti de Sitter regular black hole (the scalar curvature and the Ricci square invariants are constants proportional to $`\mathrm{\Lambda }`$ and $`\mathrm{\Lambda }^2`$.) The static (2+1)-charged black hole with cosmological constant (the static charged BTZ solution) is singular (when $`r`$ goes to zero the curvature and the Ricci square invariants blow up). Similarly as in the (3+1)-gravity, one may search for regular solutions in (2+1)-gravity incorporating nonlinear electromagnetic fields to which one imposes the weak energy conditions in order to have physically plausible matter fields. One can look for regular solutions with nonlinear electromagnetic fields of the Born-Infeld type or electrodynamics of wider spectra. In this work, we are using electromagnetic Lagrangian $`L(F)`$ depending upon a single invariant $`F=1/4F^{ab}F_{ab}`$, which we demand in the weak field limit to be equal to the Maxwell Lagrangian $`L(F)F/4\pi `$, the corresponding energy momentum tensor has to fulfill the weak energy conditions: for any timelike vector $`u^a`$, $`u^au_a=1`$ (we are using signature – + +) one requires $`T_{ab}u^au^b0`$ and $`q_aq^a0`$, where $`q^a=T_b^au^b`$. The action of the (2+1)-Einstein theory coupled with nonlinear electrodynamics is given by $`S={\displaystyle \sqrt{g}\left(\frac{1}{16\pi }(R2\mathrm{\Lambda })+L(F)\right)d^3x},`$ (1) with the electromagnetic Lagrangian $`L(F)`$ unspecified explicitly at this stage. We are using units in which $`c=G=1`$. The ambiguity in the definition of the gravitational constant (there is not Newtonian gravitational limit in (2+1)-dimensions) allows us to maintain the factor $`1/16\pi `$ in the action to keep the parallelism with (3+1)-gravity. Varying this action with respect to gravitational field gives the Einstein equations $`G_{ab}+\mathrm{\Lambda }g_{ab}=8\pi T_{ab},`$ (2) $`T_{ab}=g_{ab}L(F)F_{ac}F_b^cL_{_{,F}},`$ (3) while the variation with respect to the electromagnetic potential $`A_a`$ entering in $`F_{ab}=A_{b,a}A_{a,b}`$, yields the electromagnetic field equations $`_a\left(F^{ab}L_{_{,F}}\right)=0,`$ (4) where $`L_{_{,F}}`$ stands for the derivative of $`L(F)`$ with respect to $`F`$. Concrete solutions to the dynamical equations above we present for the static metric $`ds^2=f(r)dt^2+{\displaystyle \frac{dr^2}{f(r)}}+r^2d\mathrm{\Omega }^2,`$ (5) where $`f(r)`$ is an unknown function of the variable r. We restrict the electric field to be $`F_{ab}=E(r)\left(\delta _a^t\delta _b^r\delta _a^r\delta _b^t\right).`$ (6) The invariant $`F`$ then is given by $`2F=E^2(r),`$ (7) thus the electric field can be expressed in term of the invariant $`F`$. Substituting (6) and (7) into the electromagnetic field equations (4) we arrive at $`E(r)L_{,F}={\displaystyle \frac{e}{r}},`$ (8) where $`e`$ is an integration constant. We choose $`e=q/4\pi `$ in order to obtain the Maxwell limit. Then we have $`E(r)L_{,F}={\displaystyle \frac{q}{4\pi r}}.`$ (9) Using now (7) we express the derivative $`L_F`$ as function of $`r`$, as follows $`L_{,r}={\displaystyle \frac{q}{4\pi r}}E_{,r}.`$ (10) We rewrite the Einstein’s equations equivalently as $`R_{ab}=8\pi \left(T_{ab}Tg_{ab}\right)+2\mathrm{\Lambda }g_{ab}.`$ (11) From (3) using (6) and (7) the trace becomes $`T=3L(F)+2E^2(r)L_{,F}.`$ (12) As it was above pointed out, the Lagrangian $`L(F)`$ must satisfy: (i) correspondence to Maxwell theory, i.e. $`L(F)L/4\pi `$, and (ii) the weak energy conditions: $`T_{ab}u^au^b0`$ and $`q_aq^a0`$, where $`q^a=T_b^au^b`$ for any timelike vector $`u^a`$; in our case the first inequality requires $`(L+E^2L_{,F})0,`$ (13) which can be stated equivalently as $`LEL_{,E}L{\displaystyle \frac{q}{4\pi r}}E.`$ (14) The norm of the energy flux $`q_a`$, occurs to be always less or equal to zero; for $`u^a`$ along the time coordinate, $`u^a=\delta _t^a/\sqrt{g_{tt}}`$, one has the inequality $`q_aq^a=(L+L_{,F}E^2)^20`$. Assume now that one were taking into account additionally the scalar magnetic field $`B:=F_{\varphi r}`$, then the Maxwell equations would be $`{\displaystyle \frac{d}{dr}}[rEL_{,F}]=0,{\displaystyle \frac{d}{dr}}{\displaystyle \frac{f}{r}}BL_{,F}=0.`$ (15) On the other hand, the Ricci tensor components, evaluated for the BTZ metric (5), would yield the following relation $`A:=R_{tt}+f^2R_{rr}=0,`$ (16) while the evaluation of the same relation using the electromagnetic energy-momentum would give $`A=8\pi L_{,F}({\displaystyle \frac{f}{r}}B)^2.`$ (17) Therefore, the scalar magnetic field should be equated to zero , $`B=0`$ , thus the only case to be treated is just the one with the electric field $`E`$. As far as the Einstein equations are concerned, the $`R_{tt}(=f^2R_{_{rr}})`$ and $`R_{_{\mathrm{\Omega }\mathrm{\Omega }}}`$ components yield respectively the equations $`f_{,rr}+{\displaystyle \frac{f_{,r}}{r}}=4\mathrm{\Lambda }+16\pi \left(2L(F)+E^2L_{,F}\right),`$ (18) $`f_{,r}=2\mathrm{\Lambda }r+16\pi r\left(L(F)+E^2L_{,F}\right).`$ (19) If one replaces $`f_{,r}`$ from (19) and its derivative $`f_{,rr}`$ into (18) one arrives, taking into account the equation (10), at an identity. Therefore one can forget the equation (18) and integrate the relevant Einstein equation (19): $`f(r)=M\mathrm{\Lambda }r^2`$ (20) $`+16\pi {\displaystyle r\left[L(F(r))+E^2L_{,F}\right]𝑑r}.`$ (21) Summarizing we have obtained a wide class of solutions, depending on a Lagrangian $`L(E)`$, given by: the metric $`ds^2=f(r)dt^2+{\displaystyle \frac{dr^2}{f(r)}}+r^2d\mathrm{\Omega }^2,`$ (22) the structural function $`f(r)=M(\mathrm{\Lambda }2C)r^2`$ (23) $`+4q{\displaystyle \left[r\frac{E_{,r}}{r}𝑑rE\right]𝑑r},`$ (24) which is obtained from (20) by using (10) and (7), where $`C`$ is a constant of integration, and the Lagrangian $`L(E)`$ is constrained to $`L_{,r}={\displaystyle \frac{q}{4\pi r}}E_{,r},`$ (25) We recall that the Lagrangian and the energy momentum tensor have to fulfill the conditions quoted above. We present now various particular examples: The static charged BTZ solution is characterized by the function $`g_{tt}=f=M+{\displaystyle \frac{r^2}{l^2}}2q^2lnr,`$ (26) the Lagrangian and the electric field $`L(E)={\displaystyle \frac{1}{8\pi }}E^2={\displaystyle \frac{1}{8\pi }}{\displaystyle \frac{q^2}{r^2}},E(r)={\displaystyle \frac{q}{r}},`$ (27) where $`C`$ has been equated to zero and $`\mathrm{\Lambda }=1/l^2`$. It is worthwhile to point out that the static charged BTZ black hole is singular at $`r=0`$. Other interesting example arises in the Born-Infeld electrodynamics – nonlinear charged (2+1)–black-hole . In this case the structural function is $`g_{tt}=f=M(\mathrm{\Lambda }2b^2)r^22b^2r\sqrt{r^2+q^2/b^2}`$ (28) $`2q^2ln(r+\sqrt{r^2+q^2/b^2}),`$ (29) and the Lagrangian and the electric field are given by $`L(F)={\displaystyle \frac{b^2}{4\pi }}\left(\sqrt{1+2{\displaystyle \frac{F}{b^2}}}1\right)=`$ (30) $`{\displaystyle \frac{b^2}{4\pi }}\left({\displaystyle \frac{r}{\sqrt{r^2+q^2/b^2}}}1\right),`$ (31) $`E(r)={\displaystyle \frac{q}{\sqrt{r^2+q^2/b^2}}},`$ (32) where $`b`$ is the Born-Infeld parameter, and $`C=b^2`$. This solution fulfills the weak energy conditions and it is singular at $`r=0`$. From the Ricci and Kretschmann scalars it follows that in this case there is a curvature singularity at $`r=0`$ . A new class of solution, which is regular everywhere, is given by the structural function of the form $`f(r)=M\mathrm{\Lambda }r^2q^2ln(r^2+a^2)`$ (33) where $`M`$, $`a`$, $`q`$ and $`\mathrm{\Lambda }`$ are free parameters. The Lagrangian and the electric field are given by $`L(r)={\displaystyle \frac{q^2}{8\pi }}{\displaystyle \frac{(r^2a^2)}{(r^2+a^2)^2}},`$ (34) $`E(r)={\displaystyle \frac{qr^3}{(r^2+a^2)^2}}.`$ (35) This Lagrangian requires to set $`C=0`$. The Lagrangian and the electric field satisfy the weak energy conditions (14). To express the Lagrangian in terms of $`F`$ or equivalently $`E`$, one has to write $`r`$ in terms of $`E`$ by solving the quartic equation for $`r(E)`$, this will give rise an explicit $`r`$ containing radicals of $`E`$, which introduced in $`L(r)`$, finally will bring $`L`$ as function of $`E`$. The expression $`L(E)`$ is not quite illuminating, thus we omit it here. To establish that this solution is regular one has to evaluate the curvature invariants . The non-vanishing curvature components, which occur to be regular at $`r=0`$, are given by: $`R_{0110}={\displaystyle \frac{q^2(a^2r^2)}{(r^2+a^2)^2}}+\mathrm{\Lambda },`$ (36) $`R_{0202}=f(r)\left({\displaystyle \frac{q^2r^2}{r^2+a^2}}+\mathrm{\Lambda }r^2\right),`$ (37) $`R_{1212}=f(r)^1\left({\displaystyle \frac{q^2r^2}{r^2+a^2}}+\mathrm{\Lambda }r^2\right),`$ (38) where 0,1,2 stand respectively for t,r and $`\mathrm{\Omega }`$. Evaluating the invariants $`R`$, and $`R_{ab}R^{ab}`$ one has $`R`$ $`=`$ $`{\displaystyle \frac{2q^2(r^2+3a^2)}{(r^2+a^2)^2}}+6\mathrm{\Lambda }`$ (39) $`R_{ab}R^{ab}`$ $`=`$ $`12\mathrm{\Lambda }^2+4q^4{\displaystyle \frac{r^4+2r^2a^2+3a^4}{(r^2+a^2)^4}}`$ (41) $`+{\displaystyle \frac{8\mathrm{\Lambda }q^2(3a^2+r^2)}{(r^2+a^2)^2}}.`$ Since the metric, the electric field and these invariants behave regularly for all values of $`r`$, we conclude that this solution is curvature regular everywhere. Nevertheless, for solutions without any horizon or black hole solutions with an inner and outer horizons, at $`r=0`$ a conical singularity may arise. At $`r=0`$ the function $`f(r)`$ becomes $`f(0)=Mq^2\mathrm{ln}(a^2)`$. Thus for $`M`$ positive, $`M>0`$, and $`a`$ in the range $`0<a<1`$, the value of $`f(0)`$ will be $`f(0)=M+q^2\mathrm{ln}(1/a)^2`$ , which will be positive, say $`f(0):=\beta ^2`$, if $`\mathrm{ln}(1/a)^2>M/q^2`$. In such a case, for $`0<\beta <1`$ the solutions will show angular deficit since the angular variable $`\mathrm{\Omega }`$, which originally runs $`0\mathrm{\Omega }<2\pi `$ will now run $`0\mathrm{\Omega }<2\beta \pi `$; the parameter $`a`$ can be expressed in terms of $`\beta `$, $`q`$ and $`M`$ as $`a^2=\mathrm{exp}[(\beta ^2+M)/q^2]`$. For $`\beta =1`$, there will be no angular deficit, the ratio of the perimeter of a small circle around $`r=0`$ to its radius, as this last tends to zero, will be $`2\pi `$. If one allows $`M`$ to be negative, $`M<0`$, and $`a`$ to take values in the interval $`0<a<1`$, then $`f(0)`$ will be always positive, in this case one can adopt the following parametrization: $`M=\beta ^2\mathrm{cos}^2\alpha `$, $`q^2\mathrm{ln}(1/a)^2=\beta ^2\mathrm{sin}^2\alpha `$, therefore $`f(0)=\beta ^2`$. One will have angular deficit if $`0<\beta <1`$ , and for $`\beta =1`$ the resulting (2+1) space-time will be free of singularities. Another possibility with positive $`f(0)=\beta ^2`$ arises for $`M<0`$, and $`a>1`$, $`f(0)`$ can be parameterized as $`M=\beta ^2\mathrm{cosh}^2\alpha `$, $`q^2\mathrm{ln}(1/a)^2=\beta ^2\mathrm{sinh}^2\alpha `$. Again the values taken by $`\beta `$ will govern the existence of angular deficit, for $`\beta =1`$ the solutions will be regular. If $`f(0)`$ is negative, $`f(0)=:\beta ^2`$, the character of the coordinates $`t`$ and $`r`$ changes, the coordinate $`t`$ becomes space–like, while $`r`$ is now time–like and one could think of the singularities, if any, as causal structure singularities because they could arise at the “time” $`r=0`$. In what follows we shall treat the parameter $`a`$ as a free one, having in mind the above restrictions to have solutions free of conical singularities. To establish that this solution represents a black hole, one has to demonstrate the existence of horizons, which require the vanishing of the $`g_{tt}`$ component, i.e., $`f(r)=0`$. The roots of this equation give the location of the horizons (inner and outer in our case). The roots – at most four – of the equation $`f(r)=0`$ can be expressed in terms of the Lambert $`W(r)`$ function $`r_{1,2,3,4}=\pm [exp({\displaystyle \frac{\mathrm{\Lambda }a^2M}{q^2}}`$ (42) $`LW[{\displaystyle \frac{\mathrm{\Lambda }}{q^2}}exp\left({\displaystyle \frac{\mathrm{\Lambda }a^2M}{q^2}}\right)])a^2]^{\frac{1}{2}}.`$ (43) There arise various cases which depend upon the values of the parameters: four real roots (two positive and two negative roots, the negative roots have to be ignored), two complex and two real roots, two complex and one real positive root (the extreme case), and four complex roots ( no black holes solutions.) Although this analytical expression for the Lambert function can be used in all calculations, (we recall that Lambert function fulfills the following equation $`ln(LW(x))+LW(x)=ln(x)`$), it occurs also useful to extract information from the graphical behavior of the our $`f(r)`$ (see figures). Analytically one can completely treat the extreme black hole case; for it, the derivative of $`f(r)`$ has to be zero, $`_r(f(r))=0`$ , at the $`r_{extr}`$, this gives $`r_{extr}=\sqrt{a^2{\displaystyle \frac{q^2}{\mathrm{\Lambda }}}}>0`$ (44) for $`\mathrm{\Lambda }<0`$. From this expression one concludes that the following inequality holds: $`a^2<q^2/\mathrm{\Lambda }`$. Entering now $`r_{extr}`$ into $`f(r)=0`$ one obtains a relation between the parameters involved, which can be solved explicitly for the mass–the extreme one– $`M_{extr}=a^2\mathrm{\Lambda }+q^2\left(1+ln\left[{\displaystyle \frac{\mathrm{\Lambda }}{q^2}}\right]\right),`$ (45) this $`M_{extr}`$ varies its values depending on the values given to the parameters $`a`$, $`q`$ and $`\mathrm{\Lambda }`$. We have an extreme black hole characterized by negative cosmological constant, $`\mathrm{\Lambda }<0`$, and positive extreme mass, $`M_{extr}>0`$, if the parameter $`a`$ is restricted by the inequality $`a^2<(q^2(1+ln(\mathrm{\Lambda }/q^2)))/\mathrm{\Lambda }`$. For other values of the mass $`M`$, one distinguishes the following branches: if $`M>M_{extr}`$ one has a black hole solution, and if $`M<M_{extr}`$ there are no horizons. In FIG. 1 we draw the graph of $`f(r)`$ which corresponds to regular solutions for a fixed values of $`M`$ and changing the values of the parameters $`\mathrm{\Lambda }`$, $`q`$. In FIG. 2 we draw the graph of $`f(r)`$ corresponding to solutions which exhibit a conical singularity at $`r=0`$, for $`f(0)=1/2`$, keeping $`M`$ fixed while $`\mathrm{\Lambda }`$ and $`q`$ change. If one were interested in the thermodynamics of the obtained solution one would to evaluate the temperature of the black hole, which is given in terms of its surface gravity by $`k__BT__H={\displaystyle \frac{\mathrm{}}{2\pi }}k.`$ (46) In general, for a spherically symmetric (and for circularly symmetric in (2+1)-dimensions) system the surface gravity can be computed via (for our signature) $`k=\underset{rr__+}{lim}\left[{\displaystyle \frac{1}{2}}{\displaystyle \frac{_rg_{tt}}{\sqrt{g_{tt}g_{rr}}}}\right],`$ (47) where $`r__+`$ is the outermost horizon. For our solution we have from (23), (46) and (47) that $`k__BT={\displaystyle \frac{\mathrm{}}{2\pi }}\left(\mathrm{\Lambda }r__+{\displaystyle \frac{q^2r__+}{r__+^2+a^2}}\right).`$ (48) Since in our case there is no an analytical expression of $`r__+`$ in terms of elementary functions, one can not give a parameter dependent expression of (48). It is easy to check that when $`q=0`$, $`T`$ in (48) reduces to the BTZ temperature. In the extreme case (44), the temperature vanishes in (48). The entropy can be trivially obtained using the entropy formula $`S=4\pi r__+`$. Other thermodynamic quantities such as heat capacity and chemical potential can be computed as in . To achieve the maximal extension of our regular black solutions one has to follow step by step the procedure presented in determining first the Kruskal-Szekeres coordinates, and to proceed further to draw the Penrose diagrams. Informative discussions with Jorge Zanelli, Ricardo Troncoso, Rodrigo Aros, and Eloy Ayón-Beato are gratefully acknowledged. This work was supported in part by FONDECYT-Chile 1990601, Dirección de Promoción y Desarrollo de la Universidad del Bío-Bío through Grant No 983105-1 (M.C.), FONDECYT-Chile 1980891, CONACYT-México 3692P-E9607, 32138E (A.G.) and in part by Dicyt de la Universidad de Santiago de Chile (M.C., A.G.).
warning/0004/math0004029.html	ar5iv	text	# Non-Archimedean intersection indices on projective spaces and the Bruhat-Tits building for 𝑃⁢𝐺⁢𝐿 ## 1 Introduction In this paper we interpret non-Archimedean local intersection numbers of linear cycles in $`^{n1}`$ with the combinatorial geometry of the Bruhat-Tits building associated to $`PGL(n)`$. The ultimate motivation behind these results is to associate to a projective space a differential geometric object playing the role of a model at infinity in the sense of Manin, who constructed in \[Ma\] such an object for curves. A first step in this direction is to look for a geometric interpretation of non-Archimedean intersection numbers which has an Archimedean analogue. It turns out that the Bruhat-Tits building for $`PGL`$ is a good candidate for such a geometric framework. The goal of the present paper is to express non-Archimedean intersection numbers in terms of the building. In another work (see \[We\]) we construct the desired “model at infinity” for projective spaces as an Archimedean analogue of the building. Besides, we use results of the present paper to derive parallel geometric formulas for Archimedean and non-Archimedean Arakelov intersection numbers. Let us now describe our main results. We denote by $`X`$ the Bruhat-Tits building associated to the group $`G=PGL(V)`$, where $`V`$ is an $`n`$-dimensional vector space over a non-Archimedean local field $`K`$ of characteristic $`0`$. The vertices in $`X`$ correspond to the homothety classes $`\{M\}`$ of $`R`$-lattices $`M`$ in $`V`$, where $`R`$ is the ring of integers in $`K`$. We fix a lattice $`M`$ in $`V`$, which induces a projective space $`(M)`$ over $`R`$, and consider $`d`$ linear cycles on $`(M)`$ intersecting properly in a cycle of dimension 0. These cycles are equal to projective spaces $`(N_i)`$ for split $`R`$-submodules $`N_i`$ of $`M`$. We put $`L_j=_{ij}N_i`$, and define $`F`$ as the following set of vertices in $`X`$: $$F=\{\{\pi ^{k_1}L_1+\mathrm{}+\pi ^{k_d}L_d\}:k_1,\mathrm{},k_d\}.$$ In Theorem 5.1 we express Serre’s intersection number $`<(N_1),\mathrm{},(N_d)>`$ of our linear cycles in $`(M)`$ as: $$<(N_1),\mathrm{},(N_d)>=\text{dist}(\{M\},F),$$ where dist is what we call the combinatorial distance function in $`X`$, i.e. the minimal length of a path consisting of 1-simplices connecting $`\{M\}`$ with a vertex in $`F`$. In the case of $`^1`$, this result specializes to a formula in \[Ma\]. Moreover we investigate the case of several linear cycles meeting properly in a cycle of higher dimension. In this case, of course, we no longer have an intersection number, but we can nevertheless describe the intersection cycle (see Theorem 5.2). It consists of one component coming from the generic fibre, which appears with multiplicity 1, and one component coming from the special fibre, appearing with multiplicity $`\text{dist}(\{M\},F)`$, where $`F`$ is defined in a similar way as in the previous result. Acknowledgements: I thank Ch. Deninger, G. Kings, K. Künnemann, E. Landvogt, Y.I. Manin, P. Schneider, E. de Shalit and M. Strauch for useful and inspiring discussions. I am also grateful to the Max-Planck-Institut für Mathematik in Bonn for financial support and the stimulating atmosphere during the early stages of this project. ## 2 Intersection theory In this section we will list the definitions and results from intersection theory which are needed later on, thereby fixing our notation. Let $`\mathrm{\Omega }`$ be a scheme, of finite type and separated over a regular ring. By $`Z^p(\mathrm{\Omega })`$ we denote the codimension $`p`$ cycles on $`\mathrm{\Omega }`$, i.e. the free abelian group on the set of integral (i.e. irreducible and reduced) closed subschemes of codimension $`p`$ . We write $`CH^p(\mathrm{\Omega })`$ for the quotient of $`Z^p(\mathrm{\Omega })`$ after the group generated by the principal cycles $`\text{div}(f)`$ for rational functions $`f0`$ on a codimension $`p1`$ integral closed subscheme, see \[Fu\], 1.3. If $`T\mathrm{\Omega }`$ is a closed subset, let $`Z_T^p(\mathrm{\Omega })`$ denote the free abelian group on the set of codimension $`p`$ integral closed subschemes of $`\mathrm{\Omega }`$, which are contained in $`T`$ and $`CH_T^p(\mathrm{\Omega })`$ the Chow group of cycles supported on $`T`$, i.e. $`Z_T^p(\mathrm{\Omega })`$ modulo the subgroup generated by all $`\text{div}(f)`$ for rational functions $`f0`$ on some codimension $`p1`$ integral closed subscheme of $`\mathrm{\Omega }`$ which is contained in $`T`$. If $`T=\mathrm{}`$, we put $`CH_T^p(\mathrm{\Omega })=0`$. Now we will briefly recall Serre’s intersection pairing. It is defined on any smooth scheme $`\mathrm{\Omega }`$ over a discrete valuation ring by \[Se\], V-32. Two closed, integral subschemes $`X`$ and $`Y`$ of $`\mathrm{\Omega }`$ meet properly if for every irreducible component $`W`$ of $`XY`$ we have $`\text{codim}(X)+\text{codim}(Y)=\text{codim}(W)`$. By \[Se\] we always have the inequality $`{}_{}{}^{\prime \prime }_{}^{\prime \prime }`$. If $`X`$ and $`Y`$ meet properly, then Serre defines an intersection index $`i_W(X,Y)`$ of $`X`$ and $`Y`$ along $`W`$ by higher Tor functors, see \[Se\], V-21. Let $`W`$ be an irreducible component of $`XY`$, let $`𝒪_{\mathrm{\Omega },w}`$ be the local ring at the generic point $`w`$ of $`W`$, and let $`p_X,p_Y`$ be the ideals in $`𝒪_{\mathrm{\Omega },w}`$ corresponding to $`X`$ and $`Y`$. If $`X`$ and $`Y`$ are locally Cohen-Macaulay, then $$i_W(X,Y)=l_{𝒪_{\mathrm{\Omega },w}}(𝒪_{\mathrm{\Omega },w}/(p_X+p_Y))$$ by \[Se\], p. V-20. We define the intersection cycle of properly intersecting $`X`$ and $`Y`$ by $`XY=_Wi_W(X,Y)W`$ where the sum runs over all irreducible components of $`XY`$. We can continue this product linearly to arbitrary cycles $`X`$ and $`Y`$ meeting properly on $`\mathrm{\Omega }`$, which means that any irreducible component of $`X`$ meets all the irreducible components of $`Y`$ properly. If $`X_1,\mathrm{},X_r`$ are $`r`$ closed integral subschemes meeting properly, i.e. so that every irreducible component $`W`$ of $`_iX_i`$ satisfies $`\text{codim}W=\text{codim}X_i`$, then the intersection of $`X_1,\mathrm{}X_r`$ is defined inductively: $`X_1\mathrm{}X_r=(\mathrm{}((X_1X_2)X_3)\mathrm{})X_r`$. Again we can extend this product linearly to arbitrary cycles $`X_1,\mathrm{},X_r`$ meeting properly. Gillet and Soulé have defined an arithmetic intersection pairing for arithmetic Chow groups in \[Gi-So\]. The local contributions of this pairing at the finite places (in the smooth case) can be described as follows (see \[Gi-So\], 4.5.1): Assume that $`\pi :\mathrm{\Omega }S=\text{Spec}R`$ is smooth, separated and of finite type over the discrete valuation ring $`R`$ and that $`\mathrm{\Omega }`$ is irreducible. Let $`Z_p(\mathrm{\Omega }/S)`$ be the free abelian group on the closed integral subschemes $`Y\mathrm{\Omega }`$ of relative dimension $`p`$ over $`S`$. Here the relative dimension $`dim_S(Y)`$ of $`Y`$ over $`S`$ is defined as $$dim_SY=\text{transcendence degree of }k(Y)\text{ over }k(T)\text{codim}_S(T),$$ where $`T`$ is the closure of $`\pi (Y)`$ in $`S`$ and $`k(Y)`$, $`k(T)`$ are the function fields. This relative dimension has the property that $$dim_S(\mathrm{\Omega })=dim_S(Y)+\text{codim}_\mathrm{\Omega }(Y)$$ for all closed integral subschemes $`Y`$ of $`\mathrm{\Omega }`$ (see \[Fu\], Lemma 20.1). Let $`CH_p(\mathrm{\Omega }/S)`$ be $`Z_p(\mathrm{\Omega }/S)`$ modulo rational equivalence. For all closed subschemes $`T\mathrm{\Omega }`$ we have $`CH_T^p(\mathrm{\Omega })=CH_{dp}(T/S)`$, where $`d`$ is the relative dimension of $`\mathrm{\Omega }`$ over $`S`$. For closed subschemes $`Y`$ and $`Z`$ we can define a pairing $$CH_Y^p(\mathrm{\Omega })\times CH_Z^q(\mathrm{\Omega })CH_{YZ}^{p+q}(\mathrm{\Omega })$$ as follows: It suffices to define a pairing $$CH_{dp}(Y/S)\times CH_{dq}(Z/S)CH_{dpq}(YZ/S).$$ Let $`VY`$ and $`WZ`$ be integral closed subschemes. Then we define a cycle in $`Y\times Z`$ as follows: $$VW=\{\begin{array}{cc}0,\hfill & \text{ if }V\text{ and }W\text{ are contained in the closed fibre},\hfill \\ [V\times _SW],\hfill & \text{ otherwise}\hfill \end{array}$$ By \[Fu\], Proposition 20.2, this induces a pairing $$CH_{dp}(Y/S)\times CH_{dq}(Z/S)CH_{2dpq}(Y\times _SZ/S).$$ Since $`\mathrm{\Omega }`$ is smooth over $`S`$, the diagonal embedding $`\mathrm{\Delta }:\mathrm{\Omega }\mathrm{\Omega }\times _S\mathrm{\Omega }`$ is a regular embedding, and we have Fulton’s Gysin map (see \[Fu\], §6 and §20): $$\mathrm{\Delta }^!:CH_{2dpq}(Y\times _SZ/S)CH_{dpq}(YZ/S).$$ Hence we get the desired pairing. If $`V`$ and $`W`$ meet properly, their image under this pairing coincides with the image of Serre’s intersection pairing in $`CH_{YZ}^{p+q}(\mathrm{\Omega })`$ by \[Fu\], 7.1.2 and 20.2.2. In particular, if $`Y`$ and $`Z`$ are irreducible with $`p=\text{codim}_\mathrm{\Omega }(Y)`$ and $`q=\text{codim}_\mathrm{\Omega }(Z)`$, the image of $`(Y,Z)`$ via $`CH_Y^p(\mathrm{\Omega })\times CH_Z^q(\mathrm{\Omega })CH_{YZ}^{p+q}(X)`$ yields an intersection class $$YZ=\mathrm{\Delta }^!(YZ)CH_{YZ}^{p+q}(\mathrm{\Omega }).$$ If one cycle is given by a Cartier divisor, say $`Y=D`$, then we have a different description of the image of $`DZ`$ in $`CH_Z^{p+q}(\mathrm{\Omega })`$ (see \[Fu\], 8.1.1, 20.2.1 and 6.1c): It is equal to the class of $`E`$, where $`E`$ is any Weil divisor on $`Z`$ induced by a Cartier divisor whose line bundle is isomorphic to $`j^{}𝒪(D)`$. Here $`j:Z\mathrm{\Omega }`$ is the embedding of $`Z`$ into $`\mathrm{\Omega }`$ and $`𝒪(D)`$ is the line bundle on $`\mathrm{\Omega }`$ corresponding to the class of $`D`$. In particular, if $`D`$ and $`D^{}`$ are linear equivalent divisors on $`\mathrm{\Omega }`$, the images of $`DZ`$ and $`D^{}Z`$ in $`CH_Z^{p+q}(\mathrm{\Omega })`$ coincide. If we have several irreducible cycles $`Y_1`$, $`Y_2,\mathrm{},Y_r`$ in $`\mathrm{\Omega }`$ of codimensions $`p_1`$, $`p_2,\mathrm{},p_r`$ we can define inductively an intersection class $`Y_1\mathrm{}Y_rCH_{Y_i}^{p_1+\mathrm{}+p_r}(\mathrm{\Omega })`$. Let $`k`$ be the residue field of $`R`$. We denote by $`\mathrm{deg}`$ the degree map for 0-cycles in the special fibre $`\mathrm{\Omega }_k`$ of $`\mathrm{\Omega }`$, i.e. for all $`z=n_PPZ^d(\mathrm{\Omega }_k)`$ we put $`\mathrm{deg}z=n_P[k(P):k]`$, where $`k(P)`$ is the residue field of $`P`$. Assume additionally that $`\mathrm{\Omega }`$ is proper over $`S`$, and let $`YZ^p(\mathrm{\Omega })`$ and $`ZZ^q(\mathrm{\Omega })`$ be two irreducible closed subschemes such that $`p+q=d+1`$ which intersect properly on the generic fibre of $`\mathrm{\Omega }`$. This means that their generic fibres are disjoint, so that $`YZ`$ is contained in the special fibre $`\mathrm{\Omega }_k`$ of $`\mathrm{\Omega }`$. Hence we can define an intersection number $$<Y,Z>=\mathrm{deg}(YZ),$$ where we take the degree of the image of $`YZCH_{YZ}^{d+1}(\mathrm{\Omega })`$ in $`CH^d(\mathrm{\Omega }_k)`$. Similarly, if $`Y_iZ^{p_i}(\mathrm{\Omega })`$ for $`i=1,\mathrm{},r`$ are prime cycles with $`p_i=d+1`$ which meet properly on the generic fibre, we put $`<Y_1,\mathrm{},Y_r>=\mathrm{deg}(Y_1\mathrm{}Y_r)`$. ## 3 Hyperplanes Throughout this paper we denote by $`K`$ a finite extension of $`_p`$, by $`R`$ its valuation ring and by $`k`$ the residue class field. Besides, $`v`$ is the valuation map, normalized so that it maps a prime element to $`1`$. We write $`q`$ for the cardinality of the residue class field, and we normalize the absolute value on $`K`$ so that $`\|x\|=q^{v(x)}`$. Besides, we fix an $`n`$-dimensional vector space $`V`$ over $`K`$. Let $`(V)=\text{Proj Sym}V^{}`$ be the projective space corresponding to $`V`$, where $`V^{}`$ is the linear dual of $`V`$. Every non-zero linear subspace $`W`$ of $`V`$ defines an integral (i.e. irreducible and reduced) closed subscheme $`(W)=\text{Proj Sym}W^{}(V)`$ of codimension $`ndimW`$. By “$`R`$-lattice in $`V`$” we always mean an $`R`$-lattice in $`V`$ of full rank. Every $`R`$-lattice $`M`$ in $`V`$ defines a model $`(M)=\text{Proj Sym}_R(M^{})`$ of $`(V)`$ over $`R`$, where $`M^{}`$ is the $`R`$-linear dual of $`M`$. If the lattices $`M`$ and $`N`$ differ by multiplication by some $`\lambda K^\times `$ then the corresponding isomorphism $`(M)\stackrel{}{}(N)`$ induces the identity on the generic fibre. Throughout this paper we call a submodule $`N`$ of $`M`$ split, if the exact sequence $`0NMM/N0`$ is split, i.e. if $`M/N`$ is free (or, equivalently, torsion free). Every split $`R`$-submodule $`N`$ of $`M`$ defines a closed subscheme $`(N)=\text{Proj Sym}N^{}(M)`$. ###### Lemma 3.1 For every split $`R`$-submodule $`N`$ of $`M`$, the closed subscheme $`(N)=\text{Proj Sym}N^{}`$ of $`(M)`$ is integral, and has codimension $`n\text{rk}N`$. Proof: This follows from dualizing the sequence $`0NMM/N0`$. $`\mathrm{}`$ The cycles in $`(M)`$ induced by split submodules are called linear, and linear cycles of codimension $`1`$ are called hyperplanes. The homogeneous prime ideal corresponding to the linear cycle $`(N)(M)`$ is generated by a base of $`(M/N)^{}`$ (regarded in $`M^{}`$). In particular, it is generated by homogeneous elements of degree one. Now fix a lattice $`M`$ in $`V`$ and an $`R`$-basis $`x_1,\mathrm{},x_n`$ of $`M`$. Let $`B`$ be a matrix in $`GL(n,R)`$ which we regard as an endomorphism of $`M`$ via our fixed basis. Then $`B`$ induces an automorphism (which we also denote by $`B`$) of $`(M)`$. The following lemma can be proven easily: ###### Lemma 3.2 If the hyperplane $`H`$ in $`(M)`$ is given by the linear homogeneous element $$f=\underset{j=1}{\overset{n}{}}a_jx_j^{}M^{},$$ where $`x_1^{},\mathrm{},x_n^{}`$ is the dual basis of $`x_1,\mathrm{},x_n`$, then $`B(H)`$ is given by the homogeneous element $`_{j=1}^nb_jx_j^{}`$ where $$\left(\begin{array}{c}b_1\\ \mathrm{}\\ b_n\end{array}\right)={}_{}{}^{t}B_{}^{1}\left(\begin{array}{c}a_1\\ \mathrm{}\\ a_n\end{array}\right).$$ Now we need an easy matrix lemma. We call a quadratic matrix a permutation matrix if it contains exactly one entry $`1`$ in every line and column, and if all other entries are equal to zero. ###### Lemma 3.3 Let $`A=(a_{ij})`$ be an $`(n\times n)`$-matrix over $`R`$. Then there exist elements $`C`$ and $`D`$ in $`GL(n,R)`$, where $`D`$ is a permutation matrix, such that the matrix $`CAD=(b_{ij})_{i,j}`$ is upper triangular with $$v(b_{11})\mathrm{}v(b_{nn})\text{ and}$$ $$v(b_{ii})v(b_{ij})\text{ for all }ij.$$ Proof: We move a coefficient with minimal valuation in the upper left corner and eliminate the other entries in the first column. This can be repeated until our matrix is upper triangular. $`\mathrm{}`$ A crucial step for our geometric formulas for intersection indices (to be proven in section 5) is the expression of the intersection number of $`n`$ hyperplanes in terms of their equations. We can do this for any $`n`$ hyperplanes $`H_1,\mathrm{},H_n`$ in $`(M)`$ such that their generic fibres $`H_{1K},\mathrm{},H_{nK}`$ meet properly on $`(V)`$. ###### Theorem 3.4 Let $`M`$ be a lattice in $`V`$. We fix a basis $`x_1,\mathrm{},x_n`$ of $`M`$, and denote by $`x_1^{},\mathrm{},x_n^{}M^{}`$ the dual basis. Let $`H_1,\mathrm{},H_n`$ be hyperplanes in $`(M)`$ which intersect properly on the generic fibre. Let $`f_i=_ja_{ji}x_j^{}M^{}`$ be a linear homogeneous element generating the ideal corresponding to $`H_i`$ and put $`A=(a_{kl})_{k,l}`$. Then we have the following formula for the intersection number of $`H_1,\mathrm{},H_n`$: $$<H_1,\mathrm{},H_n>=v(detA).$$ Proof: By Lemma 3.3 we find some $`C\text{GL}(n,R)`$ and a permutation matrix $`D`$ such that $`CAD=B=(b_{ij})`$ is upper triangular and satisfies the inequalities $$v(b_{11})\mathrm{}v(b_{nn})\text{ and}$$ $$v(b_{ii})v(b_{ij})\text{ for all }ij.$$ There is a permutation $`\sigma `$ of $`\{1,\mathrm{}n\}`$ such that $`AD`$ is the coefficient matrix for the hyperplanes $`H_{\sigma (1)},\mathrm{},H_{\sigma (n)}`$. By Lemma 3.2, the linear element $`_{ji}b_{ji}x_j^{}`$ corresponds to the hyperplane $`{}_{}{}^{t}C_{}^{1}(H_{\sigma (i)})`$ for all $`i=1,\mathrm{},n`$. Now $`<H_1,\mathrm{},H_n>=`$ $`<{}_{}{}^{t}C_{}^{1}(H_{\sigma (1)}),\mathrm{},{}_{}{}^{t}C_{}^{1}(H_{\sigma (n)})>`$ and $`v(detA)=v(detB)`$. Hence we can assume that $`A`$ is upper triangular with $`v(a_{11})\mathrm{}v(a_{nn})`$ and $`v(a_{ii})v(a_{ij})`$ for $`ij`$. We can assume that $`H_i\mathrm{}`$, since otherwise our claim is trivial. For all $`aR`$ we denote by $`\overline{a}`$ its image in $`k`$. The reduction $`(H_i)_k`$ of $`H_i`$ corresponds to the homogeneous ideal generated by $`\overline{f_i}=\overline{a_{ji}}x_j^{}`$ in $`(M_k^{})`$. Let $`\overline{A}`$ be the matrix $`(\overline{a_{ij}})`$. Now $`(H_i)_k`$ is the linear cycle corresponding to the subspace $`L_kM_k`$ which is equal to $`\text{ker}{}_{}{}^{t}\overline{A}`$ via the identification of $`M_k`$ with $`k^n`$ given by the reductions of $`x_1,\mathrm{},x_n`$. We will first assume that $`H_1,\mathrm{},H_n`$ meet properly on the whole of $`(M)`$. Hence their intersection consists of one point in the special fibre, and $`v(a_{11})=\mathrm{}=v(a_{n1n1})=0`$. We can therefore assume that $`a_{11}=\mathrm{}=a_{n1n1}=1`$. Hence for all $`kn1`$ the homogeneous ideal in $`\text{Sym }M^{}`$ generated by $`f_1=x_1^{},f_2=a_{12}x_1^{}+x_2^{},\mathrm{},f_k=_{j=1}^{k1}a_{jk}x_j^{}+x_k^{}`$ is equal to the homogeneous ideal generated by $`x_1^{},\mathrm{},x_k^{}`$. Now we can compute Serre’s intersection index as follows: Note that all $`H_i`$ are isomorphic to $`_R^{n2}`$, hence they are locally Cohen-Macaulay (even regular). Assume that $`n>2`$. The closed subset $`H_1H_2`$ of $`(M)`$ is given by the homogeneous ideal $`(f_1,f_2)=(x_1^{},x_2^{})`$; the corresponding reduced closed subscheme is the linear cycle $`(N)(M)`$ for $`N=(M^{}/Rx_1^{}+Rx_2^{})^{}`$. Hence $`W=H_1H_2`$ is a prime cycle. We have $`i_W(H_1,H_2)=1`$, hence $`H_1H_2=W`$. Besides, $`W`$ is a projective space over $`R`$, hence also locally Cohen-Macaulay. The same argument (if $`n>3`$) implies that $`H_1H_2H_3`$ is equal to the cycle given by the irreducible subset $`H_1H_2H_3`$. Finally we find that $`H_1H_2\mathrm{}H_{n1}=W`$ where $`W`$ is the prime cycle corresponding to the homogeneous ideal $`(x_1^{},\mathrm{},x_{n1}^{})`$. Now it is easy to calculate $`<H_1,\mathrm{},H_n>=v(a_{nn})=v(detA)`$, which proves our claim in the case of proper intersection. Hence we can now assume that there is an $`s<n1`$ such that $`v(a_{11})=\mathrm{}=v(a_{ss})=0`$ and $`v(a_{kk})>0`$ if $`k>s`$. We write $`l_i=v(a_{ii})`$, and we can again assume that $`a_{11}=\mathrm{}=a_{ss}=1`$. For all $`m=1,\mathrm{},n1`$ let $`Y_m`$ be the cycle corresponding to the integral subscheme given by the homogeneous ideal $`(x_1^{},\mathrm{},x_m^{})`$ of $`\text{Sym }M^{}`$, and let $`Z_m`$ be the cycle corresponding to the integral subscheme given by $`(x_1^{},\mathrm{},x_{m1}^{},\pi )`$, where $`\pi `$ is a fixed prime element in $`R`$. Assume that $`m`$ is a number with $`s<m<n`$. Then the intersection of $`Y_{m1}`$ and $`H_m`$ has two irreducible components, namely $`Y_m`$ and $`Z_m`$. Since both have codimension $`m`$, the cycles $`Y_{m1}`$ and $`H_m`$ meet properly. We want to calculate $`i_{Y_m}(Y_{m1},H_m)`$ and $`i_{Z_m}(Y_{m1},H_m)`$. Let $`y`$ respectively $`z`$ be the generic points of $`Y_m`$ respectively $`Z_m`$. Since $`m<n`$, they are both contained in $`U=\{x_n^{}0\}`$. We write $`A=𝒪_{(M),y}`$. Then $`A=R[y_1,\mathrm{},y_{n1}]_{(y_1,\mathrm{},y_m)}`$ with $`y_i=x_i^{}/x_n^{}`$. Since $`Y_{m1}`$ and $`H_m`$ are locally Cohen-Macaulay, we have $$i_{Y_m}(Y_{m1},H_m)=l_A(A/(y_1,\mathrm{},y_{m1},\pi ^{l_m}y_m)).$$ As $`\pi `$ is a unit in $`A`$, this is equal to $`l_A(A/(y_1,\mathrm{},y_m))=1`$. Similarly, we put $`B=𝒪_{(M),z}`$, hence $`B=R[y_1,\mathrm{},y_{n1}]_{(y_1,\mathrm{},y_{m1},\pi )}`$, and we get $$i_{Z_m}(Y_{m1},H_m)=l_B(B/(y_1,\mathrm{},y_{m1},\pi ^{l_m}y_m)).$$ Here $`y_m`$ is a unit in $`B`$, hence this length is equal to $$l_B(B/(y_1,\mathrm{},y_{m1},\pi ^{l_m}))=l_{(R[y_m,\mathrm{},y_{n1}]/(\pi ^{l_m}))_{(\pi )}}((R[y_m,\mathrm{},y_{n1}]/(\pi ^{l_m}))_{(\pi )})=l_m,$$ since the only ideals in $`(R[y_m,\mathrm{},y_{n1}]/(\pi ^{l_m}))_{(\pi )}`$ are $`0`$ and $`(\pi ^k)`$ for $`0<k<l_m`$ and these are all distinct. Now we will prove by induction that $`H_1\mathrm{}H_m`$ is equal to the class of $`Y_m+(_{im}l_i)Z_m`$ in $`CH_{_{im}H_i}^m((M))`$ for all $`mn1`$. Since $`H_1`$ is irreducible, we have $`l_1=0`$ and $`H_1=Y_1`$, which is our claim for $`m=1`$. Now we come to the induction step. Assume that our claim holds for some $`m`$ with $`1mn2`$. If $`ms`$, we can move to $`m+1`$ using the above calculations of intersection indices. Let us now assume that $`m>s`$, hence that $`l_m`$ and $`l_{m+1}`$ are strictly positive. Then $`H_{m+1}`$ meets $`Y_m`$ properly in the components $`Y_{m+1}`$ and $`Z_{m+1}`$ and we can calculate $`Y_mH_{m+1}`$ via Serre’s intersection: $`Y_mH_{m+1}`$ is induced by the cycle $`Y_{m+1}+l_{m+1}Z_{m+1}`$ by our previous calculations. Note that $`f_{m+1}`$ is contained in $`(x_1^{},\mathrm{},x_{m1}^{},\pi )`$, since both $`l_m`$ and $`l_{m+1}`$ are strictly positive. Hence $`Z_m`$ is contained in $`H_{m+1}`$. Now we can determine the intersection $$Z_mH_{m+1}CH_{Z_mH_{m+1}}^{m+1}((M))=CH_{Z_m}^{m+1}((M))$$ by the recipe we described in section 2 for intersections with divisors. Since $`H_{m+1}`$ is linearly equivalent to the hyperplane $`H^{}`$ given by the ideal $`(x_m^{})`$, we find that $`Z_mH_{m+1}`$ is equal to the image of $`Z_mH^{}`$ in $`CH_{Z_m}^{m+1}((M))`$. Now $`Z_mH^{}=Z_{m+1}`$ and $`Z_m`$ and $`H^{}`$ meet properly in this irreducible set with $`i_{Z_{m+1}}(Z_m,H^{})=1`$, which implies that $`Z_mH^{}`$ is induced by the cycle $`Z_{m+1}`$. Altogether we find that $`H_1\mathrm{}H_{m+1}`$ is the image of $`Y_{m+1}+(l_1+\mathrm{}+l_{m+1})Z_{m+1}`$ in $`CH_{_{im+1}H_i}^{m+1}((M))`$, which finishes the proof of our claim. We know now that $`H_1\mathrm{}H_{n1}`$ is the image of $`Y_{n1}+(l_1+\mathrm{}+l_{n1})Z_{n1}`$ in $`CH_{_{in1}H_i}^{n1}((M))`$. Since $`s<n1`$, we have $`l_{n1}>0`$ and $`l_n>0`$. Now $`Y_{n1}`$ meets $`H_n`$ properly, hence we can apply our result in the case of hyperplanes meeting properly on the whole of $`(M)`$ and find $`\mathrm{deg}(Y_{n1}H_n)=l_n`$. Besides, we have $`Z_{n1}H_n=Z_{n1}`$, and as in our induction step we can show that $`Z_{n1}H_n`$ is the image of $`Z_n`$ in $`CH_{Z_{n1}}^n((M))`$. Since $`Z_n`$ is a $`k`$-rational point in the special fibre, we get $`\mathrm{deg}(Z_{n+1}H_n)=1`$. Altogether we find that $$<H_1,\mathrm{},H_n>=l_1+\mathrm{}+l_n=v(detA),$$ whence our claim. $`\mathrm{}`$ ## 4 The Bruhat-Tits building for $`PGL`$ We denote by $`X`$ the Bruhat-Tits building corresponding to the group $`G=PGL(V)`$ (see \[Br-Ti\]). $`X`$ is a metric space with a continuous $`G`$-action and a simplicial structure. For our purposes, we can think of it as the geometric realization of the following simplicial complex: We call two lattices in $`V`$ equivalent, if they differ by a factor in $`K^\times `$, and we write $`\{M\}`$ for the equivalence class of the lattice $`M`$. Two different lattice classes $`\{M^{}\}`$ and $`\{N^{}\}`$ are called adjacent, if there are representatives $`M`$ and $`N`$ of $`\{M^{}\}`$ and $`\{N^{}\}`$ such that $$\pi NMN.$$ This relation defines a flag complex, namely the simplicial complex whose vertex set is the set of all classes $`\{M\}`$, and whose simplices are the sets of pairwise adjacent lattice classes. Note that it carries a natural $`G`$-action. If $`n=2`$, then $`X`$ is an infinite regular tree, with $`q+1`$ edges meeting in every vertex. The building $`X`$ is the union of its apartments, which correspond to the maximal split tori in $`G`$. We can describe them as follows: For every decomposition $`V=_{1in}L_i`$ of $`V`$ in one-dimensional subspaces $`L_i`$ generated by some vector $`v_i`$ we define an apartment as the subcomplex of $`X`$ given by all lattices $`M`$ which can be diagonalized with respect to our decomposition, i.e. $`M=_{i=1}^nR\pi ^{k_i}v_i`$ for some integers $`k_i`$. ###### Definition 4.1 The combinatorial distance $`\text{dist}(x,y)`$ of two vertices $`x`$ and $`y`$ in $`X`$ is defined as $`\text{dist}(x,y)`$ $`=`$ $`\mathrm{min}\{k:\text{ there are vertices }x=x_0,x_1,\mathrm{},x_k=y,`$ $`\text{so that }x_i\text{ and }x_{i+1}\text{ are adjacent for all }i=0,\mathrm{},k1.\}.`$ Hence dist is the minimal number of $`1`$-simplices forming a path between $`x`$ and $`y`$. Note that dist is in general not proportional to the metric on $`X`$. ###### Lemma 4.2 Let $`x=\{M\}`$ and $`y=\{L\}`$ be two vertices in $`X`$, and define $`s`$ $`=`$ $`\mathrm{min}\{k:\pi ^kLM\}\text{and}`$ $`r`$ $`=`$ $`\mathrm{max}\{k:M\pi ^kL\}.`$ Then we have $`\text{dist}(x,y)=sr`$. Proof: Note that the term on the right hand side is independent of the choice of a representative of the lattice classes. Put $`d=\text{dist}(x,y)`$. Then we find lattices $`M=M_0,M_1,\mathrm{},M_d`$ such that $`M_d=\alpha L`$ for some $`\alpha K^\times `$ and such that $$\pi ^dM_d\pi ^{d1}M_{d1}\mathrm{}\pi M_1M_0M_1\mathrm{}M_d.$$ Hence $`\pi ^d\alpha LM`$, which implies $`sd+v(\alpha )`$, and $`M\alpha L`$, which implies $`rv(\alpha )`$. Altogether we find that $`srd`$. Let us now show that also $`srd`$ is true. We have by definition $`\pi ^sLM`$ and $`M\pi ^rL`$. Put $`L^{}=\pi ^rL`$. By the invariant factor theorem, we find an $`R`$-basis $`w_1,\mathrm{},w_n`$ of $`L^{}`$, such that $`M=\pi ^{k_i}Rw_i`$ for some integers $`k_i`$. Since $`\pi ^{sr}L^{}ML^{}`$, all $`k_i`$ are between $`0`$ and $`sr`$. Now put $`l_{ij}=\mathrm{max}\{0,k_ji\}`$ for all $`i\{0,\mathrm{},sr\}`$ and all $`j\{1,\mathrm{},n\}`$. Then we define for $`i=0,\mathrm{},sr`$ $$M_i=\pi ^{l_{i1}}Rw_1+\mathrm{}+\pi ^{l_{in}}Rw_n.$$ Note that $`M_0=M`$, and $`M_{sr}=L^{}`$. We have for all $`i=0,\mathrm{},sr1`$ the inclusions $`\pi M_{i+1}M_iM_{i+1}`$. Since either $`\{M_i\}=\{M_{i+1}\}`$, or $`\{M_i\}`$ and $`\{M_{i+1}\}`$ are adjacent, we found a chain of adjacent lattices of length $`sr`$ connecting $`\{M\}`$ and $`\{L\}`$, which implies our claim. $`\mathrm{}`$ ## 5 Intersection indices via combinatorial geometry Let us fix a lattice $`M`$ in $`V`$. We will now interpret non-Archimedean intersection numbers of linear cycles on $`(M)`$ with the combinatorial geometry of $`X`$. Take some $`d2`$ and let $`(N_1),\mathrm{},(N_d)`$ be $`d`$ linear cycles on $`(M)`$ with $$\underset{i=1}{\overset{d}{}}\text{codim}(N_i)=n,$$ which meet properly on $`(M)`$. Then $`N_1,\mathrm{},N_d`$ are split submodules of $`M`$ of rank $`r_1,\mathrm{},r_d`$ satisfying $`_{i=1}^dr_i=(d1)n`$ by 3.1. We will always assume that $`r_in`$. For all $`j=1,\mathrm{},d`$ we put $`L_j=_{ij}N_iM`$. Since all $`M/N_i`$ are torsion free, the same holds for $`M/L_j`$, so that $`L_j`$ is a split submodule of $`M`$. Now let $`F`$ be set of all vertices in $`X`$ of the form $$\{\pi ^{k_1}L_1\mathrm{}\pi ^{k_d}L_d\}\text{for some}k_1,\mathrm{},k_d.$$ (Alternatively, we can also work with the convex hull of $`F`$ in $`X`$.) Note that the intersection $`_{i=1}^dN_i`$ is zero, as $`(N_1),\mathrm{},(N_d)`$ do not meet on the generic fibre. Therefore the sum $`\pi ^{k_1}L_1\mathrm{}\pi ^{k_d}L_d`$ is direct. On the other hand, an easy calculation shows that $`dimL_{jK}nr_j`$. Since $`_jdimL_{jK}_j(nr_j)=n`$, we must have $`dimL_{jK}=nr_j`$ and $`V=_jL_{jK}`$, so that $`\pi ^{k_1}L_1\mathrm{}\pi ^{k_d}L_d`$ is indeed a lattice of full rank in $`V`$. Obviously, $`F`$ is the set of vertices contained in either a full apartment or an intersection of affine hyperplanes in some apartment. (Hence its convex hull is either an apartment of an intersection of affine hyperplanes.) Let us describe $`F`$ in two special cases: i) The case of $`n`$ hyperplanes: First note that $`n`$ hyperplanes $`H_{1K},\mathrm{},H_{nK}`$ in the generic fibre $`(V)`$ intersecting properly, i.e. not at all, in $`(V)`$, define an apartment $`A(H_{1K},\mathrm{},H_{nK})`$ in $`X`$ as follows: For all $`i`$ the hyperplane $`H_{iK}`$ is the linear cycle corresponding to an $`(n1)`$-dimensional subspace $`W_i`$ of $`V`$. Since the $`H_{iK}`$ intersect properly, $`W_i`$ is equal to $`0`$. For all $`j=1,\mathrm{},n`$ put $`U_j=_{ij}W_i`$. Then all $`U_j`$ are one-dimensional and $`V=U_j`$. We denote the apartment corresponding to this decomposition (see section 4) by $`A(H_{1K},\mathrm{},H_{nK})`$. If $`H_1,\mathrm{},H_n`$ are hyperplanes in $`(M)`$ which intersect properly on the whole of $`(M)`$, then the subset $`F`$ is just the apartment $`A(H_{1K},\mathrm{}H_{nK})`$ corresponding to the generic fibres. ii) The case of two cycles: Let $`N_1`$ and $`N_2`$ be two non-trivial split submodules of $`M`$ of rank $`p`$ respectively $`q=np`$, such that the corresponding linear cycles $`Z_1=(N_1)`$ and $`Z_2=(N_2)`$ meet properly on $`(M)`$. Then the corresponding subset $`F`$ is $$\{\{N_1+\pi ^kN_2\}:k\}.$$ In fact, one can show that $`F`$ defines a doubly infinite geodesic in the building $`X`$ whose boundary points (with respect to the Borel-Serre compactification of $`X`$) are the parabolics induced by the vector spaces $`W_i=N_i_RK`$ for $`i=1`$ and $`2`$. We will now show that the intersection number of $`d`$ linear cycles meeting properly on $`(M)`$ is the combinatorial distance of the set $`F`$ to the lattice $`M`$. If $`n`$ is equal to $`2`$ (i.e. $`X`$ is the Bruhat-Tits tree associated to $`PGL(2,K)`$), and we consider the case of two hyperplanes in $`(M)`$, then the ensuing formula is due to Manin (see \[Ma\], p. 232). Note that in this setting our special cases i) and ii) from above coincide, and the subset $`F`$ we are dealing with is a geodesic in the Bruhat-Tits tree. ###### Theorem 5.1 The intersection number of $`(N_1),\mathrm{},(N_d)`$ on $`(M)`$ can be expressed as follows with the combinatorial geometry of $`X`$: $$<(N_1),\mathrm{},(N_d)>=\text{dist}(\{M\},F),$$ where dist denotes the combinatorial distance in $`X`$, and $`\{M\}`$ is the vertex in $`X`$ defined by our fixed lattice $`M`$. Proof: Recall that $`r_i`$ is the rank of $`N_i`$ and that the numbers $`m_i=nr_i`$ satisfy $`_{i=1}^dm_i=n`$. Besides, put $`n_0=0`$ and $`n_i=m_1+\mathrm{}+m_i`$ for $`i=1,\mathrm{},d`$. Then $`n=n_d`$. For all $`i=1,\mathrm{},d`$ let $`g_{n_{i1}+1},\mathrm{},g_{n_i}`$ be a $`R`$-basis of $`(M/N_i)^{}M^{}`$. The elements $`g_{n_{i1}+1},\mathrm{},g_{n_i}`$ generate the homogeneous prime ideal corresponding to $`(N_i)(M)`$. Besides, fix a basis $`y_1,\mathrm{},y_n`$ of $`M`$, and let $`A^{}=(a_{ij}^{})`$ be the coordinate matrix of $`g_1,\mathrm{},g_n`$ with respect to the dual basis $`y_1^{},\mathrm{},y_n^{}`$, i.e. $$g_j=\underset{i=1}{\overset{n}{}}a_{ij}^{}y_i^{}.$$ Now choose an element $`a_{ij}^{}`$ with $`jn_1`$ such that $`v(a_{ij}^{})`$ is minimal among the entries of the first $`n_1`$ columns. We remove the $`j`$-th column and insert it before the first one. Then we switch rows so that this element sits in the upper left corner, and we perform some elementary row operations to eliminate $`a_{21}^{},\mathrm{},a_{n1}^{}`$. Among the $`j`$ between $`2`$ and $`n_1`$ and among the $`i2`$ we choose an entry of minimal valuation. Again, after removing the corresponding column and putting it between the first and the second one, and after elementary row operations involving only the last $`n1`$ rows we can assume that $`a_{i2}^{}=0`$ for all $`i>2`$. We continue in this way until we reach the $`n_1`$-th column. Then $`a_{ij}^{}`$ is zero for $`jn_1`$ and $`i>j`$, and the upper left corner satisfies the following divisibility conditions: $`v(a_{ii}^{})v(a_{i+1i+1}^{})\text{for }in_11\text{and}`$ $`v(a_{ii}^{})v(a_{ij}^{})\text{for }i<jn_1.`$ Now choose an element $`a_{ij}^{}`$ with $`i,jn_1+1`$ and $`jn_2`$ such that $`v(a_{ij}^{})`$ is minimal among the entries $`a_{ij}^{}`$ with $`in_1+1`$ and $`n_1+1jn_2`$. As before, we permute the columns corresponding to $`g_{n_1+1},\mathrm{},g_{n_2}`$ and perform elementary row operations involving only the rows with index bigger or equal to $`n_1+1`$ to achieve $`a_{ij}^{}=0`$ for $`jn_2`$ and $`i>j`$. We continue this process until we worked our way through the whole matrix. We see that after permuting the equations $`g_{n_{i1}+1},\mathrm{},g_{n_i}`$ corresponding to each of the $`N_1,\mathrm{},N_d`$, and after switching to another basis of $`M`$ (in order to take care of the row operations) we can assume that our coordinate matrix $`A^{}`$ is upper triangular. Let us denote the reduction of elements in $`R`$ or $`M`$ or of $`R`$-matrices by overlining. Besides we use the following notation for submatrices: For any $`n\times n`$-matrix $`D=(d_{kl})`$ and any $`ijn`$ we write $`D(ij)`$ for the $`n\times (ji+1)`$-submatrix consisting of the columns $`i,i+1,\mathrm{},j`$. Similarly we write $`D(ij)`$ for the $`(ji+1)\times n`$-submatrix consisting of the rows $`i,i+1,\mathrm{},j`$. By $`D(i_1j_1,i_2j_2)`$ we mean the submatrix where we take columns $`i_1,\mathrm{},j_1`$ followed by columns $`j_1,\mathrm{},j_2`$. Besides, $`D(i_1j_1)(i_2j_2)`$ is the submatrix consisting of all entries $`d_{kl}`$ with $`i_1kj_1`$ and $`i_2lj_2`$. For $`n\times n`$-matrices we will furthermore abbreviate $`D_{ij}=D(n_{i1}+1n_i)(n_{j1}+1n_j)`$, if $`i`$ and $`j`$ are $`d`$. If we divide $`D`$ into rectangular submatrices according to our partition $`n=m_1+\mathrm{}+m_d`$, then $`D_{ij}`$ is the rectangle in position $`(ij)`$. Since $`g_{n_{i1}+1},\mathrm{},g_{n_i}`$ form a basis of the split $`R`$-submodule $`(M/N_i)^{}`$ of $`M^{}`$, their reductions generate a vector space of rank $`m_i`$ over $`k`$, so that the coordinate matrix $`\overline{A}^{}(n_{i1}+1n_i)`$ has full rank $`m_i`$. Since $`(N_1),\mathrm{},(N_d)`$ meet properly, their intersection is empty or zero-dimensional and contained in the special fibre. Hence $`_i(N_{ik})`$ is empty or has dimension zero, so that $`_iN_{ik}`$ has dimension $`1`$. Since $`N_{ik}`$ is equal to the kernel of $`(\overline{g}_{n_{i1}+1},\mathrm{},\overline{g}_{n_i})`$, this means that $`rk\overline{A}^{}n1`$. The intersection of $`(N_1),\mathrm{},(N_d)`$ is empty iff $`\text{rk}\overline{A}^{}=n`$. In this case all elements on the diagonal of $`A^{}`$ are units. After switching to another basis of $`M`$, we can therefore assume that $`A^{}`$ is diagonal. Then $`M=L_1+\mathrm{}+L_d`$, so that $`\{M\}`$ is actually contained in $`F`$, and our formula holds. Hence we only have to deal with the case $`\text{rk}\overline{A}^{}=n1`$. Here the intersection of $`(N_1),\mathrm{},(N_d)`$ is not empty. As $`v(detA^{})>0`$, we find indices $`pqd`$ with $`v(detA_{11}^{})=\mathrm{}=v(detA_{p1p1}^{})=0,`$ $`v(detA_{pp}^{})>0,v(detA_{qq}^{})>0\text{and}`$ $`v(detA_{q+1q+1}^{})=\mathrm{}=v(detA_{dd}^{})=0.`$ Note that $`p2`$, since $`\overline{A}^{}(1n_1)`$, and thus $`\overline{A}_{11}^{}`$ has full rank. All elements on the diagonal of $`A_{11}^{},\mathrm{},A_{p1p1}^{}`$, $`A_{q+1q+1}^{},\mathrm{},A_{dd}^{}`$ are units. After performing some elementary row operations we can assume that $`A_{ij}^{}=0`$, if $`i<jp1`$ or if $`jq+1`$ and $`i<j`$. Note that the divisibility conditions in $`A_{qq}^{}`$ imply that all elements on the diagonal except possibly the last one are units. Therefore we can eliminate all entries $`a_{ij}^{}`$ in $`A^{}`$ such that $`1in_{q1}`$ and $`n_{q1}+1jn_q1`$ by elementary row operations. Hence after switching to another basis $`y_1,\mathrm{},y_n`$ of $`M`$ we may assume that our coordinate matrix $`A^{}`$ contains zeroes above $`A_{22}^{},\mathrm{},A_{p1p1}^{}`$ and $`A_{q+1,q+1}^{},\mathrm{}A_{dd}^{}`$, and that all columns above $`A_{qq}^{}`$ are zero except possibly the last one. Now recall that for $`j=1,\mathrm{},d`$ the module $`L_j=_{ij}N_i`$ is a split submodule of $`M`$ of rank $`m_j`$. For each $`j=1,\mathrm{},d`$ choose an $`R`$-basis $$w_{n_{j1}+1},\mathrm{},w_{n_j}$$ of $`L_j`$. We denote by $`B^{}`$ the transpose of the coordinate matrix of $`w_1,\mathrm{},w_n`$ with respect to $`y_1,\mathrm{},y_n`$, i.e. $`B^{}`$ is the matrix $`B^{}=(b_{ij}^{})_{ij}`$ so that $$w_i=\underset{j=1}{\overset{n}{}}b_{ij}^{}y_j.$$ For all $`j=1,\mathrm{},d`$ there are matrices $`C_j,D_jGL(m_j,R)`$ such that $`C_jB_{jj}^{}D_j`$ is a diagonal matrix with entries $`\beta _{ii}`$ such that $`v(\beta _{11})\mathrm{}v(\beta _{m_jm_j})`$. Let $`C`$ (respectively $`D`$) be the matrix with diagonal components $`C_1,\mathrm{},C_d`$ (respectively $`D_1,\mathrm{},D_d`$), and define $`B=CB^{}D`$. Then $`D`$ is the transpose of the transition matrix from a base $`x_1,\mathrm{},x_n`$ of $`M`$ to our base $`y_1,\mathrm{},y_n`$. If we put $$v_i=\underset{j=1}{\overset{n}{}}b_{ij}x_j,$$ then for all $`hd`$ the elements $`v_{n_{h1}+1},\mathrm{},v_{n_h}`$ form a basis of $`L_h`$. The matrix $`B`$ is by definition the transpose of the transition matrix from $`x_1,\mathrm{},x_n`$ to $`v_1,\mathrm{},v_n`$. Now $`D^1`$ is the transition matrix from the dual basis $`x_1^{},\mathrm{},x_n^{}`$ of $`M^{}`$ to $`y_1^{},\mathrm{},y_n^{}`$. Hence the coordinate matrix of $`g_1,\mathrm{},g_n`$ with respect to $`x_1^{},\mathrm{},x_n^{}`$ is equal to $`A^{\prime \prime }=D^1A^{}`$. The matrix $`A^{\prime \prime }`$ has the property that if a block $`A_{ij}^{}`$ is zero or zero up to the last column, then the same holds for $`A_{ij}^{\prime \prime }`$. Besides, $`v(detA_{jj}^{\prime \prime })=v(detA_{jj}^{})`$. However, the diagonal blocks $`A_{11}^{\prime \prime },\mathrm{},A_{dd}^{\prime \prime }`$ may not be upper triangular any longer. After permuting the first $`n_1=m_1`$ columns we may assume that $`v(a_{n_1n_1}^{\prime \prime })v(a_{n_1j}^{\prime \prime })`$ for all $`j=1,\mathrm{},n_1`$. By a series of elementary column operations we can eliminate $`a_{n_11}^{\prime \prime },\mathrm{},a_{n_1n_11}^{\prime \prime }`$. Now we permute the first $`n_11`$ columns to achieve $`v(a_{n_11j}^{\prime \prime })v(a_{n_11n_11}^{\prime \prime })`$ for all $`j=1,\mathrm{},n_11`$, and we clear out $`a_{n_111}^{\prime \prime },\mathrm{},a_{n_11n_12}^{\prime \prime }`$. Note that these column operations affect only the first $`n_11`$ columns, hence the first $`n_11`$ elements in the $`n_1`$-th row remain zero. We repeat this process until $`A_{11}^{\prime \prime }`$ is upper triangular. These column operations amount to passing to another basis $`f_1,\mathrm{},f_{n_1}`$ of $`(M/N_1)^{}`$ Now we work on $`A^{\prime \prime }(n_1+1n_2)`$. First we switch columns inside this block to achieve $`v(a_{n_2j}^{\prime \prime })v(a_{n_2n_2}^{\prime \prime })`$ for all $`j=n_1+1,\mathrm{},n_21`$, and we clear out $`a_{n_2n_1+1}^{\prime \prime },\mathrm{},a_{n_2n_21}^{\prime \prime }`$, then we eliminate $`a_{n_21n_1+1}^{\prime \prime },\mathrm{},a_{n_21n_22}^{\prime \prime }`$, and so on, until $`A_{22}^{\prime \prime }`$ is upper triangular. We do this with block after block until $`A^{\prime \prime }`$ is upper triangular. Since the columns were transformed block by block, we can find for each $`i=1,\mathrm{},d`$ a new $`R`$-basis $$f_{n_{i1}+1},\mathrm{},f_{n_i}$$ of $`(M/N_i)^{}`$, whose coordinate matrix $`A`$ with respect to the basis $`x_1^{},\mathrm{},x_n^{}`$ is upper triangular and has the property that the columns above $`A_{22},\mathrm{},A_{p1p1}`$ and above $`A_{q+1q+1}\mathrm{},A_{dd}`$ are zero. Besides, we still have $`v(detA_{pp})>0`$ and $`v(detA_{qq})>0`$. Since there exists a matrix $`DGL(m_q,R)`$ such that $$A(n_{q1}+1n_q)(1n_{q1})=A^{\prime \prime }(n_{q1}+1n_q)(1n_{q1})D,$$ we still know that $`\text{rk}\overline{A}(n_{q1}+1n_q)(1n_{q1})1`$. Let $`H_j`$ be the hyperplane given by the linear homogeneous element $`f_j`$. Then $`H_{n_{i1}+1},\mathrm{},H_{n_i}`$ intersect properly and $`(N_i)=H_{n_{i1}+1}\mathrm{}H_{n_i}`$. Hence by 3.4 we have $$<(N_1),\mathrm{},(N_d)>=v(detA).$$ Besides note that $$f_j(v_i)=\underset{k}{}a_{kj}b_{ik},$$ so that the entries of $`BA`$ are equal to $`(f_j(v_i))_{ij}`$. Now by definition, $`N_i`$ lies in the kernel of all $`f_{n_{i1}+1},\mathrm{},f_{n_i}`$. Since $`L_j`$ is contained in $`N_i`$ for all $`ij`$, this implies that the blocks $`(BA)_{ij}`$ are equal to $`0`$ for $`ij`$. Since all $`B_{jj}`$ are diagonal matrices, $`B`$ is therefore upper triangular. Besides we have $`B_{ij}=0`$, if $`i<jp1`$, and if $`i<j`$ and $`j>q`$. Since $`v_{n_{j1}+1},\mathrm{},v_{n_j}`$ are a basis of the split submodule $`L_jM`$, their reductions are still linear independent, hence the rank of $`\overline{B}(n_{j1}+1n_j)`$ is equal to $`m_j`$. We will now show that only the last diagonal element $`b_{n_jn_j}`$ in the block $`B_{jj}`$ may not be a unit. Indeed, assume that $`v(b_{ii})>0`$ for some $`n_{j1}+1in_j1`$. By the divisibility conditions along the diagonal of $`B_{jj}`$ we have $`v(b_{n_j1n_j1})>0`$ and $`v(b_{n_jn_j})>0`$. For all $`l=n_{j1}+1,\mathrm{},n_j`$ we find $`f_l(v_{n_j1})=b_{n_j1n_j1}a_{n_j1l}`$, so that $`\overline{f_l(v_{n_j1})}=0`$. Similarly, $`\overline{f_l(v_{n_j})}=0`$, so that $`\overline{v_{n_j1}}`$ and $`\overline{v_{n_j}}`$ are contained in $`N_{jk}`$. Since $`L_j`$ is contained in all $`N_i`$ for $`ij`$, the elements $`\overline{v_{n_j1}}`$ and $`\overline{v_{n_j}}`$ lie also in $`N_{ik}`$ for all $`ij`$. Hence we found two linear independent vectors in $`_iN_{ik}`$, which contradicts our assumption that $`(N_1),\mathrm{},(N_d)`$ meet properly. We will now show the following claim (1) There exists an index $`tp1`$ such that $`\overline{A}(1n_{t1},n_t+1n)`$ has full rank. Let us first assume that $`p<q`$. Note that $`\text{rk}\overline{A}(n_{q1}+1n_q)`$ is strictly smaller than $`m_q`$, since $`v(det(A_{qq}))>0`$. Since $`\text{rk}\overline{A}=n1`$, the matrix $`\overline{A}(n_{p1}+1n_{q1})`$ must have full rank $`(n_{q1}n_{p1})`$. Disregarding the zeros, we find that also the $`(n_{q1}n_{p1})\times (n_qn_{p1})`$-matrix $$\overline{A}(n_{p1}+1n_{q1})(n_{p1}+1n_q)$$ must have full rank $`(n_{q1}n_{p1})`$. This matrix consists of the two vertical chunks $$\overline{A}(n_{p1}+1n_{q1})(n_{p1}+1n_{q1})\text{ and }\overline{A}(n_{p1}+1n_{q1})(n_{q1}+1n_q).$$ We know that the second chunk $`\overline{A}(n_{p1}+1n_{q1})(n_{q1}+1n_q)`$ has rank $`1`$. Hence the rank of the first chunk $$\mathrm{\Omega }:=\overline{A}(n_{p1}+1n_{q1})(n_{p1}+1n_{q1})$$ must be $`(n_{q1}n_{p1}1)`$. Besides, as $`v(detA_{pp})>0`$, there must be an element on the diagonal of $`A_{pp}`$ which has positive valuation, i.e. there exists an index $`n_{p1}+1in_p`$ such that $`\overline{a}_{ii}=0`$. Since $`\mathrm{\Omega }`$ has rank $`(n_{q1}n_{p1}1)`$, the upper left corner $`\overline{A}(n_{p1}+1i1)(n_{p1}+1i)`$ must have full rank $`(i1n_{p1})`$. Now recall that $`\overline{A}(n_{p1}+1n_p)`$ has full rank, hence the first $`in_{p1}`$ columns of this matrix, namely $`\overline{A}(n_{p1}+1i)`$, also have full rank $`in_{p1}`$. Put $$\lambda _j=\overline{A}(n_{p1}+1i)(jj),$$ which is just the $`j`$-th row of this matrix. We have seen that $`\lambda _{n_{p1}+1},\mathrm{},\lambda _{i1}`$ are linear independent, hence there exists a row $`\lambda _{j_0}`$ with $`j_0n_{p1}`$ so that $`\lambda _{j_0},\lambda _{n_{p1}+1},\mathrm{},`$$`\lambda _{i1}`$ is a full linear independent set of rows in $`\overline{A}(n_{p1}+1i)`$. Now let $`tp1`$ be the index of the block in which $`\lambda _{j_0}`$ lies, i.e. $`n_{t1}+1j_0n_t`$. We want to show that $`\overline{A}(1n_{t1},n_t+1n)`$ has full rank. It obviously suffices to show that $`\overline{A}(n_t+1n)(n_{t1}+1n)`$ has full rank. Removing the first block of rows $`\overline{A}(n_t+1n)(n_{t1}+1n_t)`$ and putting it before the block $`\overline{A}(n_t+1n)(n_{p1}+1n_p)`$, we see that it remains to show that the matrix $`\mathrm{\Lambda }=\left(\begin{array}{ccc}\overline{A}_{tp}& \mathrm{}& \overline{A}_{td}\\ \overline{A}_{pp}& \mathrm{}& \overline{A}_{pd}\\ & \mathrm{}& \mathrm{}\\ 0& & \overline{A}_{dd}\end{array}\right).`$ has full rank $`nn_{p1}`$. Recall that we fixed an element $`a_{ii}`$ on the diagonal of $`A_{pp}`$ with positive valuation. Now $`\mathrm{\Lambda }=\left(\begin{array}{ccccccc}\overline{a}_{n_{t1}+1,n_{p1}+1}& \mathrm{}& \hfill \overline{a}_{n_{t1}+1,i}\|& & & & \\ & \mathrm{}& \hfill \|& & & & \\ \overline{a}_{n_t,n_{p1}+1}& \mathrm{}& \hfill \overline{a}_{n_ti}\|& & & & \\ \overline{a}_{n_{p1}+1,n_{p1}+1}& \mathrm{}& \hfill \overline{a}_{n_{p1}+1,i}\|& & & & \\ & \mathrm{}& \hfill \|& & & & \\ \overline{a}_{i1,n_{p1}+1}& \mathrm{}& \hfill \overline{a}_{i1i}\|& & & & \\ & & \hfill \|& \hfill & & & \\ & & \hfill \|& \overline{a}_{ii+1}\hfill & \mathrm{}& \overline{a}_{in}& \\ & & \hfill \|& \overline{a}_{i+1i+1}\hfill & \mathrm{}& \overline{a}_{i+1n}& \\ & 0& \hfill \|& & \mathrm{}& \mathrm{}& \\ & & \hfill \|& & & \overline{a}_{nn}& \end{array}\right).`$ Hence in the upper left corner of $`\mathrm{\Lambda }`$ we find the rows $`\lambda _j`$ for $`j=n_{t1}+1,\mathrm{},n_t`$ and $`j=n_{p1}+1,\mathrm{},i1`$, and we know that there are $`(in_{p1})`$ linear independent ones among them. So the upper left corner of $`\mathrm{\Lambda }`$ has full rank $`in_{p1}`$. The lower right corner (which is a $`(ni+1)\times (ni)`$-matrix) has the same rank as the $`(ni+1)\times n`$-matrix $`\overline{A}(in)`$. Since $`\text{rk}(\overline{A})=n1`$, the rank of the lower right corner must therefore be bigger or equal to $`ni`$ which means that this corner has also full rank. Hence $`\mathrm{\Lambda }`$ has full rank. This finishes the proof of claim (1) in the case $`p<q`$. If $`p`$ is equal to $`q`$, similar arguments can be used to prove $`(1)`$. Now we claim that we can calculate our intersection number as follows: For the index $`t`$ we found in (1) we have (2) $`<(N_1),\mathrm{},(N_d)>=v(b_{n_tn_t})`$. Put $`w_1=v_{n_{t1}+1},\mathrm{},w_{m_t}=v_{n_t}`$, which is a basis of $`L_t=_{it}N_i`$. Since $`L_t`$ is a split submodule of $`M`$, we can complete this basis to a basis $`w_1,\mathrm{},w_n`$ of $`M`$. Since $`f_j=_{i=1}^nf_j(w_i)w_i^{}`$, the coordinate matrix of $`f_1,\mathrm{},f_n`$ with respect to the dual basis $`w_1^{},\mathrm{},w_n^{}`$ is equal to $`(f_j(w_i))_{i,j}`$. Using 3.4, we can therefore calculate the intersection number as follows: $$<(N_1),\mathrm{},(N_d)>=v(det(f_j(w_i)_{i,j=1,\mathrm{},n})).$$ Now $`f_j(w_i)=0`$ for $`i=1,\mathrm{},m_t`$ and $`jJ_t:=\{n_{t1}+1,\mathrm{},n_t\}`$. Hence after permuting columns the matrix $`(f_j(w_i))_{i,j=1,\mathrm{},n}`$ looks like this $`\left(\begin{array}{cc}& 0\\ & \end{array}\right)`$, and we get: $$v(det(f_j(w_i)_{i,j=1,\mathrm{},n}))=v(det(f_j(w_i)_{i=1,\mathrm{},m_t,jJ_t}))+v(det(f_j(w_i)_{i=m_t+1,\mathrm{},n,jJ_t})).$$ Note that $`(f_j(w_i))_{i=1,\mathrm{},m_t,jJ_t}=B_{tt}A_{tt}`$. As $`t`$ is strictly smaller than $`p`$, the determinant of $`A_{tt}`$ is a unit. Besides, as we have shown above, only the last element on the diagonal of $`B_{tt}`$ may not be a unit, so that we can calculate the first term as follows: $$v(det(f_j(w_i)_{i=1,\mathrm{},m_t,jJ_t}))=v(b_{n_tn_t}).$$ It remains to show that the second term is zero, hence that $`(\overline{f_j(w_i)})_{i=m_t+1,\mathrm{},n,jJ_t}`$ is an invertible $`(nm_t)\times (nm_t)`$-matrix over $`k`$. Since $`f_j(w_i)=0`$ for $`i=1,\mathrm{},m_t`$ and $`jJ_t`$ we can as well show that the matrix $`(\overline{f_j(w_i)})_{i=1,\mathrm{},n,j=1,jJ_t}`$ has full rank $`nm_t`$. In order to prove this we may change the base of $`M`$ and show that $`\overline{A}(1n_{t1},n_t+1n)=(\overline{f_j(x_i)})_{i=1,\mathrm{},n,j=1,jJ_t}`$ has full rank $`nm_t`$, which was done in (1). Let us fix some $`h=1,\mathrm{},d`$. We will first show by induction that for all $`in_h`$ we have (3) $`b_{n_hn_h}\text{ divides }b_{n_hi}_{j=n_h+1}^ia_{jj},`$ where the empty product is equal to $`1`$. For $`i=n_h`$ this is trivial. So let us suppose our claim is true for all $`i`$ with $`n_hi<i_0`$ for some $`i_0n_h+1`$. Then $$0=f_{i_0}(v_{n_h})=\underset{i=n_h}{\overset{i_0}{}}b_{n_hi}a_{ii_0}.$$ We multiply by $`_{j=n_h+1}^{i_01}a_{jj}`$ and get $$0=\underset{i=n_h}{\overset{i_01}{}}b_{n_hi}(\underset{j=n_h+1}{\overset{i_01}{}}a_{jj})a_{ii_0}+b_{n_hi_0}(\underset{j=n_h+1}{\overset{i_01}{}}a_{jj})a_{i_0i_0},$$ hence the induction hypothesis implies our claim for $`i_0`$. Note that (3) implies that $`v(b_{n_tn_t})=\mathrm{max}_{i=1,\mathrm{},n}v(b_{ii})`$. Indeed, if $`i`$ is not equal to $`n_h`$ for some $`h=1,\mathrm{},d`$, then $`b_{ii}`$ is a unit. Since $`\overline{B}`$ contains no zero lines, for every $`n_h`$ there must be an index $`in_h`$ such that $`v(b_{n_hi})=0`$. By (3) we find that $$v(b_{n_hn_h})\underset{j=n_h+1}{\overset{i}{}}v(a_{jj})v(detA).$$ But $`v(detA)`$ is equal to the intersection number $`<(N_1),\mathrm{},(N_d)>`$, hence equal to $`v(b_{n_tn_t})`$ by $`(2)`$. Now we can prove (4) $`\text{dist}(\{M\},F)=v(b_{n_tn_t}).`$ We take a lattice class $`\{L\}`$ corresponding to a vertex in $`F`$, i.e. $`L=\pi ^{k_1}L_1+\mathrm{}+\pi ^{k_{d1}}L_{d1}+L_d`$ for some integers $`k_1,\mathrm{},k_{d1}`$. Then $`\mathrm{min}\{k:\pi ^kLM\}k_t`$ and $`\mathrm{max}\{k:M\pi ^kL\}k_tv(b_{n_tn_t})`$, so that 4.2 implies $`\text{dist}(\{M\},\{L\})v(b_{n_tn_t})`$. Hence we are done if we show that $$\text{dist}(\{M\},\{L_1+\mathrm{}+L_d\})v(b_{n_tn_t}).$$ Recall that $`BA`$ is the matrix consisting of the diagonal blocks $`B_{11}A_{11},\mathrm{},B_{dd}A_{dd}`$. Now for all $`h=1,\mathrm{},d`$ the determinant of $`A_{hh}`$ is $`_{j=n_{h1}+1}^{n_h}a_{jj}`$, and $`v(detB_{hh})=v(b_{n_hn_h})`$. Applying again (3) and the fact that in every line in $`B`$ there must be a unit, we find $$v(b_{n_hn_h})v\left(\underset{j=n_h+1}{\overset{n}{}}a_{jj}\right),$$ so that $$v(detB_{hh}A_{hh})v\left(\underset{j=n_h+1}{\overset{n}{}}a_{jj}\right)+v\left(\underset{j=n_{h1}+1}{\overset{n_h}{}}a_{jj}\right)v(detA)=v(b_{n_tn_t}),$$ hence $`\pi ^{v(b_{n_tn_t})}(A_{hh}^1B_{hh}^1)`$ has $`R`$-coefficients for all $`h`$. Therefore $`\pi ^{v(b_{n_tn_t})}B^1`$ has $`R`$-coefficients, which implies $`M\pi ^{v(b_{n_tn_t})}(L_1+\mathrm{}+L_d)`$, as $`{}_{}{}^{t}BM=L_1+\mathrm{}+L_d`$. Now it is easy to see that indeed $`\text{dist}(\{M\},\{L_1+\mathrm{}+L_d\})v(b_{n_tn_t})`$. $`\mathrm{}`$ One may wonder if our assumption that the linear cycles meet properly on the whole of $`(M)`$ is really necessary. In fact, in 3.4 we have proven a formula for the intersection number of $`n`$ hyperplanes meeting properly only on $`(V)`$. Hence it is tempting to try and use this as a starting point to generalize Theorem 5.1 at least to the case of $`n`$ hyperplanes $`H_1,\mathrm{},H_n`$ in $`(M)`$ meeting properly only on the generic fibre $`(V)`$. Unfortunately, the geometric expression $`\text{dist}(\{M\},A(H_{1K},\mathrm{},H_{nK}))`$ does in general no longer coincide with the intersection number $`<H_1,\mathrm{},H_n>`$ in this case. In fact, one can show as in the proof of 5.1 that for $`H_1,\mathrm{},H_n`$ with coordinate matrix $`A=(a_{ij})`$ (upper triangular and subject to $`v(a_{ii})v(a_{i+1i+1})`$ and $`v(a_{ii})v(a_{ij})`$ for $`i<j`$) we have $`\text{dist}(\{M\},A(H_{1K},\mathrm{},H_{nK}))=v(a_{nn})`$, whereas the intersection number $`<H_1,\mathrm{},H_n>`$ $`=_{i=1}^nv(a_{ii})`$ (see 3.4) may be bigger. We will conclude this paper by generalizing Theorem 5.1 in another direction. Namely, let us consider the case of $`d2`$ linear cycles on $`(M)`$ meeting properly in a cycle of dimension bigger than zero. So let $`N_1,\mathrm{},N_d`$ be split submodules of $`M`$ of rank $`r_1\mathrm{},r_d`$, so that $`(N_1),\mathrm{},(N_d)`$ are linear cycles meeting properly on $`(M)`$, i.e. any irreducible component in the intersection $`_i(N_i)`$ has codimension $$\underset{i}{}\text{codim}(N_i)=dn\underset{i=1}{\overset{d}{}}r_i.$$ Note that only the case $`_i\text{codim}(N_i)n`$ is interesting, since otherwise $`(N_i)`$ must be empty. Since we already dealt with the case $`_i\text{codim}(N_i)=n`$ in 5.1, let us now assume that $`_i\text{codim}(N_i)<n`$. Put $`r_0=n(dn_ir_i)=_ir_i(d1)n>0`$, and let $`L_0`$ be the intersection $`L_0=_{i=1}^dN_i`$. Since all $`M/N_i`$ are torsion free, the same holds for $`M/L_0`$, so that $`L_0`$ is a split submodule of $`M`$. Hence $`(L_0)`$ is a linear cycle contained in the intersection $`_i(N_i)`$, so it has codimension $`nr_0`$, which implies that $`\text{rk}L_0r_0`$. On the other hand, we can calculate $$\text{rk}L_0=dim(N_{iK})\underset{i=1}{\overset{d}{}}dim(N_{iK})(d1)n=\underset{i=1}{\overset{d}{}}r_i(d1)n=r_0,$$ so that we find $`\text{rk}L_0=r_0`$. For $`j=1,\mathrm{}d`$ we put $`L_j=_{ij}N_iM`$. Then $`L_0`$ is also a split submodule of $`L_j`$, hence there exists a free $`R`$-module $`L_j^{}`$ with $`L_j=L_0L_j^{}`$. Let $`F`$ be the set of vertices in $`X`$ of the form $$\{\pi ^{k_1}L_1^{}\mathrm{}\pi ^{k_d}L_d^{}\pi ^{k_{d+1}}L_0\}\text{for some }k_1,\mathrm{},k_{d+1}.$$ Of course, $`F`$ depends on the choice of $`L_1^{},\mathrm{},L_d^{}`$. It is easy to see that the sum $`\pi ^{k_1}L_1^{}+\mathrm{}+\pi ^{k_d}L_d^{}+\pi ^{k_{d+1}}L_0`$ is indeed direct. Since $$dimL_{iK}\underset{ji}{}dimN_{jK}(d2)n=\underset{ji}{}r_j(d2)n=nr_i+r_0$$ we find that $`{\displaystyle \underset{i=1}{\overset{d}{}}}dimL_{iK}^{}+dimL_{0K}=`$ $`{\displaystyle \underset{i=1}{\overset{d}{}}}dimL_{iK}(d1)r_0dn{\displaystyle \underset{i=1}{\overset{d}{}}}r_i+r_0=n,`$ so that $`\pi ^{k_1}L_1^{}\mathrm{}\pi ^{k_d}L_d^{}\pi ^{k_{d+1}}L_0`$ is indeed a lattice of full rank in $`V`$. Besides we see that $`\text{rk}L_i^{}=nr_i`$. Then the strategy of the proof of 5.1 can be modified to prove the following result: ###### Theorem 5.2 Let $`(N_1),\mathrm{},(N_d)`$ be linear cycles meeting properly on $`(M)`$ such that $`_{i=1}^d\text{codim}(N_i)<n`$. For any choice of $`F`$ as above, we can describe the Serre intersection cycle $`(N_1)\mathrm{}(N_d)`$ as follows: $$(N_1)\mathrm{}(N_d)=\overline{(_{i=1}^dN_{iK})}+\text{dist}(\{M\},F)(_{i=1}^dN_{ik}).$$ Here the first cycle $`\overline{(_{i=1}^dN_{iK})}`$ is the closure in $`(M)`$ of the linear cycle $`(_{i=1}^dN_{iK})`$ on $`(V)`$, which is just the Serre intersection cycle $`(N_1)_K\mathrm{}(N_d)_K`$ of the generic fibres. The second cycle $`(_{i=1}^dN_{ik})`$ is the linear cycle on the special fibre corresponding to the subspace $`_{i=1}^dN_{ik}`$ of $`M_k`$, which we regard as a cycle on $`(M)`$. Both our results 5.1 and 5.2 are proven by a complicated series of direct computations once the correct formulas are found. It would be desirable to give a more conceptual proof providing deeper insights in the nature of our formulas and allowing generalizations.
warning/0004/cond-mat0004481.html	ar5iv	text	# Spatial Inhomogeneities in Disordered d-Wave Superconductors: Effect on Density of States and Superfluid Stiffness Recently there has been a great deal of theoretical and experimental interest in the problem of disordered d-wave superconductors (SC) given the d-wave symmetry of the high Tc cuprates. A single unitary scatterer was predicted by Balatsky and coworkers to lead to a low-energy resonance with a characteristic four-fold symmetric wavefunction about the impurity site. This was recently observed in an STM study of a Zn-doped cuprate. The $`T`$-matrix approximation used in Ref. is very accurate for the one-impurity problem, and the order parameter suppression near the impurity, which it neglects, does not lead to any qualitative changes. On the other hand, the problem of a finite density of unitary scatterers is more subtle . There is a large body of theoretical work using the self-consistent $`T`$-matrix approximation leading to very interesting predictions . However, the impurity averaging procedure used in this method, and in other approaches that go beyond it, basically washes out the inhomogeneous structures that the disorder potential gives rise to. We use the Bogoliubov-deGennes (BdG) approach to study the spatial inhomogeneity induced by unitary scatterers in a short coherence length ($`\xi _0`$) superconductor, and to understand how these affect the low energy properties of the system such as the one-particle density of states (DOS) $`N(\omega )`$ and the superfluid stiffness $`D_s`$. Our main results can be summarized as follows: (1) The low energy DOS is considerably reduced relative to the T-matrix result. The low lying excitations are found to be generated by the interference of individual impurity resonances. Such excitations, with their nontrivial spatial structure, cannot be adequately described within the T-matrix formulation. (2) The superfluid stiffness $`D_s`$ is found to be significantly larger than that obtained within the T-matrix analysis. The larger $`D_s`$ directly correlates with the lower DOS for “normal” excitations. (3) We find that off-diagonal long range order (ODLRO) and finite superfluid stiffness survive to impurity concentrations much higher than the critical concentration of the T-matrix approximation. This relative insensitivity of short coherence length d-wave superconductors to impurities is shown to be closely tied to the inhomogeneity of the pairing amplitude on the scale of $`\xi _0`$ in response to a random potential. In contrast, the T-matrix approach assumes a uniform amplitude which then gets globally suppressed to zero at a critical disorder. Several authors have previously used the BdG approach for dirty d-wave systems. Tc reduction, superfluid density and localization of excitations was studied in Ref. and more recently the density of states has been studied in refs. . Our calculations differ from these in several aspects, some more important (e.g., working at fixed density rather than fixed chemical potential and inclusion of inhomogeneous Hartree-Fock shifts) than others (such as choice of Hamiltonian, parameters, and particle-hole asymmetry). While our results are broadly consistent with those obtained previously, what is new here is our emphasis on understanding the BdG results for $`N(\omega )`$, $`D_s`$ and ODLRO in terms of two different effects: (A) the inhomogeneity in the local pairing amplitude which characterizes the disordered ground state, and (B) the spatial structures characterizing the low-lying excitations in the disordered system. This provides a deeper insight into our BdG results, and also highlights the shortcomings of the $`T`$-matrix approach. We model the 2D disordered d-wave SC by the Hamiltonian $`=𝒦+_{\mathrm{int}}+_{\mathrm{dis}}`$. The kinetic energy $`𝒦=t_{<ij>,\alpha }(c_{i\alpha }^{}c_{j\alpha }+h.c.)`$ describes electrons, with spin $`\alpha `$ at site $`i`$ created by $`c_{i\alpha }^{}`$, hopping between nearest-neighbors $`<ij>`$ on a square lattice. The interaction term $`_{\mathrm{int}}=J_{<ij>}\left(𝐒_i𝐒_jn_in_j/4\right)+U_in_in_i`$ is chosen to lead to a $`d`$-wave SC ground state in the disorder-free system. The spin operator $`S_i^a=c_{i\alpha }^{}\sigma _{\alpha \beta }^ac_{i\beta }`$, where the $`\sigma ^a`$ are Pauli matrices, and the density $`n_{i\alpha }=c_{i\alpha }^{}c_{i\alpha }`$ with $`n_i=n_i+n_i`$. Finally, $`_{\mathrm{dis}}=_i\left(V(i)\mu \right)n_i`$ where $`\mu `$ is the chemical potential and the disorder potential $`V(i)`$ is an independent random variable at each site which is either $`+V_0`$, with a probability $`n_{\mathrm{imp}}`$ (impurity concentration), or zero. We believe that such a simple model is adequate to describe the strongly-correlated cuprates at low temperatures because their SC state has sharp quasiparticle excitations. The BdG equations are given by: $$\left(\begin{array}{cc}\widehat{\xi }& \widehat{\mathrm{\Delta }}\\ \widehat{\mathrm{\Delta }}^{}& \widehat{\xi }^{}\end{array}\right)\left(\begin{array}{c}u_n\\ v_n\end{array}\right)=E_n\left(\begin{array}{c}u_n\\ v_n\end{array}\right)$$ (1) where $`\widehat{\xi }u_n(j)=_\delta (t+W_j)u_n(j+\delta )+(V(j)\stackrel{~}{\mu }_j)u_n(j)`$ and $`\widehat{\mathrm{\Delta }}u_n(j)=_\delta \mathrm{\Delta }(j+\delta ;\delta )u_n(j+\delta )`$, and similarly for $`v_n(j)`$. The pairing amplitude on a bond $`(j;\delta )`$, where $`\delta =\pm \widehat{𝐱},\pm \widehat{𝐲}`$, is defined by $`\mathrm{\Delta }(j;\delta )=Jc_{j+\delta }c_j+c_jc_{j+\delta }/2`$. The inhomogeneous Hartee-Fock shifts are given by $`\stackrel{~}{\mu }_j=\mu Un_j/2+J_\delta n_{j+\delta }`$ and $`W_j=Jc_{j,\alpha }^{}c_{j+\delta ,\alpha }`$ We numerically solve for the BdG eigenvalues $`E_n0`$ and eigenvectors $`(u_n,v_n)`$ on a lattice of $`N`$ sites with periodic boundary conditions. We then calculate the pairing amplitude $`\mathrm{\Delta }(j;\delta )=J_n\left[u_n(j+\delta )v_n^{}(j)+u_n(j)v_n^{}(j+\delta )\right]/2`$ at $`T=0`$, the density $`n_j=2_n\|v_n(j)\|^2`$, and Fock shift $`W_j=J_nv_n(j+\delta )v_n^{}(j)`$. These are fed back into the BdG equation, and the process iterated until self consistency is achieved for each of the (local) variables defined on the sites and bonds of the lattice. The chemical potential $`\mu `$ is chosen to obtain a given average density $`n=_in_i/N`$, The d-wave pairing amplitude is given by $`\mathrm{\Delta }(j)=\left[\mathrm{\Delta }(j;+\widehat{x})\mathrm{\Delta }(j;+\widehat{y})+\mathrm{\Delta }(j;\widehat{x})\mathrm{\Delta }(j;\widehat{y})\right]/4`$. We have studied the model for a range of parameters and lattice sizes. Here we focus on $`J=U=1.15`$, in units of $`t=1`$, with $`n=0.875`$ (similar to the parameters used in refs. ) on systems of size up to $`26\times 26`$. For these parameters, and $`n_{\mathrm{imp}}=0`$, the DOS $`N_00.21`$ and $`\mathrm{\Delta }_00.077`$ corresponding to a maximum gap of $`0.31`$. For the impurity potential we choose $`V_0=100`$, close to the unitary limit. The results are averaged over 15 - 40 different realizations of the random potential. Let us first study the density of states (DOS) $`N(\omega )=\frac{1}{N}_{n,i}\left[\|u_n(i)\|^2\delta (\omega E_n)+\|v_n(i)\|^2\delta (\omega +E_n)\right]`$ (where we broaden the delta functions with a width comparable to average level spacing). In Fig. 1 we plot $`N(\omega )`$ for several impurity concentrations on a small energy scale; for comparison, the maximum energy gap in the disorder-free system is $`0.31`$ and the T-matrix self energy scale $`\gamma =\sqrt{n_{\mathrm{imp}}\mathrm{\Delta }/2N_0}0.25`$ for the parameters chosen; ($`\mathrm{\Delta }`$ is the T-matrix gap). In the T-matrix theory $`N(\omega )`$ is a constant for $`\omega \gamma `$, while we find a sharp dip in the DOS close to the chemical potential, consistent with ref. . In fact, we found $`N(0)=0`$ for each impurity configuration at every concentration that we studied. The scale of the sharp dip at finite $`n_{\mathrm{imp}}`$ was found to be the same as the energy of an isolated impurity resonance. It is very clear that the low energy DOS in the BdG calculations is considerably smaller than that in the T-matrix approximation (even though we do not have the spectral resolution to quantify the asymptotic form of $`N(\omega )`$ as $`\omega 0`$). To highlight this, we compare in Fig. 2 the finite $`N(0)`$ of the T-matrix analysis with the BdG $`\overline{N}(0)`$, which is the average of $`N(\omega )`$ over the (arbitrarily chosen) range $`\|\omega \|0.05\gamma `$. To gain further insight into this difference between the T-matrix and BdG results, we study the wavefunctions of the low-lying excitations for individual disorder realizations. The probability density $`\|u_n(i)\|^2+\|v_n(i)\|^2`$ corresponding to the lowest energy states at various impurity concentrations are plotted in right hand panels of Fig. 3. The resonance for a single unitary impurity shows characteristic powerlaw tails along diagonal directions . From Fig. 3, and other low lying excitations not shown here, we see that for finite $`n_{\mathrm{imp}}`$ these wave functions are generated by the hybridization of individual impurity resonances. The effects of constructive and destructive interference between the “diagonal tails” of individual resonances are apparent. The importance of such states was suggested in ref. ; however, their analysis assumed that the resonance energies are randomly distributed over a scale $`W\mathrm{\Delta }_0`$, which is not the case in the physical situation obtained here. We emphasize that excitations with such non-trivial spatial structures cannot be described by T-matrix theory, which treats the scattering of quasiparticles in a homogeneous (impurity averaged) medium off a single impurity in a self-consistent fashion. The resulting constant $`N(0)`$ then arises from a constant broadening $`\gamma `$ (defined above) of states near the d-wave nodes. In contrast, the low energy DOS in the BdG theory comes from new states arising out of hybridization of impurity resonances. We already see from Fig. 2 that at and beyond the critical concentration of the T-matrix approach, $`n_{\mathrm{imp}}^c0.08`$ for our choice of parameters, the BdG DOS does not approach the non-disordered value $`N_0`$. This raises the questions: does SC persist beyond $`n_{\mathrm{imp}}^c`$, and if so, how? To address these issues we calculate the superfluid stiffness using the linear response result: $`D_s/\pi =k_x\mathrm{\Lambda }_{xx}(q_x=0,q_y0,\omega =0)`$. The diamagnetic term $`k_x`$, is one-half (in 2D) the kinetic energy $`𝒦`$, and the paramagnetic term $`\mathrm{\Lambda }_{xx}`$, is the long wavelength limit of the transverse current-current correlation averaged over disorder realizations. We see from Fig. 4b that the superfluid stiffness $`D_s`$ is much larger than the T-matrix result, consistent with Ref. , and does not vanish up to $`n_{\mathrm{imp}}=0.12`$ which is $`50\%`$ larger than $`n_{\mathrm{imp}}^c`$ within the T-matrix approximation. (We did not go to higher impurity concentrations because of the increase in computational time to reach self-consistency.) In any case, we expect that once $`D_s`$ is sufficiently small, phase fluctuations neglected within the BdG mean field approach will drive the transition to the non-superconducting state ; this is left for a future investigation. Here we wish to gain insight into how the system manages to exhibit $`D_s>0`$, even when T-matrix theory predicts it to be non-superconducting. One way to think about this is to correlate $`D_s`$ and $`N(\omega )`$. A smaller DOS for low-lying excitations in the BdG approach implies fewer “normal fluid” excitations and hence a larger superfluid density compared to the T-matrix approximation. A complementary approach, which we find very illuminating, relates the $`D_s`$ to the inhomogeneous pairing amplitude $`\mathrm{\Delta }(i)`$ in the disordered ground state, shown in the left panels of Fig. 3. Notice that the d-wave pairing amplitude is suppressed in the vicinity of an impurity on the scale of the coherence length $`\xi _0`$ which is 3 to 4 lattice units. (In addition, a small extended s-wave component, not shown, also develops nearby). The regions of suppressed pairing amplitude give the appearance of “swiss cheese” at finite $`n_{\mathrm{imp}}`$ in Fig. 3. In the T-matrix approach the order parameter is forced to be spatially uniform and it vanishes for $`n_{\mathrm{imp}}n_{\mathrm{imp}}^c`$. However, by allowing the pairing amplitude to vary on the scale of $`\xi _0`$, in response to the impurity potential, the BdG solution permits a non-vanishing order parameter $`\overline{\mathrm{\Delta }}`$ which is larger than that obtained within T-matrix theory for all $`n_{\mathrm{imp}}`$; see Fig. 4a. ($`\overline{\mathrm{\Delta }}`$ is formally defined in terms of the long distance behavior of the appropriate reduced two-particle density matrix). We note that both $`\overline{\mathrm{\Delta }}/\mathrm{\Delta }_0`$ and $`D_s/D_{s,0}`$ are linear functions of $`n_{\mathrm{imp}}\xi _0^2`$ for a substantial range of impurity concentration. To qualitatively understand the superfluid stiffness $`D_s`$ consider applying an external phase twist to the inhomogeneous ground state. Despite the fact that at large $`n_{\mathrm{imp}}`$ there are large regions where the amplitude vanishes, there are still paths that permit phase information to be conveyed from one edge of the system to the other, thus leading to a non-vanishing $`D_s`$. Thus the spatial inhomogeneity of the pairing amplitude, which is particularly important in short coherence length superconductors, is crucial in understanding the relative insensitivity of the system to unitary impurities, in that the order parameter and superfluid stiffness are much larger than one might have guessed from the T-matrix approximation. This lack of sensitivity of the high Tc cuprates to disorder has been seen in numerous experiments . Despite the quantitative results on finite systems and their detailed qualitative understanding, many questions remain open. The first one relates to Tc suppression. While it is easy to calculate the “mean field Tc”, a more reliable estimate should include the effect of both phase fluctuations and quasiparticles. Another important question is thermal transport in the SC state. Why does it not reflect the low energy structure of the DOS and why is it consistent with the universal behavior predicted by T-matrix theory , when the superfluid density shows deviations from it. A full understanding of the asymptotic DOS of the low-energy excitations, their localization properties and the study of SC state transport on a network of hybridized resonances are all topics for future research. Acknowledgements: We would like to thank A. V. Balatsky, P. J. Hirschfeld, A. Paramekanti, S. H. Pan and G. P. Das for illuminating discussions. M. R. was supported in part by the Department of Science and Technology through the Swarnajayanti scheme.
warning/0004/hep-ph0004263.html	ar5iv	text	# Towards a solution of the charmonium production controversy: 𝑘_⟂-factorization versus color octet mechanism ## Abstract The cross section of $`\chi _{cJ}`$ hadroproduction is calculated in the $`k_{}`$ -factorization approach. We find a significant contribution of the $`\chi _{c1}`$ state due to non-applicability of the Landau-Yang theorem because of off-shell gluons. The results are in agreement with data and in contrast to the collinear factorization show a dominance of the color singlet part and a strong suppression of the color octet contribution. Our results could therefore lead to a solution of the longstanding controversy between the color singlet model and the color octet mechanism. 20.09.2000 The production of heavy quarkonia received a lot of attention from both theory and experiment in recent years. It is e.g. the most prominent signal in the search for the quark gluon plasma. Its usefullness is, however, questionable as long as the charmonium production process is not understood. For a review we refer to . Originally heavy quarkonium production was described in the color singlet model (CSM) . Calculations based on this model and standard collinear factorization show however disagreement with the experimental data. For example the next-to-leading order (NLO) QCD collinear results for direct $`J/\mathrm{\Psi }`$ hadroproduction underestimate the measured cross section at Tevatron by a factor of $``$ 50 (see fig.4 in and Ref.). The proposed solution to this strong discrepancy is the so called color-octet-mechanism (COM) , according to which a color octet $`q\overline{q}`$-pair which has been produced at short distances can evolve into a physical quarkonium state by radiating soft gluons. The COM introduces uncalculable non-perturbative parameters, the color octet matrix elements, which have to be determined by a fit to the data . The inclusion of the COM into NLO QCD collinear calculations leads in the case of hadroproduction to a reasonable agreement with experiment . In these calculations the color octet contribution dominates. On the other hand up to now the COM suffers at least from two unsolved problems. When the, supposedly universal, color octet matrix elements are applied to electroproduction of heavy quarkonium the theoretical predictions fail to describe the data . Furthermore the results of the COM for polarized heavy quarkonium hadroproduction seem to be incompatible with recent data from Tevatron . The longstanding discrepancy between the results based on the CSM together with collinear factorization and the experimental data shows up especially strongly in the $`k_{}`$-dependent cross sections from Tevatron . Thus one can wonder if the collinear approximation, in which in NLO the only transverse momentum of the produced quarkonium comes from an additional final state gluon, is suitable at all. The aim of our paper is to clarify this question by a study of $`\chi _{cJ}`$ hadroproduction within the $`k_{}`$ -factorization approach, which takes the nonvanishing transverse momenta of the colliding t-channel gluons into account. Generically this corresponds to taking into account new regions of the phase space of the colliding gluons which is mandatory for the description of hard processes in the Regge region. More precisely we calculate the production of $`J/\mathrm{\Psi }`$’s originating from radiative $`\chi _{cJ}`$ decays. In a recent study of open $`b\overline{b}`$ hadroproduction we found that $`k_{}`$ -factorization gives far better results than NLO collinear QCD calculations and we expect a similar improvement for heavy quarkonium production. The main ingredients of our calculations in are the unintegrated gluon distribution and the effective next-to-leading-logarithmic-approximation (NLLA) $`q\overline{q}`$ -BFKL production vertex which we use in this article as well. The projection of the heavy quark-antiquark pair onto the corresponding charmonium state is described in the standard way within the non-relativistic-quarkonium-model . We study the production of $`\chi _{cJ}`$ whose lowest Fock state component is $`q\overline{q}(^3P_J)`$. For $`J/\mathrm{\Psi }`$ (a $`q\overline{q}(^3S_1)`$ state) the LO production amplitude is zero. In order to get a nonzero $`q\overline{q}(^3S_1)`$-amplitude one has (in NLO in $`\alpha _S`$) to emit an additional gluon. The amplitude for the production of a $`q\overline{q}`$ -pair plus a gluon within the BFKL approach would in our case require an effective three-particle production vertex which still has to be derived. In contrast the production of a $`\chi _{c1}`$ can be calculated in our approach in LO because the Landau-Yang theorem which usually forbids the production of a $`{}_{}{}^{3}P_{1}^{}`$ state is not valid for off-mass-shell gluons. We use the following definition of the light cone coordinates $`k^+=k^0+k^3,\text{ }k^{}=k^0k^3,\text{ }k_{}=(0,k^1,k^2,0)=(0,𝐤,0).`$ In the c.m. frame the momenta of the scattering hadrons are given by $`P_1^+=P_2^{}=\sqrt{s},\text{ }P_1^{}=P_2^+=P_1=P_2=0,`$ where the Mandelstam variable $`s`$ is as usual the c.m.s. energy squared. The momenta of the t-channel gluons are $`q_1`$ and $`q_2`$ (see Fig.1). The on-shell quark and antiquark (with mass $`m`$) have momentum $`k_1`$ respectively $`k_2`$ with $`k_1^{}={\displaystyle \frac{(m^2k_1^2)}{k_1^+}},\text{ }k_2^{}={\displaystyle \frac{(m^2k_2^2)}{k_2^+}}.`$ In the high energy (large $`s`$) regime we have $`P^+=q_1^+q_2^+q_1^+,P^{}=q_1^{}q_2^{}q_2^{},q_{1/2}^2q_{1/2}^2,`$ where $`P=k_1+k_2`$ is the momentum of the heavy quarkonium with $`P^2=4m^2`$. The longitudinal momentum fractions of the gluons are $`x_1=q_1^+/P_1^+`$, $`x_2=q_2^{}/P_2^{}`$. The heavy quarkonium hadroproduction cross section in the $`k_{}`$-factorization approach is , $`\sigma _{P_1P_2\chi X}={\displaystyle \frac{1}{8(2\pi )}}{\displaystyle \frac{d^3P}{P^+}d^2q_1d^2q_2\delta ^2(q_1q_2P_{})}`$ (1) $`(x_1,q_1){\displaystyle \frac{1}{(q_1^2)^2}}\left\{{\displaystyle \frac{\psi _\chi ^{c_2c_1}\psi _\chi ^{c_2c_1}}{(N_C^21)^2}}\right\}{\displaystyle \frac{1}{(q_2^2)^2}}(x_2,q_2).`$ (2) The factor $`(N_C^21)^2`$ comes from the projection on color singlet in the t-channel. $`(x,q_{})`$ is the unintegrated gluon distribution. The heavy quarkonium production amplitude $`\psi _\chi ^{c_2c_1}(x_1,x_2,q_1,q_2,P)`$ is factorized (see below) in a hard part which describes the production of the $`q\overline{q}`$ pair and an amplitude describing the binding of this pair into a physical charmonium state. We choose the scale $`\mu ^2`$ for $`\alpha _S(\mu ^2)`$ in the amplitude $`\psi _\chi ^{c_2c_1}`$ to be $`𝐪_1^2=q_1^2`$ respectively $`𝐪_2^2=q_2^2`$ . The amplitude for the production of the charmonium state can be written as $`\psi _\chi ^{c_2c_1}=𝒫(q\overline{q}\chi _{cJ})\mathrm{\Psi }^{c_2c_1}.`$ (3) The $`q\overline{q}`$ production vertex $`\mathrm{\Psi }^{c_2c_1}`$ derived in for massless QCD, appropriately generalized for massive quarks, has the form $`\mathrm{\Psi }^{c_2c_1}=g^2\left(t^{c_1}t^{c_2}b(k_1,k_2)t^{c_2}t^{c_1}b^T(k_2,k_1)\right),`$ where $`t^c`$ are the colour group generators in the fundamental representation. The operator $`𝒫(q\overline{q}\chi _{cJ})`$ projects the $`q\overline{q}`$ pair onto the charmonium bound state, see below. The functions $`b(k_1,k_2)`$ and $`b^T(k_2,k_1)`$ are illustrated in Fig.2 and their explicit form can be found in . One important property of the charmonium production amplitude for on-mass-shell quark and antiquark states (3), which is related to the gauge invariance of the whole approach, is its vanishing in the limit $`q_10`$ (or $`q_20`$). The relation between the usual gluon distribution $`xg(x,𝐪^2)`$ and the unintegrated gluon distribution $`(x,𝐤)`$ is given by $`xg(x,𝐪^2)={\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{d𝐤^2}{𝐤^2}}\mathrm{\Theta }(𝐪^2𝐤^2)(x,𝐤).`$ (4) $`(x,𝐤)`$ includes the evolution in $`x`$ and $`𝐤^2`$ described by the BFKL and DGLAP equation. In the non-perturbative region of small $`𝐤^2`$ the unintegrated gluon distribution is not known, therefore we write (4) according to as $`xg(x,𝐪^2)=xg(x,𝐪_0^2)+{\displaystyle _{𝐪_0^2}^{\mathrm{}}}{\displaystyle \frac{d𝐤^2}{𝐤^2}}\mathrm{\Theta }(𝐪^2𝐤^2)(x,𝐤),`$ which introduces the a priori unknown initial scale $`𝐪_0`$ and the initial gluon distribution $`xg(x,𝐪_0^2)`$. Following , we neglect the momentum dependence of the hard cross section in the soft region $`\|𝐪\|<\|𝐪_0\|`$, so that $`{\displaystyle \frac{1}{q_1^2}}\left\{{\displaystyle \frac{\psi _\chi ^{c_2c_1}\psi _\chi ^{c_2c_1}}{(N_C^21)^2}}\right\}{\displaystyle \frac{1}{q_2^2}}S(q_1,q_2)`$ $`\left[S(q_1,q_2)\mathrm{\Theta }(𝐪_2^2𝐪_0^2)+S(q_1,0)\mathrm{\Theta }(𝐪_0^2𝐪_2^2)\right]\mathrm{\Theta }(𝐪_1^2𝐪_0^2)`$ $`+\left[S(0,q_2)\mathrm{\Theta }(𝐪_2^2𝐪_0^2)+S(0,0)\mathrm{\Theta }(𝐪_0^2𝐪_2^2)\right]\mathrm{\Theta }(𝐪_0^2𝐪_1^2),`$ see also the discussion of this expression in . One important point is the proper choice of the unintegrated gluon distribution function. We use the results of Kwiecinski, Martin and Stas̀to . They determined it using a combination of DGLAP and BFKL evolution equations. With the initial conditions $`𝐪_0^2=1\text{ GeV},xg(x,𝐪_0^2)=1.57(1x)^{2.5}.`$ (5) they obtained an execellent fit to $`F_2(x,Q^2)`$ data over a large range of $`x`$ and $`Q^2`$. In order to see the effect of off-shell gluons and the inapplicability of the Landau-Yang theorem as well as to perform calculations which do not require a fit to the data we start with calculation of the color singlet part of the amplitude. This is most easily done by adapting the method of . The projection of the hard amplitude onto the charmonium bound state is given by $`\psi _\chi ^{c_2c_1}=𝒫(q\overline{q}\chi _{cJ})\mathrm{\Psi }^{c_2c_1}`$ (8) $`={\displaystyle \underset{i,j}{}}{\displaystyle \underset{L_z,S_z}{}}{\displaystyle \frac{1}{\sqrt{m}}}{\displaystyle \frac{d^4q}{(2\pi )^4}\delta \left(q^0\frac{\stackrel{}{q}^2}{M}\right)\mathrm{\Phi }_{L=1,L_z}(\stackrel{}{q})}`$ $`L=1,L_z,S=1,S_z\|J,J_z3i,\overline{3}j\|1Tr\left\{\mathrm{\Psi }_{ij}^{c_2c_1}𝒫_{S=1,S_z}\right\},`$ (9) where $`\mathrm{\Phi }_{L=1,L_z}(\stackrel{}{q}=\stackrel{}{k}_1\stackrel{}{k}_2)`$ is the momentum space wave function of the charmonium, and the projection operator $`𝒫_{S=1,S_z}`$ for a small relative momentum $`q=k_1k_2`$ has the form $`𝒫_{S=1,S_z}={\displaystyle \frac{1}{2m}}(\mathit{}_2m){\displaystyle \frac{\mathit{ϵ̸}(S_z)}{\sqrt{2}}}(\mathit{}_1+m).`$ The Clebsch-Gordan coefficient in color space is given by $`3i,\overline{3}j\|1=\delta _{ji}/\sqrt{N_C}`$. Since $`P`$-waves vanish at the origin, one has to expand the trace in (9) in a Taylor series around $`\stackrel{}{q}=0`$. This yields an expression proportional to $`{\displaystyle \frac{d^3\stackrel{}{q}}{(2\pi )^3}q^\alpha \mathrm{\Phi }_{L=1,L_z}(\stackrel{}{q})}=i\sqrt{{\displaystyle \frac{3}{4\pi }}}ϵ^\alpha (L_z)^{}(0),`$ with the derivative of the $`P`$-wave radial wave function at the origin $`^{}(0)`$ whose numerical values can be found in . For the individual $`\chi _{cJ=1}`$ and $`\chi _{cJ=2}`$ amplitudes we use $`{\displaystyle \underset{L_z,S_z}{}}1,L_z,1,S_z\|1,J_zϵ^\mu (L_z)ϵ^\nu (S_z)=i\sqrt{{\displaystyle \frac{1}{2}}}\epsilon ^{\mu \nu \alpha \beta }{\displaystyle \frac{P_\alpha }{M}}ϵ_\beta (J_z)`$ $`{\displaystyle \underset{L_z,S_z}{}}1,L_z,1,S_z\|2,J_zϵ^\mu (L_z)ϵ^\nu (S_z)=ϵ^{\mu \nu }(J_z)`$ where we introduce the spin 1 and spin 2 polarization tensors $`ϵ_\beta (J_z)`$ and $`ϵ^{\mu \nu }(J_z)`$ of the produced charmonium $`\chi _{cJ=1}`$ respectively $`\chi _{cJ=2}`$. In the unpolarized case the squared amplitudes are further evaluated using $`{\displaystyle \underset{J_z}{}}ϵ^\mu (J_z)ϵ^\nu (J_z)`$ $`=`$ $`g^{\mu \nu }+{\displaystyle \frac{P^\mu P^\nu }{M^2}}=P^{\mu \nu },`$ $`{\displaystyle \underset{J_z}{}}ϵ^{\mu \nu }(J_z)ϵ^{\alpha \beta }(J_z)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(P^{\mu \alpha }P^{\nu \beta }+P^{\nu \alpha }P^{\mu \beta }\right){\displaystyle \frac{1}{3}}P^{\mu \nu }P^{\alpha \beta }.`$ The cross section for $`J/\mathrm{\Psi }`$ production from radiative $`\chi _{cJ}`$ decays is then given by $`\sigma _{J/\mathrm{\Psi }from\chi _c}={\displaystyle \underset{J=0,1,2}{}}\sigma _{P_1P_2\chi _{cJ}X}Br(\chi _{cJ}J/\mathrm{\Psi }+\gamma ),`$ with the $`\chi _{cJ}`$ hadroproduction cross section $`\sigma _{P_1P_2\chi _{cJ}X}`$ (2). Because of the small branching ratio $`Br(\chi _{cJ=0}J/\mathrm{\Psi }+\gamma )=𝒪(10^3)`$ the contribution from $`\chi _{cJ=0}`$ is negligible. For the numerical computation we use the values $`m_c=1.48\text{ GeV,}\|^{}(0)\|^2=0.075\text{ GeV}^5.`$ The pseudorapidity is defined as $`\eta =\frac{1}{2}\mathrm{ln}\left((\sqrt{P_0^2M^2}+P_3)/(\sqrt{P_0^2M^2}P_3)\right)`$. To compare with data we multiply our cross sections with the braching ratio $`Br(J/\mathrm{\Psi }\mu ^+\mu ^{})`$. The resulting $`P_{}`$-dependent cross section for $`J/\mathrm{\Psi }`$’s from radiative deacays of $`\chi _c`$’s produced in $`pp`$-collisions is shown in Fig.3 together with the data from the CDF Collaboration and a NLO QCD collinear result (see Fig.7 in ). The individual contributions from $`\chi _{c1}`$ and $`\chi _{c2}`$ are shown in Fig.4. The description of the data by the color singlet part alone is very satisfactory and becomes even better if the difference of the transverse momentum of $`J/\mathrm{\Psi }`$ (which is measured experimentally) and $`\chi _c`$ (which enters our calculation) is taken into account. (Due to the radiative decay the transverse momentum of $`J/\mathrm{\Psi }`$ is typically larger by an amount of $`300`$ MeV than the corresponding $`\chi _c`$ one which leads to a shift of the theoretical curve to the right.) The typical scale of the gluon off-shellness is given by the transverse momentum of the produced quarkonium. We emphasize that the result has been obtained without fitting any of the parameters involved: The unintegrated gluon distribution has been adopted from Kwiecinski et al. . The parameters of the quarkonium bound state are the ones given by Eichten and Quigg . For the $`\chi _{c1}`$ state it is crucial that the gluons are off-shell in $`k_{}`$ -factorization. Now we proceed with the calculation of $`J/\mathrm{\Psi }`$ production by $`\chi _c`$ radiative decays adopting the colour octet mechanism. The infrared stability of higher order corrections to the cross section requires the existence of a color octet contribution, without fixing its size . The $`\chi _c`$ state can be written in a velocity expansion as $`\|\chi _c=𝒪(1)\|q\overline{q}\left[{}_{}{}^{3}P_{J}^{\text{ }1}\right]+𝒪(𝐯)\|q\overline{q}\left[{}_{}{}^{3}S_{1}^{8}\right]g+\mathrm{}.`$ Following the formalism of the resulting cross section is then proportional to the color octet matrix element $`0\left\|O_8^{\chi _{c1}}(^3S_1)\right\|0`$ which has to be fitted to data. Using the results for the color singlet part and adding the color octet contribution we obtain as value of the color octet matrix element $`0\left\|O_8^{\chi _{c1}}(^3S_1)\right\|0=(9.0\pm 2.0)\times 10^4.`$ Comparing this with the result obtained in the collinear factorization we find a suppression of the matrix element due to the flat $`P_{}`$-dependence of the color octet contribution by roughly one order of magnitude, resulting in a violation of the velocity scaling rules. These scaling rules are derived rigorously in the framework of non-relativistic QCD (NRQCD) . It is, therefore, natural to assume that the charm quark is simply not heavy enough for the velocity scaling rules of NRQCD to be valid. This is also suggested by other observations, see e.g. the very recent study . In contrast the description of bottom systems in NRQCD should be more accurate. This shows the importance of a detailed analysis of bottomonia production in the $`k_{}`$-factorization approach. Let us conclude. The $`k_{}`$-factorization approach relying on an unintegrated gluon distribution compatible with the small $`x`$ behaviour of the structure function $`F_2`$ together with the BFKL NLLA fermion production vertices describes correctly $`\chi _c`$ production in the central rapidity region. Whereas the standard collinear factorization approach in NLO can describe the data in the TeV range only by introducing a dominant octet contribution, we have shown that in the $`k_{}`$-factorization approach such a contribution gives an improved description of the data but is suppressed by its $`P_{}`$ behaviour. Our main conclusion is therefore that the correct way to improve the standard QCD calculations for quarkonium production in the TeV range is to abandon the collinear approximation. The contributions disregarded in the collinear approximation of strong transverse momentum ordering become essential in the small-$`x`$ range. The relative merits of the $`k_{}`$ -factorization as the standard approach for other processes in high energy hadronic collisions still has to be investigated. L.Sz. and O.V.T. thank A. Tkabladze for discussion. The work is supported by the DFG.
warning/0004/quant-ph0004035.html	ar5iv	text	# Collective versus local measurements on two parallel or antiparallel spins ## Abstract We give a complete analysis of covariant measurements on two spins. We consider the cases of two parallel and two antiparallel spins, and we consider both collective measurements on the two spins, and measurements which require only Local Quantum Operations and Classical Communication (LOCC). In all cases we obtain the optimal measurements for arbitrary fidelities. In particular we show that if the aim is determine as well as possible the direction in which the spins are pointing, it is best to carry out measurements on antiparallel spins (as already shown by Gisin and Popescu), second best to carry out measurements on parallel spins and worst to be restricted to LOCC measurements. If the the aim is to determine as well as possible a direction orthogonal to that in which the spins are pointing, it is best to carry out measurements on parallel spins, whereas measurements on antiparallel spins and LOCC measurements are both less good but equivalent. One of the central problems of quantum measurements is how to best estimate the state of an unknown quantum system. This problem has been addressed by many authors, using many different approaches, see for reviews. In the present paper we take a new look at a particular example in which the task is to determine the direction of polarization of two identical spin 1/2 particles. We suppose that the polarization direction is completely unknown, ie. is uniformly distributed on the sphere. This problem, generalized to the case of an arbitrary number N of spin 1/2 has already been studied by several authors. The particular case of 2 spins has the advantage of being sufficiently simple that a complete solution can be obtained. Furthermore it allows for several twists where particular features of quantum mechanics related to entanglement reveal themselves. The first twist on the original problem was suggested by Peres and Wootters who asked whether there is a difference between collective, as compared to local, measurements on two particles. Technically, in the first case one allows arbitrary quantum operations on both spins, whereas in the second case one restricts oneself to Local Quantum Operations on each particle and Classical Communication between the particles (LOCC). Peres and Wootters gave numerical evidence that even if two particles are in the same state, collective measurements can be better than measurements using LOCC. An analytical proof was given in in the case of two identical spin 1/2 whose polarization direction is uniformly distributed on the sphere, at least in the case where there are only a finite number of rounds of communication between the two parties. A remarkable example was exhibited in which consists of a basis of separable states, ie. states that can be prepared using LOCC, but which nevertheless cannot be distinguished unambiguously using LOCC. In the phrase “non locality without entanglement” was coined for this property. Another result is that of where it was shown that in the limit of an infinite number of identical spin 1/2, LOCC measurements perform as well as collective measurements if the spins are in a pure state, but not as well if the spins are in a mixed state. A second twist on the original scenario was recently proposed by Gisin and Popescu who considered the case of two antiparallel spins. They showed that one can determine better the direction of polarization of two antiparallel spins than of two parallel spins. Gisin and Popescu’s result is related to “non locality without entanglement” because if LOCC measurements are carried out on the two spins, it does not make any difference whether the spins are parallel or antiparallel. Thus collective measurements on two antiparallel spins is an example of non locality without entanglement. Mathematically the passage from two parallel to two antiparallel spins, that is the flip of one of the spins, is the same operation that Peres used to distinguish whether two states are entangled or not. The present paper aims at providing an exhaustive analysis of measurements on two spins 1/2 particles in the three cases of collective measurements on parallel spins, collective measurements on antiparallel spins, and LOCC measurements. The spin flip operation will play a central role in this analysis because it will allow us to treat all three cases in the same framework. Our analysis allows one to find the optimal measurements for arbitrary fidelities. As an illustration we consider two such fidelities. The first fidelity is $`f=(1+\mathrm{cos}\theta )/2`$ where $`\theta `$ is the angle between the direction in which the spins are polarized and the direction in which one guesses that they are polarized. In this case we recover the results of that if the spins are parallel the maximal average fidelity is $`f=0.75`$. If the spins are antiparallel we show that the maximal fidelity is $`f=0.788`$. This fidelity was already obtained in but it was not known whether it is optimal. Finally we shall show that if one restricts oneself to LOCC measurements, then the maximal fidelity is $`f=0.736`$, which is 1.4 % lower than for measurements on parallel spins. That this is the optimal value for LOCC measurements was already found by D.G. Fischer, S.H. Kienle, and M. Freyberger . Thus even in the limit of an infinite number of rounds of communication collective measurements are better than LOCC measurements in the case of two parallel spins. The second fidelity is $`f=1cos^2\theta `$. In this case it is most advantageous to guess a direction orthogonal to that in which the spins are pointing ($`\theta =\pi /2`$) and most disadvantageous to guess in the direction in which the spins are pointing ($`\theta =0`$) or in the orthogonal direction ($`\theta =\pi `$). Geometrically this can be rephrased as a situation in which the spins encode the orientation of a plane by pointing in the direction normal to the plane and the aim is to find a vector lying in the plane. In this case the highest fidelity $`f=0.8`$ is obtained when the spins are parallel. Antiparallel spins or LOCC measurements both give the same optimal fidelity $`f=0.733`$. We now turn to the proof of these results. Essential to our analysis will be the spin flip operation which we denote by $`\stackrel{~}{}`$. For a single spin 1/2 it takes the form $$\rho =\frac{I}{2}+\stackrel{}{\alpha }\stackrel{}{\sigma }\rho \stackrel{~}{}=\frac{I}{2}\stackrel{}{\alpha }\stackrel{}{\sigma }$$ (1) where $`I`$ is the identity operator and $`\sigma _i`$ the Pauli spin operators. In the case of two spins, we will be interested in the operation, denoted $`\stackrel{~}{}^2`$ which flips only the second spin. If we write the state as $$\rho =\frac{I}{4}+\stackrel{}{\alpha }\stackrel{}{\sigma }\frac{I}{2}+\stackrel{}{\beta }\frac{I}{2}\stackrel{}{\sigma }+\underset{i,j}{}\gamma _{ij}\sigma _i\sigma _j,$$ (2) then $`\rho \stackrel{~}{}^2`$ is given by $$\rho \stackrel{~}{}^2=\frac{I}{4}+\stackrel{}{\alpha }\stackrel{}{\sigma }\frac{I}{2}\stackrel{}{\beta }\frac{I}{2}\stackrel{}{\sigma }\underset{i,j}{}\gamma _{ij}\sigma _i\sigma _j$$ (3) The $`\stackrel{~}{}^2`$ operation is equivalent, up to a unitary operation acting on particle 2 only, to the partial transpose introduced in . As an application of the $`\stackrel{~}{}^2`$ operation consider the state of two parallel spin 1/2 particles both pointing in the $`\stackrel{}{m}`$ direction $`\rho (\stackrel{}{m},\stackrel{}{m})`$ $`=`$ $`\|_\stackrel{}{m}_\stackrel{}{m}\|\|_\stackrel{}{m}_\stackrel{}{m}\|`$ (4) $`=`$ $`{\displaystyle \frac{I}{4}}+{\displaystyle \underset{i}{}}m_i(\sigma _i{\displaystyle \frac{I}{2}}+{\displaystyle \frac{I}{2}}\sigma _i)+{\displaystyle \underset{i,j}{}}m_im_j\sigma _i\sigma _j`$ (5) and the state of two antiparallel spins $$\rho (\stackrel{}{m},\stackrel{}{m})=\|_\stackrel{}{m}_\stackrel{}{m}\|\|_\stackrel{}{m}_\stackrel{}{m}\|.$$ (7) We have the relation $$\rho (\stackrel{}{m},\stackrel{}{m})=\rho (\stackrel{}{m},\stackrel{}{m})\stackrel{~}{}^2.$$ (8) We can also consider the dual of the $`\stackrel{~}{}^2`$ operation, that is how it acts on operators. Suppose that $`\rho `$ is a state and $`a`$ an operator, then $`a\stackrel{~}{}^2`$ is defined by the relation $$Tra\rho \stackrel{~}{}^2=Tra\stackrel{~}{}^2\rho .$$ (9) One finds that it takes exactly the same form for operators as it does for states. If the operator $`a`$ is expressed as $$a=wI+\stackrel{}{x}\stackrel{}{\sigma }I+\stackrel{}{y}I\stackrel{}{\sigma }+\underset{i,j}{}z_{ij}\sigma _i\sigma _j,$$ (10) then the operator $`a\stackrel{~}{}^2`$ takes the form $$a\stackrel{~}{}^2=wI+\stackrel{}{x}\stackrel{}{\sigma }I\stackrel{}{y}I\stackrel{}{\sigma }\underset{i,j}{}z_{ij}\sigma _i\sigma _j.$$ (11) The $`\stackrel{~}{}^2`$ operation for operators allows us to put a restriction on the Positive Operator Valued Measures (POVM) acting on the space of two spin 1/2 particles that can be realized by local Quantum Operations and Classical Communication (LOCC). Indeed it was shown in that such a POVM, defined by its elements $`a_i0`$, $`_ia_i=1`$, must obey $`a_i\stackrel{~}{}^20`$ for all $`i`$. To proceed with the proof, consider a set of operators $`a_i`$ that sum to the identity $`_ia_i=1`$. We are interested in the following 3 positivity conditions on $`a_i`$: 1. $`a_i0`$. In this case the $`a_i`$ constitute a POVM. The probability of getting outcome $`i`$ if the state is $`\rho (\stackrel{}{m},\stackrel{}{m})`$ is $`P_{}(i\|\stackrel{}{m})=Tr\rho (\stackrel{}{m},\stackrel{}{m})a_i`$. 2. $`a_i\stackrel{~}{}^20`$. In this case the $`a_i\stackrel{~}{}^2`$ constitute a POVM. The probability of getting outcome $`i`$ if the state is $`\rho (\stackrel{}{m},\stackrel{}{m})`$ is $`P_{}(i\|\stackrel{}{m})=Tr\rho (\stackrel{}{m},\stackrel{}{m})a_i\stackrel{~}{}^2`$. Using equation (9) we have $`P_{}(i\|\stackrel{}{m})=Tr\rho (\stackrel{}{m},\stackrel{}{m})a_i`$. 3. $`a_i0`$ and $`a_i\stackrel{~}{}^20`$. In this case both $`a_i`$ and $`a_i\stackrel{~}{}^2`$ constitute a measurement which can be realized by LOCC. The probability $`P_{}(i\|\stackrel{}{m})=Tr\rho (\stackrel{}{m},\stackrel{}{m})a_i`$ of obtaining outcome $`i`$ if the spins are parallel and the measurement is $`a_i`$ equals the probability $`P_{}(i\|\stackrel{}{m})=Tr\rho (\stackrel{}{m},\stackrel{}{}m)a_i\stackrel{~}{}^2`$ of of obtaining outcome $`i`$ if the spins are antiparallel and the measurement is $`a_i\stackrel{~}{}^2`$. The equality of $`P_{}(i\|\stackrel{}{m})`$ and $`P_{}(i\|\stackrel{}{m})`$ shows that in this case there is no difference between making measurements on parallel and antiparallel spins. Thus the $`\stackrel{~}{}^2`$ operation relates measurements on parallel spins (given by $`P_{}(i\|\stackrel{}{m})`$), measurements on antiparallel spins (given by $`P_{}(i\|\stackrel{}{m})`$), and measurements that can be realized by LOCC. The central idea is that by using the $`\stackrel{~}{}^2`$ operation all these quantities can be expressed in terms of the same trace $`Tr\rho (\stackrel{}{m},\stackrel{}{m})a_i`$, but with operators $`a_i`$ which obey the different positivity conditions enumerated above. To further explicitise these different positivity conditions we shall suppose that the aim of the measurement is to distinguish along which direction the spins are pointing. We can then label the POVM elements $`a_\stackrel{}{n}`$ by the direction $`\stackrel{}{n}`$ along which one guesses the spins are pointing. Furthermore we shall suppose that the spins are polarized in a random direction uniformly distributed on the sphere. We can then, without loss of generality , suppose that we are dealing with covariant measurements, that is measurements for which the guessed direction $`\stackrel{}{n}`$ spans the whole sphere and which satisfy $$Tra_\stackrel{}{n}\rho (\stackrel{}{m},\stackrel{}{m})=Tra_{R(\stackrel{}{n})}\rho (R(\stackrel{}{m}),R(\stackrel{}{m}))$$ (12) where $`R`$ is an arbitrary rotation, ie. an element of $`SO(3)`$. Using the fact that $`\rho (R(\stackrel{}{m}),R(\stackrel{}{m}))=RR\rho (\stackrel{}{m},\stackrel{}{m})R^{}R^{}`$ where $`R`$ is the corresponding element of $`SU(2)`$, we have $$a_{R(\stackrel{}{n})}=RRa_\stackrel{}{n}R^{}R^{}.$$ (13) We can also without loss of generality suppose that the measurement is symmetric with respect to interchanges the two spins. This implies that $$a_\stackrel{}{n}=wI+\stackrel{}{x}(\stackrel{}{\sigma }I+I\stackrel{}{\sigma })+\underset{i,j}{}z_{ij}\sigma _i\sigma _j$$ (14) with $`z_{ij}`$ a symmetric matrix. The covariance condition (13) implies a considerable simplification on the coefficients $`w,\stackrel{}{x},z_{ij}`$ in (14). Consider the POVM element $`a_\stackrel{}{z}`$ corresponding to guessing the spins are polarized along the $`+z`$ direction. Let $`R_{\varphi ,z}`$ be a rotation of angle $`\varphi `$ around the $`z`$ axis. We have $`a_\stackrel{}{z}=R_{\varphi ,z}R_{\varphi ,z}a_\stackrel{}{z}R_{\varphi ,z}^{}R_{\varphi ,z}^{}`$ for all $`\varphi `$. Using (14), this implies that $`a_\stackrel{}{z}`$ has the form $`a_\stackrel{}{z}`$ $`=`$ $`wI+\alpha (\sigma _zI+I\sigma _z)+\beta \sigma _z\sigma _z`$ (16) $`+\gamma (\sigma _x\sigma _x+\sigma _y\sigma _y)`$ where $`\alpha `$, $`\beta `$, $`\gamma `$ are three real numbers. A final simplification results if we recall that the operators $`a_\stackrel{}{n}`$ must sum to the identity: $$𝑑\stackrel{}{n}a_\stackrel{}{n}=_{SU(2)}𝑑RRRa_{+\stackrel{}{z}}R^{}R^{}=I.$$ (17) Using (16) this implies that $`a_\stackrel{}{z}`$ $`=`$ $`I+\alpha (\sigma _zI+I\sigma _z)`$ (19) $`+\gamma (2\sigma _z\sigma _z\sigma _x\sigma _x\sigma _y\sigma _y)`$ which only depends on two parameters $`\alpha `$ and $`\gamma `$. It is now easy to compute the restriction on the two parameters $`\alpha `$ and $`\gamma `$ which result from each of the three positivity conditions enumerated above: $`a_\stackrel{}{n}0`$ $``$ $`\gamma 1,1+\alpha +\gamma /20,`$ (21) $`1\alpha +\gamma /20`$ $`a_\stackrel{}{n}\stackrel{~}{}^20,`$ $``$ $`\gamma 2,1+\gamma \alpha ^20,`$ (22) $`a_\stackrel{}{n}\text{and}a_\stackrel{}{n}\stackrel{~}{}^20`$ $``$ $`\gamma 1,1+\gamma \alpha ^20.`$ (23) These constitute convex sets. The extremal points of these convex sets will be the optimal measurements. To understand what extremal point corresponds to what optimal measurement we introduce a fidelity function $`f`$. We now study different fidelity functions. The covariance of the measurement set up implies that $`f`$ is a function only of the angle between the direction in which the spins are polarized $`\stackrel{}{m}`$ and the direction guessed by the POVM element $`a_\stackrel{}{n}`$. Thus in the case of outcome $`+z`$, $`f`$ is a function of $`m_z`$ only. It is convenient to expand $`f`$ in Legendre polynomials $`f(+z,\stackrel{}{m})`$ $`=`$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}f_nP_n(m_z)`$ (24) $`=`$ $`f_0+f_1m_z+f_2{\displaystyle \frac{3m_z^21}{2}}+\mathrm{}.`$ (25) To compute the average fidelity we need the probability of each outcome. Suppose that the spins point in direction $`\stackrel{}{m}`$ and that the measurement outcome is $`+z`$. This occurs with probability $`P(+z\|\stackrel{}{m})=Tr\rho (\stackrel{}{m},\stackrel{}{m})a_\stackrel{}{z}`$ $`=`$ $`1+\alpha m_z+{\displaystyle \frac{\gamma }{2}}{\displaystyle \frac{3m_z^21}{2}}.`$ (26) The average fidelity is therefore $`F`$ $`=`$ $`{\displaystyle _1^{+1}}{\displaystyle \frac{dm_z}{2}}{\displaystyle _0^{2\pi }}{\displaystyle \frac{d\varphi }{2\pi }}f(+z,\stackrel{}{m})P(+z\|\stackrel{}{m})`$ (27) $`=`$ $`f_0+{\displaystyle \frac{\alpha }{3}}f_1+{\displaystyle \frac{\gamma }{10}}f_2.`$ (28) Thus only the first three coefficients enter into the average fidelity. (In the case of covariant measurements on $`N`$ parallel spins only the $`N+1`$ first coefficients of the expansion of $`f`$ will enter into the average fidelity). Using eqs. (21, 22, 23) and (28) it is straightforward to find the optimal measurement for an arbitrary fidelity function in the case of parallel spins, antiparallel spins and LOCC measurements. As a first illustration, let us consider the example studied in and in which the fidelity has the form $`f(\stackrel{}{n}\|\stackrel{}{m})=\|_\stackrel{}{m}\|_\stackrel{}{n}\|^2=(1+\mathrm{cos}\theta )/2`$. Thus in this example $`f_0=1/2`$, $`f_1=1/2`$, $`f_2=0`$ and therefore $`F=1/2+\alpha /6`$. In this case the largest fidelity is obtained by taking for $`\alpha `$ the largest possible value. In the case of two parallel spins the largest possible value of $`\alpha `$ is $`\alpha _{max}=3/2`$ corresponding to $`F_{}=3/4=0.75`$, a result already obtained in and . In the case of two antiparallel spins $`\alpha _{max}=\sqrt{3}`$ corresponding to $`F_{}=1/2+1/(2\sqrt{3})0.788`$, a result already obtained in . In the case of measurements carried out using only LOCC, $`\alpha _{max}=\sqrt{2}`$ corresponding to $`F_{LOCC}=1/2+1/(3\sqrt{2})0.736`$. Thus for this fidelity collective measurements on antiparallel spins are better than collective measurements on parallel spins which are themselves better than LOCC measurements on parallel (or antiparallel) spins. Note that if the spins are parallel or antiparallel, optimal measurements that use a 1 dimensional ancilla (which could be the singlet state) have been exhibited in and . In the case of LOCC measurements it is easy to check that a simple optimal strategy consists in Alice making a von Neumann measurement of spin along some direction $`\stackrel{}{a}`$ and Bob making a von Neumann measurement of spin along an orthogonal direction $`\stackrel{}{b}`$ ($`\stackrel{}{a}\stackrel{}{b}=0`$). Denote by $`\stackrel{}{\alpha }=\pm \stackrel{}{a}`$ and $`\stackrel{}{\beta }=\pm \stackrel{}{b}`$ the results of the two measurements. Then the guessed direction is the bisectrix $`\stackrel{}{\alpha }+\stackrel{}{\beta }`$ of the two results. Thus in all cases the optimal measurements can be implemented using rather simple strategies which necessitate low dimensional ancillas. This should be compared with the covariant measurements which, although useful for the theoretical analysis, require infinite dimensional ancillas. As a second illustration consider the case where the fidelity is $`f=\mathrm{sin}^2\theta =1\mathrm{cos}^2\theta `$. In this case $`f_0=2/3`$, $`f_1=0`$, $`f_2=2/3`$, hence $`F=2/3\gamma /15`$ and the best measurement is that which has the smallest value of $`\gamma `$. In the case of collective measurements on parallel spins the smallest value is $`\gamma _{min}=2`$ yielding a fidelity $`F_{}=4/5=0.8`$. For collective measurements on antiparallel spins or LOCC measurements the minimum value is $`\gamma _{min}=1`$ yielding an optimal fidelity $`F_{,LOCC}=11/150.733`$. In this case measurements on two parallel spins are better than measurements on antiparallel spins or LOCC measurements which are both equivalent. Thus for some fidelities measurements on parallel spins are better, for other fidelities measurements on antiparallel spins are better, and in all cases LOCC measurements are the worst. In conclusion the present article gives explicitly the optimal measurements and optimal fidelity for all possible fidelity functions in the cases of parallel spins, antiparallel spins, and LOCC measurements. This provides an interesting target for future experiments since it provides a criterion for putting “non locality without entanglement” into evidence. I would like to thank N. Gisin, N. Linden and S. Popescu for stimulating and helpful discussions, and Dietmar Fisher for pointing out his recent work to me. I am a research associate of the Belgian National Fund for Scientific Research.
warning/0004/hep-ph0004027.html	ar5iv	text	# QCD analysis of near-to-planar 3-jet events ## 1 Introduction In recent years the standards were established for the perturbative QCD description of various characteristics of hadron jets produced in $`e^+e^{}`$ annihilation. These standards include: * all-order resummation of double- (DL) and single-logarithmic (SL) contributions due to soft and collinear gluon radiation effects, * two-loop analysis of the basic gluon radiation probability and * matching the resummed logarithmic expressions with the exact $`𝒪\left(\alpha _\text{s}^2\right)`$ results. Such programmes were carried out for a number of jet shape observables such as Thrust ($`T`$) and heavy jet mass ($`M_H`$, $`C`$-parameter and jet Broadenings (total $`B_T`$ and wide-jet broadening $`B_W`$. Perturbative description of jets produced in processes other than $`e^+e^{}`$ annihilation poses more difficulties, as the structure of jets (both average jet characteristics and distributions) becomes sensitive to details of the underlying hard interaction. For example, characteristics of quark jets produced in the current fragmentation region in Deep Inelastic Scattering (DIS) depend on the Bjorken $`x`$. In spite of these complications a steady progress is being made in this domain as well . Given a (relative) perfection of perturbative QCD technologies, it became possible to aim at deviations of the measured hadronic characteristics from the corresponding perturbative predictions, in a search for genuine non-perturbative confinement effects. As is well known, these deviations, for a broad variety of jet observables, amount to sizable $`1/Q`$ corrections, with $`Q`$ the hardness scale (the total $`e^+e^{}`$ annihilation energy), for reviews see . As far as physics of $`e^+e^{}`$ annihilation is concerned, till now these developments, both in the perturbative and non-perturbative sectors, were confined to two-jet ensembles which constitute the bulk of events. Little (if any) attention has been paid to multi-jet ensembles, in particular to three-jet events. In spite of being obviously more rare ($`\sigma _N/\sigma _{tot}\alpha _\text{s}^{N2}`$, with $`N`$ the number of jets), selected multi-jet configurations are of special interest as they, so to speak, are subject to more quantum mechanics than unrestricted (mainly two-jet) hadron production events. Indeed, in the latter case it is known that the gross inclusive features of particle production can be described in probabilistic terms by imposing angular ordering on successive $`12`$ intra-jet parton decays. It suffices to implement strict angular ordering, $`\theta _{i+1}\theta _i`$, the proper running coupling, $`\alpha _\text{s}(k_{})`$, and the standard parton splitting probabilities to ensure the next-to-leading accuracy of the perturbative description of inclusive energy spectra, mean multiplicities and the multiplicity distribution . Thus, a transparent and powerful probabilistic technology exists for predicting (with next-to-leading accuracy) inclusive observables. At the same time, the very selection of events (say, three-jet events) necessarily makes a measurement “less inclusive” and destroys, generally speaking, the probabilistic picture, and one’s hold on the SL accuracy with it. Of special interest for two-jet configurations is the kinematical region of two narrow jets ($`1T,M_H^2/Q^2,C,B_{T,W}1`$) where multiple radiation (Sudakov suppression) effects are essential. An analogue of this region for three-jet events is near-to-planar kinematics. Three hard partons — $`q\overline{q}`$ and a gluon $`g`$ — form a plane. Secondary gluon radiation off the three-prong QCD antenna brings in aplanarity. In what follows we choose the out-of-plane transverse momentum $`K_{\text{out}}`$ as an aplanarity measure. The distribution of events in $`K_{\text{out}}`$ is subject to DL suppression in the region of small aplanarity, $`K_{\text{out}}/Q1`$, where normal accompanying radiation is vetoed. The corresponding suppression factor, and thus the resulting $`K_{\text{out}}`$-distribution is easy to predict in a DL approximation which takes care of the leading contributions $`𝒪\left([\alpha _\text{s}\mathrm{ln}^2(Q/K_{\text{out}})]^n\right)`$ in all orders while disregarding subleading SL corrections of the order of $`[\alpha _\text{s}\mathrm{ln}(Q/K_{\text{out}})]^m`$. The answer is given simply by the product of three proper QCD Sudakov factors, $`F_q^2F_g`$, which veto radiation of gluons with out-of-plane momentum components exceeding $`K_{\text{out}}`$, off the three hard partons treated as independent emitters. The DL approximation is known to be too rough to be practically reliable. Improving it proves to be nontrivial a quest: at the level of SL terms the geometry of the underlying three-jet configuration enters the game: in addition to intra-jet parton multiplication one has to take into consideration inter-jet particle production. The latter, however, does not admit “classical” probabilistic interpretation: gluon radiation in-between jets results from the coherent action of all three elements of the antenna. This is a general feature which complicates the analyses of other, more simple, characteristics of three-jet ensembles as well. For example, contribution of coherent inter-jet particle flows enters, at the level of the next-to-leading $`𝒪\left(\sqrt{\alpha _\text{s}}\right)`$ correction, into perturbative prediction of the mean particle multiplicity in three-jet events, making it event-geometry-dependent . In this paper we attempt, for the first time, the all-order perturbative analysis of the $`K_{\text{out}}`$-distribution in three-jet $`e^+e^{}`$ annihilation events, aiming at SL accuracy. In what follows we will single out and resum logarithmically-enhanced DL and SL contributions and systematically neglect relative corrections of the order $`𝒪\left(\alpha _\text{s}\right)`$. The latter belong to the non-logarithmic normalization factor (“coefficient function”) $`1+c\alpha _\text{s}+\mathrm{}`$, whose first coefficient, $`c`$, can be found by comparing the approximate resummed result with numerical calculation based on the exact $`\alpha _\text{s}^2`$ matrix element . To accommodate all essential SL contributions one has 1. to take into account soft inter-jet gluon radiation, 2. to analyse corrections due to hard intra-jet parton decays, 3. to define an event plane for a multi-parton system and properly treat kinematical effects due to parton recoil, 4. to prove, for three-jet environment, soft gluon exponentiation (at the two-loop level) and the prescription for the argument of the running coupling that enters the basic gluon emission probability. In the present paper we address these issues. The answer we derive has the following key features: * Kinematical constraints which determine an event plane are rather complicated but can be resolved with a help of multiple Fourier–Mellin representation which allows for exponentiation of multiple radiation in the parameter space. * Multiple soft radiation off the three-parton system can be resummed and exponentiated in terms of three colour dipoles that together determine the colour structure of (and accompanying particle flows in) the event. * For the cumulative $`K_{\text{out}}`$-distribution (as well as for other sufficiently inclusive observables) inclusive treatment of the two-parton decay of a gluon emitted by the $`q\overline{q}g`$ system results in the running of the coupling constant. * The running $`\alpha _\text{s}`$ which describes the intensity of gluon emission off each dipole ($`qg`$, $`\overline{q}g`$, $`q\overline{q}`$) depends on the invariant transverse momentum of the gluon with respect to two partons that form the corresponding dipole. * The exponent can be represented as a sum of three basic parton “radiators” each of which describes one-gluon emission off a single hard parton, weighted by the colour factor of this parton. This radiation can be treated as independent provided a proper hardness scale is ascribed to the parton radiator. * Essential SL contribution (“hard” intra-jet and coherent inter-jet corrections) can be conveniently embodied into the scales. Effects of the large-angle soft (inter-jet) radiation make these scales event-geometry-dependent and different for each of the three primary partons. * The structure of the hardness scales of the parton radiators entering in the total $`K_{\text{out}}`$-distribution has a clear geometrical interpretation: the scale $`Q_a`$ for the parton $`a`$ is proportional to the invariant transverse momentum of this parton, $`p_{ta}`$, with respect to the hyper-plane formed by the other two hard partons. * Equivalent expressions for the radiators of the total $`K_{\text{out}}`$-distributions can be constructed which smoothly interpolate between 3- and 2-jet configurations. As has been already said, kinematics of three-jet observables complicates the analysis. As a result, the final expressions are rather cumbersome as they involve 4- and 5-dimensional integrals. The kinematical constraints, and thus the final formulae, are somewhat simpler for the distribution of $`K_{\text{out}}`$ accumulated in the right hemisphere, i.e. the one that has the smallest transverse momentum with respect to the thrust axis (for a review see ). Quantitative analysis of the predictions and numerical results are discussed in a separate publication . The paper is organised as follows. In Section 2 we define the aplanarity measure $`K_{\text{out}}`$, derive kinematical relations defining the event plane and discuss resummation of soft gluon radiation for near-to-planar 3-jet events. Section 3 is devoted to perturbative analysis of the $`K_{\text{out}}`$-distribution in the right hemisphere. Here we develop technique for analysing and embodying all necessary SL corrections into the all-order perturbative result. In Section 4 we apply this technique to derive, with SL accuracy, the perturbative prediction for the total $`K_{\text{out}}`$-distribution in 3-jet events with given kinematics (thrust $`T`$ and thrust-major $`T_M`$). Technical details are confined to Appendices. ## 2 Aplanarity and soft parton resummation We consider $`e^+e^{}`$ annihilation events with almost planar configuration of final state particles. Such event are characterised by the inequality $$TT_MT_m,$$ (1) with $`T`$ the thrust, $`T_M`$ and $`T_m`$ the so-called thrust-major and thrust-minor. We study the distribution of $`T_m`$. We shall refer to $`e^+e^{}`$ events in this phase space region as $`3`$-jet events. Thrust is defined as $$TQ=\underset{\stackrel{}{n}}{\mathrm{max}}\left\{\underset{h}{}\left\|\stackrel{}{n}\stackrel{}{p}_h\right\|\right\}=\underset{h}{}\left\|\stackrel{}{n}_T\stackrel{}{p}_h\right\|=\underset{h}{}\left\|p_{hz}\right\|,$$ (2) where the sum runs over all particles $`h`$ produced in a given event with the total center of mass energy $`Q`$. Hereafter we choose the $`z`$-axis to lie along the thrust axis, $`\stackrel{}{n}_T`$. The thrust-major is defined analogously; it maximises the sum of the particle momenta projections in the two-dimensional plane orthogonal to the thrust axis: $$T_MQ=\underset{\stackrel{}{n}\stackrel{}{n}_T=0}{\mathrm{max}}\left\{\underset{h}{}\left\|\stackrel{}{n}\stackrel{}{p}_h\right\|\right\}=\underset{h}{}\left\|\stackrel{}{n}_M\stackrel{}{p}_h\right\|=\underset{h}{}\left\|p_{hy}\right\|.$$ (3) The direction $`\stackrel{}{n}_M`$ we shall identify with the $`y`$-axis. Finally, for the thrust-minor we have $$T_mQ=\underset{h}{}\left\|p_{hx}\right\|K_{\text{out}}.$$ (4) We choose $`\stackrel{}{n}_T`$ in such a way that the most energetic particle in the event has a positive $`z`$-component and $`\stackrel{}{n}_M`$ in such a way that the second most energetic particle has a positive $`y`$-component. Hereafter we attribute $`p_h`$ to the right hemisphere $`C_R`$ (left hemisphere $`C_L`$) if $`p_{hz}>0`$ ($`p_{hz}<0`$). Similarly $`p_h`$ is in the up hemisphere $`C_U`$ (down hemisphere $`C_D`$) if $`p_{hy}>0`$ ($`p_{hy}<0`$). We shall also consider separately the right-hemisphere cumulative out-of-plane transverse momentum: $$K_{\text{out}}^R\underset{hR}{}\left\|p_{hx}\right\|.$$ (5) It can be easily shown that from the definition (23) the kinematical constraints follow: $$\begin{array}{c}\hfill \underset{hR}{}\stackrel{}{p}_{ht}=\underset{hL}{}\stackrel{}{p}_{ht}=0,\underset{hU}{}p_{hx}=\underset{hD}{}p_{hx}=0,\end{array}$$ (6) where $`\stackrel{}{p}_t`$ is the two-dimensional vector transversal to the $`z`$-axis. We introduce the three-dimensional vector $`\stackrel{}{P}_1`$ and the two-dimensional vector $`\stackrel{}{P}_{2t}`$ that define the “event plane” $`\{z,y\}`$: $$\stackrel{}{P}_1=(P_{1x},P_{1y},P_{1z})=(0,0,TE),\stackrel{}{P}_{2t}=(P_{2x},P_{2y})=(0,T_ME),Q=2E.$$ (7) From the definition of $`T`$ and $`T_M`$ we have $$\stackrel{}{P}_1=\underset{hR}{}\stackrel{}{p}_h,\stackrel{}{P}_{2t}=\underset{hU}{}\stackrel{}{p}_{ht}.$$ (8) In what follows we study the distribution of events in the cumulative out-of-plane transverse momentum $`K_{\text{out}}`$ defined in (4). The integrated $`K_{\text{out}}`$-distribution is defined as $$\frac{d\sigma (K_{\text{out}})}{dTdT_M}=Q^5\underset{m}{}𝑑\sigma _m\mathrm{\Theta }\left(K_{\text{out}}\underset{h=1}{\overset{m}{}}\|p_{hx}\|\right)\delta ^3\left(\underset{hR}{}\stackrel{}{p}_h\stackrel{}{P}_1\right)\delta ^2\left(\underset{hU}{}\stackrel{}{p}_{ht}\stackrel{}{P}_{2t}\right),$$ (9) where $`m`$ denotes the number of final particles in an event. The last two delta-functions fix the event plane and the theta-function defines the observable. Analogously we define the right distribution by restricting the sum over particles in the theta-function of (9) to those belonging to the right hemisphere (see (5)). ### 2.1 Jet momenta and hard parton recoil At the parton level the events in the region (1) can be treated as three-jet events generated by a system of energetic quark, antiquark and a gluon accompanied by an ensemble of soft partons. At the Born level, $`𝒪\left(\alpha _\text{s}\right)`$, in the absence of accompanying radiation the 3-parton system is truly planar, $`T_m0`$. The kinematical configuration of $`q,\overline{q}`$ and $`g`$ treated as massless partons is then uniquely fixed by the values of $`T`$ and $`T_M`$. Denoting by $`P_1,P_2`$ and $`P_3`$ the energy ordered parton momenta, $`P_{10}>P_{20}>P_{30}`$, we have $`\stackrel{}{P}_1`$ lying along the thrust axis. The vectors $`\stackrel{}{P}_1`$ and $`\stackrel{}{P}_{2t}=\stackrel{}{P}_{3t}`$ are given by the event plane momenta defined in (7). There are various kinematical configurations of the Born system. The three parton momenta $`P_q,P_{\overline{q}},P_g`$ can belong to three configurations $`𝒞_\delta `$, namely (see Figure 1) $$\begin{array}{cc}\hfill (P_1,P_2,P_3)& =(P_g,P_q,P_{\overline{q}})𝒞_1\hfill \\ & =(P_q,P_g,P_{\overline{q}})𝒞_2\hfill \\ & =(P_q,P_{\overline{q}},P_g)𝒞_3.\hfill \end{array}$$ (10) Notice that the index $`\delta `$ labelling the configuration $`𝒞_\delta `$ coincides with the index of the gluon momentum. (Interchanging the quark and the antiquark does not affect the accompanying radiation and the distributions under study.) For three massless partons the values of $`T`$ and $`T_M`$ are restricted to the region $$\frac{2(1T)}{T}\sqrt{2T1}<T_M<\sqrt{1T}.$$ (11) (For the kinematics of the Born system momenta see Appendix A.) In the following we will analyse the distribution inside this region — the 3-jet region — in which three skeleton parton momenta $`P_a`$ can be reconstructed from the $`T`$ and $`T_M`$ values. Beyond the Born approximation, in the presence of secondary gluon radiation, $`T_m`$ is no longer vanishing. The $`x`$-components of bremsstrahlung gluon momenta are logarithmically distributed over a broad range, so that $`T_m\alpha _\text{s}T_M\alpha _\text{s}T`$. In the region (1) the PT expansion develops logarithmically enhanced contributions which need to be resummed in all orders. Our aim is to perform this resummation with SL accuracy. This means that we shall keep both DL, $`\alpha _\text{s}\mathrm{log}^2T_m`$, and SL, $`\alpha _\text{s}\mathrm{log}T_m`$, contributions to the exponent (“radiator”, see below) while neglecting non-logarithmic corrections $`𝒪\left(\alpha _\text{s}\right)`$. With account of the running coupling effect, the DL and SL contributions formally expand into series of terms $`\alpha _\text{s}^n\mathrm{log}^{n+1}T_m`$ and $`\alpha _\text{s}^n\mathrm{log}^nT_m`$, respectively. It is important to stress that exponentiation of the SL correction $`\alpha _\text{s}\mathrm{log}T_m`$ makes sense only if it is supported by the calculation of the non-logarithmic “coefficient function”. Indeed, a “cross-talk” between an $`𝒪\left(\alpha _\text{s}\right)`$ correction to the overall normalization of the resummed distribution and the leading DL term, symbolically, $$(1+c\alpha _\text{s})\times \mathrm{exp}\left\{\alpha _\text{s}\mathrm{log}^2+s\alpha _\text{s}\mathrm{log}\right\}=1+\mathrm{}+c\alpha _\text{s}\times \alpha _\text{s}\mathrm{log}^2+\frac{s^2}{2!}(\alpha _\text{s}\mathrm{log})^2+\mathrm{},$$ (12) gives rise to the correction of the same order as the SL term squared. The coefficient $`c`$ is not known analytically. It can be determined by comparing the $`𝒪\left(\alpha _\text{s}^2\right)`$ term of the log-resummed expression with the result of a numerical calculation based on the exact $`\alpha _\text{s}^2`$ matrix element . We denote by $`p_a`$ the momenta of the three hard partons, $`q`$, $`\overline{q}`$, $`g`$, that in general no longer lie in the event plane defined by the vectors (7). In the region (1) they differ from $`P_a`$ by “soft recoil parts” $`q_a`$, $$p_a=P_a+q_a,$$ (13) whose $`x`$-components contribute to $`K_{\text{out}}`$ together with the soft parton momenta $`k_i`$: $$K_{\text{out}}=\|q_{1x}\|+\|q_{2x}\|+\|q_{3x}\|+\|k_{ix}\|.$$ (14) To SL accuracy, as in the case of the broadening distribution , it suffices to keep the recoil momenta $`q_a`$ only in the phase space, in particular their contribution to $`K_{\text{out}}`$ and to the constraints defining the event plane (see (9)). At the same time, $`q_a`$ can be neglected in the radiation matrix element, that is, in the emission distributions we can substitute the skeleton momenta $`P_a`$ for actual parton momenta $`p_a`$. We remark that, as in the case of jet broadening , for the study of non-perturbative power-suppressed corrections the approximation $`p_aP_a`$ in the radiation matrix element is not valid . In order to resum the PT series for (9) in the region (1) we need to use the factorization property of the soft radiation matrix element and to factorize the multi-parton phase space. We discuss these points in succession after giving a brief description of the kinematics. For more details see Appendix A. ### 2.2 Kinematics of soft or collinear emission Introducing the recoil momenta $`q_a`$ defined in (13), the right, left and total $`K_{\text{out}}`$ are given by $$K_{\text{out}}^R=\|q_{1x}\|+\underset{R}{}\|k_{ix}\|,K_{\text{out}}^L=\|q_{2x}\|+\|q_{3x}\|+\underset{L}{}\|k_{ix}\|,K_{\text{out}}=K_{\text{out}}^R+K_{\text{out}}^L.$$ (15) The event plane momenta (7) are given by $$\begin{array}{c}\hfill \stackrel{}{P}_1=\stackrel{}{p}_1+\underset{R}{}\stackrel{}{k}_i,\stackrel{}{P}_{2t}=\stackrel{}{p}_{2t}+\stackrel{}{p}_{1t}\vartheta (p_{1y})+\underset{U}{}\stackrel{}{k}_{it}.\end{array}$$ (16) In the region in which the emitted partons are soft or collinear the components of $`q_a`$ that accompany non-vanishing components of $`P_a`$ can be neglected in the calculation. Therefore we consider only those four components that vanish in the Born approximation namely, the $`y`$-component of the leading parton momentum, $`q_{1y}=p_{1y}`$, and the $`x`$-components $`q_{ax}=p_{ax}`$, $`a=1,2,3`$. Then, from the kinematical relations (16) and (5) we have $$\begin{array}{cc}& q_{1y}+\underset{R}{}k_{iy}=q_{1x}+\underset{R}{}k_{1x}=0,\hfill \\ & q_{2x}+q_{1x}^++\underset{U}{}k_{1x}=q_{3x}+q_{1x}^{}+\underset{D}{}k_{1x}=0,q_{1x}^\pm =q_{1x}\vartheta (\pm q_{1y}).\hfill \end{array}$$ (17) Notice that this implies the following kinematical relations $$\begin{array}{c}\hfill \left\{q_{1x}^++\underset{iC_1}{}k_{ix}\right\}=\left\{q_{1x}^{}+\underset{iC_4}{}k_{ix}\right\}=\left\{q_{2x}+\underset{iC_2}{}k_{ix}\right\}=\left\{q_{3x}+\underset{iC_3}{}k_{ix}\right\}.\end{array}$$ (18) ### 2.3 Matrix element factorization At the parton level the distribution is given by the sum of the partial cross sections $$d\sigma _n=\frac{1}{n!}M_n^2d\mathrm{\Phi }_n,$$ (19) where $`M_n`$ is the matrix element for the emission of $`n`$ secondary partons off the $`q\overline{q}`$, $`g`$ system, and $`d\mathrm{\Phi }_n`$ the corresponding phase-space factor, $$\begin{array}{c}\hfill d\mathrm{\Phi }_n=(2\pi )^4\delta ^4(\underset{a}{}p_a+\underset{i}{}k_iQ)\underset{i=1}{\overset{n}{}}[dk_i]\underset{a=1}{\overset{3}{}}\frac{d^3p_a}{(2\pi )^32E_a},[dk]=\frac{d^3k}{\pi \omega }.\end{array}$$ (20) As it is explained in detail in Appendix B.2.2 the contribution of collinear non soft emission is process independent and can be taken into account a posteriori introducing in the computed distribution the hard part of the splitting functions. Therefore the accompanying partons $`k_i`$ are assumed to be soft and their distribution can be treated as independent. Assembling together the contributions from the three configurations $`𝒞_\delta `$ defined in (10), the distribution can be presented as $$M_n^2=\underset{\delta =1}{\overset{3}{}}M_0^2(𝒞_\delta )\underset{i}{\overset{n}{}}W_\delta (k_i),$$ (21) where $`M_0(𝒞_\delta )`$ is the Born $`q\overline{q}g`$ matrix element and $`W_\delta `$ is the distribution of the soft gluon radiation off the hard three-parton antenna in the momentum configuration $`𝒞_\delta `$. Here we discuss the real emission contribution to $`M_n`$; the virtual corrections will be accommodated later. For the configuration $`\delta =3`$, for example, the squared Born matrix element reads $$M_0^2(𝒞_3)=\frac{C_F\alpha _\text{s}}{2\pi }\frac{x_1^2+x_2^2}{(1x_1)(1x_2)},x_a\frac{2P_aQ}{Q^2},$$ (22) where $`P_3`$ is the gluon momentum and $`P_1,P_2`$ the quark-antiquark momenta (see (10)). For this configuration the single soft gluon radiation pattern is given, at the one-loop level, by $$W_3(k)=C_Fw_{12}+\frac{N_c}{2}\left(w_{13}+w_{23}w_{12}\right)=\frac{N_c}{2}\left(w_{13}+w_{23}\frac{1}{N_c^2}w_{12}\right).$$ (23) Here $`w_{ab}`$ is the standard two-parton antenna of the $`ab`$-dipole, which, within the normalization convention prescribed by (20), is given by $$\begin{array}{c}\hfill w_{ab}(k)=\frac{\alpha _\text{s}}{\pi }\frac{(P_aP_b)}{2(P_ak)(kP_b)}=\frac{\alpha _\text{s}}{\pi k_{t,ab}^2}.\end{array}$$ (24) Here $`k_{t,ab}`$ is the invariant gluon transverse momentum with respect to the hyper-plane defined by the $`P_a,P_b`$ momenta. The first term $`C_Fw_{12}`$ in (23) is the “Abelian” contribution describing soft gluon emission off the $`q\overline{q}`$ pair. The second term proportional to $`N_c`$ is its “non-Abelian” counterpart that describes radiation off the hard gluon $`P_3`$. Similar expressions for two other kinematical configurations ($`\delta =1,2`$) is straightforward to write down by properly adjusting the parton indices, $$\begin{array}{cc}& W_1(k)=\frac{N_c}{2}\left(w_{13}+w_{12}\frac{1}{N_c^2}w_{23}\right),\hfill \\ & W_2(k)=\frac{N_c}{2}\left(w_{23}+w_{12}\frac{1}{N_c^2}w_{13}\right).\hfill \end{array}$$ (25) To reach SL accuracy it is necessary to treat multi-parton emission at the two-loop level. This involves allowing secondary gluon to split into two gluons or into a $`q\overline{q}`$ pair. In principle, perturbative analysis of a system consisting of three hard partons and two softer partons (with comparable energies) can be found in the literature (see, e.g., ). However, to the best of our knowledge, the most important feature of the result has never been explicitly stressed, namely, that the colour structure and geometrical properties of emission of a soft two-parton system off a 3-jet ensemble is identical to those for single gluon radiation. It can be shown that after subtracting the uncorrelated radiation of two soft gluons, $`W_3(k_1)W_3(k_2)`$, the correlated two-parton production is given by the expression $$W_3^{(2)}(k_1,k_2)=C_Fw_{12}^{(2)}(k_1,k_2)+\frac{N_c}{2}\left(w_{13}^{(2)}(k_1,k_2)+w_{23}^{(2)}(k_1,k_2)w_{12}^{(2)}(k_1,k_2)\right),$$ (26) where $`w_{ab}^{(2)}`$ is the standard distribution known from the two-loop analysis of 2-jet event shapes (the first term on the right-hand side of (26), see ). It describes decay into $`q\overline{q}`$\- or $`gg`$-pair of a virtual parent gluon radiated, in our case, by one of the three two-parton dipoles $`ab`$. Given this remarkably simple dipole structure, the analysis of the two-loop effects in 3-jet events reduces to the known 2-jet case. As shown in Ref. , for a sufficiently inclusive observable (and the $`K_{\text{out}}`$-distribution under interest in particular) the two-loop refinement results in (and reduces to) substituting the proper argument for the running coupling describing the intensity of the parent gluon emission, $$w_{ab}(k)=\frac{\alpha _\text{s}}{\pi k_{t,ab}^2}w_{ab}(k)+𝑑k_1𝑑k_2\delta (kk_1k_2)w_{ab}^{(2)}(k_1,k_2)\frac{\alpha _\text{s}(k_{t,ab})}{\pi k_{t,ab}^2}.$$ (27) Here $`k_{t,ab}`$ is the invariant transverse momentum defined in (24) and $`\alpha _\text{s}`$ is taken in the physical scheme . It is worthwhile to notice that the arguments of the coupling are different for the three dipoles that participate in gluon radiation according to (23). We remark that at the level of the leading power-suppressed non-perturbative correction the rôle of two-loop effects is more involved as they give rise to the Milan factor . ### 2.4 Phase space factorisation We discuss here factorization of the phase space in the soft region, which is needed for resummation. The phase space factor reads $$\begin{array}{cc}\hfill d\mathrm{\Gamma }_n& =d\mathrm{\Phi }_n\delta ^3\left(\stackrel{}{p}_1+\underset{R}{}\stackrel{}{k}_i\stackrel{}{P}_1\right)\delta ^2\left(\stackrel{}{p}_{2t}+\stackrel{}{p}_{1t}\vartheta (p_{1y})+\underset{U}{}\stackrel{}{k}_{it}\stackrel{}{P}_{2t}\right)\hfill \\ & =\underset{i=1}{\overset{n}{}}[dk_i]\underset{a=1}{\overset{3}{}}\frac{d^3p_a}{(2\pi )^32E_a}D_n,[dk]=\frac{d^3k}{\pi \omega },\hfill \end{array}$$ (28) where the factor $`D_n`$ takes care of the kinematical relations, $$D_n=(2\pi )^4\delta ^4(\underset{a}{}p_a+\underset{i}{}k_iQ)\delta ^3(\stackrel{}{p}_1+\underset{R}{}\stackrel{}{k}_i\stackrel{}{P}_1)\delta ^2\left(\stackrel{}{p}_{2t}+\stackrel{}{p}_{1t}\vartheta (p_{1y})+\underset{U}{}\stackrel{}{k}_{it}\stackrel{}{P}_{2t}\right).$$ (29) The first delta-function stands for the energy–momentum conservation, while the last two (three- and two-dimensional) delta-functions set the event plane. Now we single out from $`d\mathrm{\Gamma }_n`$ small recoil components of hard partons momenta, $`q_{1y}=p_{1y}`$ and the three $`q_{ax}=p_{ax}`$, which have to satisfy four event-plane constraints given in (17). Then, introducing the unity, $$1=𝑑q_{1y}\underset{a}{}dq_{ax}S_n,$$ where $$S_n\delta \left(q_{1y}+\underset{R}{}k_{iy}\right)\delta \left(q_{1x}+\underset{R}{}k_{ix}\right)\delta \left(q_{2x}+q_{1x}^++\underset{U}{}k_{ix}\right)\delta \left(q_{3x}+q_{1x}^{}+\underset{D}{}k_{ix}\right),$$ (30) we neglect hard parton recoils $`q_{ai}`$ in all but these four components and approximate the kinematical factor $`D_n`$ as follows $$\begin{array}{c}\hfill D_n=𝑑q_{1y}\underset{a}{}dq_{ax}S_nD_nD_0𝑑q_{1y}\underset{a}{}dq_{ax}S_n.\end{array}$$ (31) Here $`D_0`$ is a trivial phase space factor which corresponds to the Born three-parton kinematics, $$D_0=(2\pi )^4\delta ^4\left(\underset{a}{}p_aQ\right)\delta ^3\left(\stackrel{}{p}_1\stackrel{}{P}_1\right)\delta ^2\left(\stackrel{}{p}_{2t}\stackrel{}{P}_{2t}\right).$$ We then have $$d\mathrm{\Gamma }_n\mathrm{\Gamma }_0\underset{i=1}{\overset{n}{}}[dk_i]dh_n,dh_n=dq_{1y}\underset{a=1}{\overset{3}{}}dq_{ax}S_n,\mathrm{\Gamma }_0=\underset{a=1}{\overset{3}{}}\frac{d^3p_a}{(2\pi )^32E_a}D_0,$$ (32) where $`\mathrm{\Gamma }_0`$ is the Born phase space given in Appendix A. Neglecting bremsstrahlung in the $`D_0`$-factor proves to be legitimate: it can be shown that to achieve SL accuracy it suffices to take care of accompanying parton momenta in the $`S`$-factor (30) and in the observable itself. Finally, in order to “exponentiate” the multiple radiation we need to factorize dependence on individual secondary parton momenta contained in the delta-functions in $`S_n`$ and in the $`K_{\text{out}}`$-observable. This is achieved in a standard way by means of Mellin and Fourier representations. In the following we apply this procedure to the integrated $`K_{\text{out}}`$-distribution defined in (9). In what follows we shall separately consider $`K_{\text{out}}`$ accumulated in the right hemisphere and the total $`K_{\text{out}}`$ of the event. We start from a simpler case of the $`K_{\text{out}}`$-distribution in the right (one-jet) hemisphere. The total $`K_{\text{out}}`$-distributions will be considered in Sec. 4. ## 3 Right $`K_{\text{out}}`$-distribution Consider the differential three-jet cross section with given $`T`$ and $`T_M`$ and with accumulated out-of-event-plane momentum in the right hemisphere smaller than a given $`K_{\text{out}}`$. Using the soft-factorization formula (21) we can write the cross section (9) for small $`K_{\text{out}}`$ in the form $$\begin{array}{cc}\hfill \frac{d\sigma ^R(K_{\text{out}})}{dTdT_M}& =\underset{n}{}\frac{1}{n!}M_n^2𝑑\mathrm{\Gamma }_n\vartheta (K_{\text{out}}\left\|q_{1x}\right\|\underset{R}{}\|k_{ix}\|)\hfill \\ & =\underset{\delta =1}{\overset{3}{}}\frac{d\sigma _\delta ^{(0)}}{dTdT_M}\mathrm{\Sigma }_\delta ^R(K_{\text{out}}),\hfill \end{array}$$ (33) where $`\sigma _\delta ^{(0)}`$ is the three-jet differential Born cross section for the parton configuration $`𝒞_\delta `$ defined in (10), $$\frac{d\sigma _\delta ^{(0)}}{dTdT_M}\mathrm{\Gamma }_0M_0^2(𝒞_\delta ),$$ calculated in Appendix A. The accompanying radiation factor $`\mathrm{\Sigma }`$, a function of $`K_{\text{out}}`$, $`T`$ and $`T_M`$, reads $$\begin{array}{cc}& \mathrm{\Sigma }_\delta ^R(K_{\text{out}})=\underset{n}{}\frac{1}{n!}\underset{i}{\overset{n}{}}[dk_i]W_\delta (k_i)H^R(K_{\text{out}}),\hfill \\ & H^R(K_{\text{out}})𝑑h_n\vartheta (K_{\text{out}}\left\|q_{1x}\right\|\underset{R}{}\|k_{ix}\|).\hfill \end{array}$$ (34) We recall that within the adopted PT accuracy the hard parton momenta $`p_a`$ entering into the soft gluon distribution factor $`W_\delta (k)`$ can be approximated by the event plane vectors $`P_a`$ so that integration over the recoil variables $`q_a=p_aP_a`$ can be easily performed. Since the observable involves only momenta in the right hemisphere, the recoil momentum components in the left hemisphere $`q_{2x}`$ and $`q_{3x}`$ can be freely integrated out with use of the last two delta-functions in (30) and leave no trace in the distribution under consideration. Then, a non-trivial dependence on $`q_{1y}`$ (via $`q_{1x}^\pm )`$) also disappears, and the $`q_{1y}`$–integration trivializes as well. We are left with a single delta-function in (30): $$H^R(K_{\text{out}})=𝑑q_{1x}\vartheta (K_{\text{out}}\left\|q_{1x}\right\|\underset{R}{}\|k_{ix}\|)\delta \left(q_{1x}+\underset{R}{}k_{ix}\right).$$ To factorize the dependence on secondary parton momenta we use the Mellin representation for the theta-function in $`K_{\text{out}}`$ and the Fourier representation for the delta-function in $`q_{1x}`$: $`\vartheta (K_{\text{out}}\left\|q_{1x}\right\|{\displaystyle \underset{R}{}}\|k_{ix}\|)`$ $`=`$ $`{\displaystyle _𝒞}{\displaystyle \frac{d\nu }{2\pi i\nu }}\mathrm{exp}\left\{\nu (K_{\text{out}}\left\|q_{1x}\right\|{\displaystyle \underset{R}{}}\|k_{ix}\|)\right\},`$ (35) $`\delta (q_{1x}+{\displaystyle \underset{R}{}}k_{ix})`$ $`=`$ $`\nu {\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{d\beta }{2\pi }}\mathrm{exp}\left\{i\nu \beta (q_{1x}+{\displaystyle \underset{R}{}}k_{ix})\right\},`$ (36) where the contour $`𝒞`$ in (35) runs parallel to the imaginary axis at $`\text{Real}\nu >0`$. Defining the Fourier integral (36) we have extracted, for the sake of convenience, the factor $`\nu `$ as if it were a real parameter: for a complex value of $`\nu `$ it implies rotating the $`\beta `$-contour by $`Arg\nu `$, $`\left\|Arg\nu \right\|<\pi /2`$. (Then, using the analytic continuation, the $`\beta `$-integral can be transformed to run along the real axis.) Integrating over $`q_{1x}`$ we obtain $$H^R(K_{\text{out}})=_𝒞\frac{d\nu }{2\pi i\nu }e^{\nu K_{\text{out}}}_{\mathrm{}}^{\mathrm{}}\frac{d\beta }{\pi (1+\beta ^2)}\underset{R}{}e^{\nu (\|k_{ix}\|+i\beta k_{ix})}.$$ (37) The limit of small $`K_{\text{out}}`$ (and thus of small $`q_{1x}`$) corresponds to the region of large values of the conjugate variables, i.e. $`\nu `$ and $`\nu \beta `$ respectively. Therefore in what follows we will concentrate on the limit of large $`\nu `$ and neglect the contributions of the order of $`\nu ^1`$ which correspond to $`𝒪\left(K_{\text{out}}\right)`$ corrections to the distribution. At the same time, by examining (37) it is easy to see that the characteristic values of the rescaled Fourier variable $`\beta `$ are of the order of unity. Substituting (37) into (34) we get $$\mathrm{\Sigma }_\delta ^R(K_{\text{out}})=_𝒞\frac{d\nu }{2\pi i\nu }e^{\nu K_{\text{out}}}_{\mathrm{}}^{\mathrm{}}\frac{d\beta }{\pi (1+\beta ^2)}e^{_\delta ^R(\nu ,\beta )},$$ (38) where we have introduced the “radiator” $$_\delta ^R(\nu ,\beta )=_R[dk]W_\delta (k)\left[\mathrm{\hspace{0.25em}1}e^{\nu (\|k_x\|+i\beta k_x)}\right].$$ (39) Here the unity in the square brackets has been included to account for the virtual corrections. For example, for the most probable jet configuration $`\delta =3`$ (with gluon the least energetic parton), the soft distribution (23) gives $$_3^R=\frac{N_c}{2}\left(r_{13}^R+r_{23}^R\frac{1}{N_c^2}r_{12}^R\right),r_{ab}^R=_R[dk]w_{ab}(k)\left[\mathrm{\hspace{0.25em}1}e^{\nu (\|k_x\|+i\beta k_x)}\right].$$ (40) With the DL accuracy only gluons collinear to $`P_1`$ (the thrust axis direction) contribute to the right-hemisphere $`K_{\text{out}}`$. Therefore the (identical) DL contributions are contained in $`r_{12}^R`$ and $`r_{13}^R`$. According to (40), they combine into the expression proportional to $`C_F`$ — the colour charge of the quark $`P_1`$. Single logarithmic corrections to these two dipoles, as well as SL contribution of the third dipole, $`r_{23}^R`$, are calculated in Appendix B.4. Here we report the result: $$\begin{array}{c}\hfill r_{1a}^R=r(\overline{\mu },Q^2)+F_{1a}_{1/\overline{\mu }}^Q\frac{dk_x}{k_x}\frac{\alpha _\text{s}(2k_x)}{\pi },r_{23}^R=F_{23}_{1/\overline{\mu }}^Q\frac{dk_x}{k_x}\frac{\alpha _\text{s}(2k_x)}{\pi },\end{array}$$ (41) where the variable $`\overline{\mu }`$ originates from an approximate evaluation of the characteristic momentum integral which is explained in Appendix B.6, $$\left[\mathrm{\hspace{0.25em}1}e^{\nu (\|k_x\|+i\beta k_x)}\right]\theta \left(\|k_x\|\frac{1}{\overline{\mu }}\right),\overline{\mu }=\overline{\nu }\mu ,\mu =\sqrt{1+\beta ^2},$$ (42) with $$\overline{\nu }\nu e^{\gamma _E}.$$ (43) In (41) $`r(\overline{\mu },Q^2)`$ is the DL function $$r(\overline{\mu },Q^{}_{}{}^{}2)=_{1/\overline{\mu }}^Q^{}\frac{dk_x}{k_x}\frac{\alpha _\text{s}(2k_x)}{\pi }\mathrm{ln}\frac{Q^{}_{}{}^{}2}{k_x^2},$$ (44) and the factors $`F_{ab}=F_{ab}(T,T_M)`$ are independent of the integration variables $`\nu `$ and $`\beta `$. The origin of an essential subleading correction embodied into the precise argument of the running coupling, $`\alpha _\text{s}(2k_x)`$, is explained in Appendix B.5. Combining these results we obtain the radiators, evaluated with SL accuracy, for each of the three kinematical jet configurations, $$_\delta ^R(\overline{\mu })=C_\delta r(\overline{\mu },Q_\delta ^2).$$ (45) Here $`C_\delta `$ is the colour charge of the hard parton along the thrust axis, that is $`C_\delta =C_F`$ for $`\delta =2,3`$ and $`C_1=N_c`$. SL corrections in (45) were absorbed into the definition of the hard scale. The corresponding scale $`Q_\delta `$ depends on the event configuration $`\delta `$ and, though the functions $`F_{ab}(T,T_M)`$, on the event kinematics. These scales are given by $$\begin{array}{cc}& N_C\mathrm{ln}Q_1^2=N_C\mathrm{ln}(Q^2e^{\frac{\beta _0}{2N_c}})+\frac{N_c}{2}(F_{12}+F_{13}\frac{1}{N_c^2}F_{23}),\hfill \\ & C_F\mathrm{ln}Q_2^2=C_F\mathrm{ln}(Q^2e^{\frac{3}{2}})+\frac{N_c}{2}(F_{12}+F_{23}\frac{1}{N_c^2}F_{13}),\hfill \\ & C_F\mathrm{ln}Q_3^2=C_F\mathrm{ln}(Q^2e^{\frac{3}{2}})+\frac{N_c}{2}(F_{13}+F_{23}\frac{1}{N_c^2}F_{12}).\hfill \end{array}$$ (46) They also include the factors $`e^{3/4}`$ and $`e^{\beta _0/4N_c}`$ coming from SL corrections due to the “hard” parts of the quark and gluon splitting functions. As we shall see below in Sec. 4, for the case of the total $`K_{\text{out}}`$-distribution the corresponding scales are simply related with geometry of the 3-jet event. At the same time, the $`T`$\- and $`T_M`$-dependence of the scales entering into one-hemisphere distribution does not have a simple geometrical interpretation. This is due to the fact that the kinematical constraint restricting the observable to a single hemisphere is foreign to the structure of the soft-gluon radiation pattern. The radiator (45) depends on the Mellin-Fourier moments $`\nu `$ and $`\beta `$ only via the variable $`\overline{\mu }`$. The $`\beta `$-dependence can be further simplified at SL accuracy by expanding the radiator for large $`\nu `$, $$_\delta ^R(\overline{\mu })=_\delta ^R(\overline{\nu })+C_\delta r^{}(\overline{\nu })\mathrm{ln}\sqrt{1+\beta ^2},$$ where $$r^{}(\overline{\nu })=\frac{\alpha _\text{s}(k_x)}{\pi }\mathrm{ln}\frac{Q^2}{k_x^2},k_x1/\overline{\nu }.$$ (47) Since $`r^{}`$ constitutes a SL correction, we can use $`Q`$ as a common scale (neglecting $`𝒪\left(\alpha _\text{s}\right)`$ mismatch) and omit the factor $`2`$ in the running coupling argument (as producing a negligible correction $`𝒪\left(\alpha _\text{s}^2\mathrm{log}\overline{\nu }\right)`$) in such subleading terms. Thus, the right-hemisphere $`K_{\text{out}}`$-distribution to SL accuracy takes the form $$\begin{array}{cc}& \mathrm{\Sigma }_\delta ^R(K_{\text{out}})=\frac{d\nu }{2\pi i\nu }e^{\nu K_{\text{out}}}e^{_\delta ^R(\overline{\nu })}_\delta ^R(\overline{\nu }),\hfill \\ & _\delta ^R(\overline{\nu })=_{\mathrm{}}^{\mathrm{}}\frac{d\beta }{\pi (1+\beta ^2)^{1+\eta }}=\frac{\mathrm{\Gamma }(\frac{1}{2}+\eta )}{\sqrt{\pi }\mathrm{\Gamma }(1+\eta )},\eta =\frac{1}{2}C_\delta r^{}(\overline{\nu }).\hfill \end{array}$$ (48) Integration over the Mellin variable $`\nu `$ can be performed by steepest descent or by the operator method introduced in . We follow the last method which exploits the following identities: $$\begin{array}{cc}\hfill f(\overline{\nu })=& f(e^_z)\left(\overline{\nu }\right)^z\|_{z=0};\hfill \\ \hfill \frac{d\nu }{2\pi i\nu }\left(\overline{\nu }\right)^ze^{\nu K_{\text{out}}}=& \frac{(\overline{K}_{\text{out}})^z}{\mathrm{\Gamma }(1+z)},\overline{K}_{\text{out}}e^{\gamma _E}K_{\text{out}}.\hfill \end{array}$$ The following approximation is valid: $$e^{R(e^_z)}(e^_z)\frac{(\overline{K}_{\text{out}})^z}{\mathrm{\Gamma }(1+z)}\|_{z=0}e^{R(\overline{K}_{\text{out}}^1)}(\overline{K}_{\text{out}}^1)\mathrm{\Gamma }^1\left(1+R^{}(\overline{K}_{\text{out}}^1)\right),$$ where we have neglected relative corrections of the order $`R^{\prime \prime }(x)x_xR^{}(x)=𝒪\left(\alpha _\text{s}\right)`$ and assumed that $``$ is a smooth function of $`R^{}`$ which is true for $`_\delta ^R(\overline{\nu })`$ in (48), $`d\mathrm{ln}/d\eta =𝒪\left(1\right)`$. Applying these relations to our distribution in (48) we derive the final answer to SL accuracy, $$\mathrm{\Sigma }_\delta ^R(K_{\text{out}})=e^{_\delta ^R\left(\overline{K}_{\text{out}}^1\right)}\frac{_\delta ^R\left(\overline{K}_{\text{out}}^{}{}_{}{}^{1}\right)}{\mathrm{\Gamma }\left(1+C_\delta r^{}(\overline{K}_{\text{out}}^1)\right)},\overline{K}_{\text{out}}=e^{\gamma _E}K_{\text{out}}.$$ (49) The first and second factors resum DL and SL contributions, respectively. We remark that the precise argument $`\overline{K}_{\text{out}}`$ is essential to keep in the first factor, while substituting $`\overline{K}_{\text{out}}`$ by, say, $`K_{\text{out}}`$ in the second factor would amount to a negligible $`𝒪\left(\alpha _\text{s}\right)`$ correction. Comparing the $`\nu `$-integrand in (48) with the final result we conclude that the SL factor $`\mathrm{\Gamma }^1`$ in (49) accounts for a mismatch between the Mellin-conjugated $`\nu `$ and $`K_{\text{out}}`$ values. It can be looked upon as a next-to-leading order prefactor of the WKB (steepest descent) approximation. ## 4 Total $`K_{\text{out}}`$-distribution The calculation is similar to the previous case except that now radiation in all four quadrants contributes to $`K_{\text{out}}`$. The $`K_{\text{out}}`$-integrated distribution (9) for small $`K_{\text{out}}`$ is given by $$\begin{array}{cc}& \frac{d\sigma (K_{\text{out}})}{dTdT_M}=\underset{𝒞_\delta }{}\frac{d\sigma _\delta ^{(0)}}{dTdT_M}\mathrm{\Sigma }_\delta ^T(K_{\text{out}}),\hfill \\ & \mathrm{\Sigma }_\delta ^T(K_{\text{out}})=\underset{n}{}\frac{1}{n!}\underset{i}{\overset{n}{}}[dk_i]W_\delta (k_i)H^T(K_{\text{out}}),\hfill \\ & H^T(K_{\text{out}})𝑑h_n\vartheta (K_{\text{out}}\underset{a=1}{\overset{3}{}}\left\|q_{ax}\right\|\underset{i}{}\|k_{ix}\|).\hfill \end{array}$$ (50) To factorize the soft momenta in $`H^T`$ we proceed as before by using Mellin and Fourier representation for the theta- and delta-functions. Again we denote by $`\nu `$ the variable conjugate to $`K_{\text{out}}`$ and study the region $`\left\|\nu \right\|Q1`$. We rescale by $`\nu `$ the variables conjugate to the soft recoil variables $`q_{ax}`$ ($`a=1,2,3`$) and $`q_{1y}`$ in $`dh_n`$ to arrive at $$\begin{array}{cc}& 𝑑h_n\underset{a=1}{\overset{3}{}}e^{\nu \left\|q_{ax}\right\|}\underset{i=1}{\overset{n}{}}e^{\nu \left\|k_{ix}\right\|}=\hfill \\ & _{\mathrm{}}^{\mathrm{}}\frac{d\gamma }{2\pi }\left(\underset{a=1}{\overset{3}{}}\frac{d\beta _a}{\pi }\right)I(\beta ,\gamma )\left\{\underset{C_1}{}u_{12}(k_i)\underset{C_4}{}u_{13}(k_i)\underset{C_2}{}u_2(k_i)\underset{C_3}{}u_3(k_i)\right\}.\hfill \end{array}$$ (51) Here we have introduced the probing functions $`u_\alpha (k)`$ for each of the quadrants $`C_{\mathrm{}}`$: $$\begin{array}{cc}& u_{12}(k)=u(\beta _{12},\gamma ),u_{13}(k)=u(\beta _{13},\gamma ),u_2(k)=u(\beta _2,0),u_3(k)=u(\beta _3,0),\hfill \end{array}$$ (52) where $`\beta _{12}=\beta _1+\beta _2`$, $`\beta _{13}=\beta _1+\beta _3`$ and the “source function” $`u`$ is $$\begin{array}{c}\hfill u(\beta ,\gamma )\mathrm{exp}\left\{\nu \left(\left\|k_{ix}\right\|+i\beta k_{ix}+i\gamma \|k_{iy}\|\right)\right\}.\end{array}$$ (53) The function $`I(\beta ,\gamma )`$ in (51) is the result of integrations over the recoil momenta, namely, $`q_{1y}`$ (conjugate to $`\gamma `$) and three $`q_{ax}`$ (conjugate to $`\beta _a`$): $$\begin{array}{c}\hfill I(\beta ,\gamma )=\frac{1}{1+\beta _2^2}\frac{1}{1+\beta _3^2}\left(\frac{1}{1+\beta _{12}^2}\frac{1}{i\gamma +ϵ}+\frac{1}{1+\beta _{13}^2}\frac{1}{i\gamma +ϵ}\right).\end{array}$$ (54) Notice that the large variable $`\nu `$ enters only in the sources. Now we are in a position to resum multiple accompanying radiation. The result reads $$\begin{array}{c}\hfill \mathrm{\Sigma }_\delta (K_{\text{out}})=\frac{d\nu }{2\pi i\nu }e^{\nu K_{\text{out}}}\frac{d\gamma }{2\pi }\underset{a=1}{\overset{3}{}}\frac{d\beta _a}{\pi }I(\beta ,\gamma )e^{_\delta (\nu ,\beta ,\gamma )},\end{array}$$ (55) with the radiator given by $$\begin{array}{c}\hfill _\delta (\nu ,\beta ,\gamma )=[dk]W_\delta (k)\left[1\underset{\mathrm{}=1}{\overset{4}{}}u_{\mathrm{}}(k)\mathrm{\Theta }_{\mathrm{}}(k)\right].\end{array}$$ (56) As before, the unity in the square brackets has been included to account for the virtual correction contribution. Here $`\mathrm{\Theta }_{\mathrm{}}(k)`$ is the support function for a parton $`k`$ emitted in the quadrant $`C_{\mathrm{}}`$, and we have denoted $`u_1=u_{12}`$ and $`u_4=u_{13}`$. In Appendix B the radiators $`_\delta `$ are evaluated with SL accuracy and we obtain $$\begin{array}{c}\hfill _\delta =C_2^{(\delta )}r(\overline{\mu }_2,Q_2^2)+C_3^{(\delta )}r(\overline{\mu }_3,Q_3^2)+\frac{C_1^{(\delta )}}{\pi }_0^{\mathrm{}}\frac{dy}{1+y^2}\left[r(\overline{\mu }_{12},Q_1^2)+r(\overline{\mu }_{13},Q_1^2)\right],\end{array}$$ (57) where $`r`$ is the DL function defined in (44). Some comments on the colour charges, the scales and the various $`\overline{\mu }`$-variables are in order. * The radiator consists of three “independent” contributions from the radiation off each of three hard partons. Here $$C_a^{(a)}=N_c;C_a^{(b)}=C_F,\text{for}ab,$$ (58) is the colour charge of the parton $`P_a`$. * The hardness scale $`Q_a^2`$ in each term has a simple structure: it is determined by the invariant transverse momentum of the hard parton $`a`$ with respect to the dipole $`bc`$: $$\begin{array}{c}\hfill Q_a^2=\frac{p_{ta}^2}{4}e^{g_a},p_{ta}^2=2\frac{(P_bP_a)(P_aP_c)}{(P_bP_c)}.\end{array}$$ (59) This makes $`Q_a`$ depending on the event geometry. The scale $`Q_a^2`$ also includes an additional, geometry-independent, factor $`e^{g_a}`$ depending on the nature of the parton $`a`$. This factor takes into account a SL correction due to hard parton splitting, with $`g_a=3/2`$ for a quark ($`a=1,2`$ in (57)) and $`g_a=\beta _0/2N_c`$ for a gluon. Due to this factor the scales depend on the configuration $`\delta `$. Recall that the derivation of the exact form of the hard scales required a SL analysis which takes a due care of the inter-jet regions where the soft distribution does not have collinear singularities. As we have seen in the previous section, the hard scales for the right distribution (46) do not have such a simple interpretation because selecting one hemisphere is unnatural for the soft radiation pattern. As shown in Appendices B.2.1 and B.2.2, the boundary effects due to radiation at 90<sup>o</sup> which complicate the scales, cancel with SL accuracy in the total distribution. * The various functions $`\overline{\mu }`$ originate from the following substitutions in the integral over $`k_x`$ $`1u_a`$ $``$ $`\vartheta (k_x1/\overline{\mu }_a),\overline{\mu }_a=\overline{\nu }\sqrt{1+\beta _a^2}\overline{\nu }\mu _a,a=2,3`$ (60) $`1u_{12}`$ $``$ $`\vartheta (k_x1/\overline{\mu }_{12}),\overline{\mu }_{12}=\overline{\nu }\sqrt{(1i\gamma y)^2+\beta _{12}^2}\overline{\nu }\mu _{12}.`$ (61) $`1u_{13}`$ $``$ $`\vartheta (k_x1/\overline{\mu }_{13}),\overline{\mu }_{13}=\overline{\nu }\sqrt{(1+i\gamma y)^2+\beta _{13}^2}\overline{\nu }\mu _{13},`$ (62) with $`\overline{\nu }`$ given in (43). As in the case of the right radiator, such substitution is valid within SL accuracy. The $`\gamma `$-dependence enters only in the two $`\overline{\mu }_{1a}`$ parameters which are associated with the gluon emission off the parton $`P_1`$ in the right hemisphere. This is in agreement with the fact that the variable $`\gamma `$ is conjugate to $`q_{1y}`$. Opposite signs of $`\gamma `$ in $`\overline{\mu }_{12}`$ and $`\overline{\mu }_{13}`$ reflect the fact that the recoil momentum $`q_{1y}`$ is positive (negative) in the right-up (right-down) quadrant. All $`\nu `$-, $`\beta `$\- and $`\gamma `$-dependence is contained in the parameters $`\overline{\mu }`$. As before, we can simplify the $`\beta `$\- and $`\gamma `$-dependence by expanding the various terms to SL accuracy for large $`\nu `$. We can write $$_\delta =R_\delta (\overline{\nu })+r^{}(\overline{\nu })S_\delta (\beta ,\gamma ),$$ (63) where $`R_\delta `$ is the DL contribution, given by the sum of three standard antenna terms, $`r^{}`$ is the SL function defined in (47), and the coefficient $`S_\delta `$ carries all the $`\beta _a`$ and $`\gamma `$ dependence. $`R_\delta `$ and $`S_\delta `$ are given by $$\begin{array}{cc}& R_\delta (\overline{\nu })=C_1^{(\delta )}r(\overline{\nu },Q_1^2)+C_2^{(\delta )}r(\overline{\nu },Q_2^2)+C_3^{(\delta )}r(\overline{\nu },Q_3^2),\hfill \\ & S_\delta (\beta ,\gamma )=C_2^{(\delta )}\mathrm{ln}\mu _2+C_3^{(\delta )}\mathrm{ln}\mu _3+\frac{C_1^{(\delta )}}{\pi }_0^{\mathrm{}}\frac{dy}{1+y^2}\mathrm{ln}\left[\mu _{12}(y)\mu _{13}(y)\right],\hfill \end{array}$$ (64) with the charges given in (58). In detail the DL terms for the three configurations are $`R_3(\overline{\nu })`$ $`=`$ $`C_Fr(\overline{\nu },Q_1^2)+C_Fr(\overline{\nu },Q_2^2)+N_cr(\overline{\nu },Q_3^2),`$ (65) $`R_2(\overline{\nu })`$ $`=`$ $`C_Fr(\overline{\nu },Q_1^2)+N_cr(\overline{\nu },Q_2^2)+C_Fr(\overline{\nu },Q_3^2),`$ (66) $`R_1(\overline{\nu })`$ $`=`$ $`N_cr(\overline{\nu },Q_1^2)+C_Fr(\overline{\nu },Q_2^2)+C_Fr(\overline{\nu },Q_3^2).`$ (67) There are two sources of SL corrections in (63). The first is due to different hard scales in three terms $`r(\overline{\nu },Q_1)`$ in (65)–(67), which depend on the geometry of the three-jet event, that is on the values of $`T`$ and $`T_M`$. The second is the contribution proportional to $`r^{}`$ given by the sum of $`\mathrm{ln}\mu `$–terms which depend on $`\beta _a`$ and $`\gamma `$. These contributions are specific for a three-jet topology. At the same time, within the SL accuracy this correction is insensitive to the details of the event geometry, that is to the values of $`T,T_M1`$. In conclusion, the total $`K_{\text{out}}`$-distribution, to SL accuracy, can be expressed by the following Mellin integral: $$\mathrm{\Sigma }_\delta (K_{\text{out}})=\frac{d\nu }{2\pi i\nu }e^{\nu K_{\text{out}}}e^{R_\delta (\overline{\nu })}_\delta (\overline{\nu }),$$ where the SL prefactor $`_\delta (\overline{\nu })`$ is given by $$_\delta (\overline{\nu })=_{\mathrm{}}^{\mathrm{}}\frac{d\gamma }{2\pi }\underset{a=1}{\overset{3}{}}_{\mathrm{}}^{\mathrm{}}\frac{d\beta _a}{\pi }I(\beta ,\gamma )e^{r^{}(\overline{\nu })S_\delta (\beta ,\gamma )}.$$ (68) For example, the explicit expression for $`_3`$ reads $$\begin{array}{cc}\hfill _3(\overline{\nu })=& \frac{d\gamma }{2\pi }\underset{a=1}{\overset{3}{}}\frac{d\beta _a}{\pi }I(\beta ,\gamma )\left(1+\beta _2^2\right)^{{\scriptscriptstyle \frac{1}{2}}C_Fr^{}(\overline{\nu })}\left(1+\beta _3^2\right)^{{\scriptscriptstyle \frac{1}{2}}N_cr^{}(\overline{\nu })}\hfill \\ & \times \mathrm{exp}\left\{C_Fr^{}(\overline{\nu })\left\{\frac{1}{\pi }_0^{\mathrm{}}\frac{dy}{1+y^2}\mathrm{ln}\left[\mu _{12}(y)\mu _{13}(y)\right]\right\}\right\}.\hfill \end{array}$$ The function $`_\delta `$ is analysed in Appendix C. In the first order in $`\alpha _\text{s}\mathrm{ln}\nu `$ the factors $``$ become (see Appendix C.1) $$\begin{array}{cc}& _\delta =1(3C_F+N_c)\mathrm{ln}2r^{}(\overline{\nu }),\delta =2,3,\hfill \\ & _1=1(2C_F+2N_c)\mathrm{ln}2r^{}(\overline{\nu }).\hfill \end{array}$$ Different weights of the quark and gluon colour factors for the two cases have a simple explanation. Due to the kinematics of parton recoil (see (17)) contribution of the gluon radiation off the most energetic (right-hemisphere) parton $`P_1`$ is twice that off the left-hemisphere partons $`P_2`$ and $`P_3`$. As a result, the SL correction is proportional to $`2C_F+C_F+N_c`$ when $`P_1`$ is a quark/antiquark ($`\delta =2,3`$) and to $`2N_c+C_F+C_F`$ when it is a gluon ($`\delta =1`$). Integration over the Mellin variable $`\nu `$ can be done as before and we obtain $$\mathrm{\Sigma }_\delta (K_{\text{out}})e^{R_\delta \left(\overline{K}_{\text{out}}^1\right)}\frac{_\delta \left(\overline{K}_{\text{out}}^{}{}_{}{}^{1}\right)}{\mathrm{\Gamma }\left(1+R^{}(\overline{K}_{\text{out}}^1)\right)},\overline{K}_{\text{out}}e^{\gamma _E}K_{\text{out}},$$ (69) where $`R^{}`$ is the logarithmic derivative of $`R_\delta `$. To SL accuracy we can take $$R^{}(\overline{K}_{\text{out}}^1)=(2C_F+N_c)r^{}(\overline{K}_{\text{out}}^1),$$ (70) which is the same function for all configurations $`\delta `$. This is possible since, $`r^{}`$ being a SL function, the difference between the hard scales can be neglected at the level of the next-to-next-to-leading $`𝒪\left(\alpha _\text{s}\right)`$ correction. ### 4.1 Radiators in the quasi-2-jet limit In the 3-jet kinematics we have been considering, $`T_M\begin{array}{c}<\hfill \\ \hfill \end{array}T\begin{array}{c}<\hfill \\ \hfill \end{array}1`$, jet energies are comparable and relative angles between jets are large. In these circumstances three basic scales in (59) are of the same order, $`Q_1Q_2Q_3\begin{array}{c}<\hfill \\ \hfill \end{array}Q`$. Still, keeping precise scales in (65)–(67) is essential, since their deviation from the overall hardness parameter $`Q`$ in the DL radiator function $`r(\overline{\nu },Q_a^2)`$ produced a SL correction $`\delta ^{(1)}Rr^{}(\overline{\nu },Q)\mathrm{ln}(Q_a/Q)=𝒪\left(\alpha _\text{s}\mathrm{log}\nu \right)`$. At the same time, since the next order expansion terms are negligible, $`\delta ^{(2)}Rr^{\prime \prime }(\overline{\nu },Q)\mathrm{ln}^2(Q_a/Q)=𝒪\left(\alpha _\text{s}\right)`$, it is perfectly legitimate to present the answer in a form different from (65)–(67). For example, expanding the scales in the first two (quark) terms in $`R_3`$ around $`q^2=2P_1P_2`$ we obtain, instead of (65) $$R_3\mathrm{\hspace{0.25em}2}C_Fr\left(\frac{P_1P_2}{2}e^{\frac{3}{2}}\right)+N_cr(Q_3^2),Q_3^2=\frac{2(P_1P_3)(P_2P_3)}{(P_1P_2)}e^{\frac{\beta _0}{2N_c}},$$ (71) where we have suppressed the first argument $`\overline{\nu }`$ and used $`\mathrm{ln}\frac{p_{t1}^2}{2P_1P_2}=\mathrm{ln}\frac{p_{t2}^2}{2P_1P_2}`$ to cancel the linear expansion terms, see (59). Mismatch between the radiator in (65) and the right-hand side of (71) is $$R_3(\text{71})R_3(\text{65})=𝒪\left(\alpha _\text{s}\mathrm{ln}^2\frac{P_1P_3}{P_2P_3}\right).$$ Being equivalent within the adopted accuracy, the representations (65) and (71) start to significantly differ, however, when the jet configuration becomes 2-jet-like. Indeed, when the gluon jet $`P_3`$ becomes relatively soft and/or collinear to the quark $`P_2`$, we should expect the answer to correspond to the $`q\overline{q}`$-dominated radiation pattern, and the gluon contribution to disappear. The latter representation (71) correctly describes this situation: the quark-antiquark antenna contribution with the hardness scale equal to the invariant squared mass of the $`q\overline{q}`$ pair, $`q^2=2P_1P_2`$, takes on the job, while radiation off the gluon vanishes with decrease of the gluon transverse momentum, $`Q_3^2p_{t3}^2q^2Q^2`$. At the same time, the first representation goes hay-wire. For example, in the collinear limit, $`P_3P_2P_3P_1\begin{array}{c}<\hfill \\ \hfill \end{array}P_1P_2`$, one of the quark scales, $`Q_2^2`$, vanishes (together with $`Q_3^2`$) while the other formally goes to infinity, $`Q_1^2Q^2`$. The geometric mean of the quark scales, $`2P_1P_2=\sqrt{p_{1t}^2p_{2t}^2}`$, employed in (71) cures this unphysical behaviour. For two other jet configurations, with the gluon having an intermediate or the largest of the three parton energies, analogous representations of the radiators read $`R_2`$ $``$ $`2C_Fr\left({\displaystyle \frac{P_1P_3}{2}}e^{\frac{3}{2}}\right)+N_cr(Q_2^2),Q_2^2={\displaystyle \frac{2(P_1P_2)(P_2P_3)}{(P_1P_3)}}e^{\frac{\beta _0}{2N_c}};`$ (72) $`R_1`$ $``$ $`N_c\left[r\left({\displaystyle \frac{P_1P_2}{2}}e^{(\frac{3}{4}+\frac{\beta _0}{4N_c})}\right)+r\left({\displaystyle \frac{P_1P_3}{2}}e^{(\frac{3}{4}+\frac{\beta _0}{4N_c})}\right)\right]{\displaystyle \frac{1}{N_c}}r\left({\displaystyle \frac{P_2P_3}{2}}e^{\frac{3}{2}}\right).`$ (73) These expressions, contrary to the original ones (66) and (67), survive the limit of the left-hemisphere jets becoming quasi-collinear, $`P_2P_3P_1P_3\begin{array}{c}<\hfill \\ \hfill \end{array}P_1P_2`$. Indeed, the configuration $`\delta =2`$ is then identical, colour-wise, to $`\delta =3`$, and (72) is dominated by the $`q\overline{q}`$ two-jet contribution, while the non-Abelian part vanishes with the gluon transverse momentum, $`Q_2^2p_{2t}^2Q^2`$. A rare but interesting configuration $`\delta =1`$, where the gluon $`P_1`$ in the right hemisphere is balanced by a quasi-collinear $`q\overline{q}`$ pair in the left hemisphere, see (72), corresponds to the gluon-gluon system: radiation off the colour-octet $`q\overline{q}`$ pair is dominated by large-angle coherent bremsstrahlung proportional to the gluon charge, $`N_c`$, with the colour-suppressed $`1/N_c`$ correction term vanishing in the collinear limit. Thus, the radiators in the form of (65)–(67) are inapplicable in the quasi-two-jet kinematics. An attentive reader could have noticed that the modified expressions (71)–(72), though better behaved, cannot pretend to uniformly preserve the desired SL accuracy. Accommodating correctly SL effects due to large-angle soft bremsstrahlung, these expressions fail, however, to properly account for “hard” SL corrections in the collinear limit. For example, the right-hand side of (71) for the most natural configuration $`\delta =3`$ has a perfect soft–gluon limit, when $`2(P_1P_2)=Q^2`$ becomes the proper 2-jet scale. At the same time, when $`P_3`$ remains energetic but collinear, $`E_3\begin{array}{c}<\hfill \\ \hfill \end{array}E_2`$, $`(P_2P_3)0`$, the quark-jet scale in (71) remains smaller than the total annihilation energy, $`2(P_1P_2)<2(P_1(P_2+P_3))=Q^2`$. The mismatch amounts to a SL correction $`𝒪\left(r^{}\mathrm{ln}x_2\right)`$. ## 5 Discussion and conclusions In this paper we performed the all-order perturbative analysis of the out-of-plane transverse momentum distributions in three-jet $`e^+e^{}`$ annihilation events. We considered the total $`K_{\text{out}}`$ of the event and $`K_{\text{out}}^R`$ accumulated in one event-hemisphere which contains the most energetic of three jets. The perturbative expression for the integrated $`K_{\text{out}}`$-distribution in the right (one-jet) hemisphere is given in (44)–(49) and (124). The $`K_{\text{out}}`$-total distribution is determined by expressions (57)–(70). These answers resum all double- and single-logarithmic contributions to the exponent (the so-called “radiator”) of the distribution in the standard Mellin-Fourier parameter space. DL and SL contributions to the radiator can be formally represented, respectively, as series $`\alpha _\text{s}^n\mathrm{log}^{n+1}(Q/K_{\text{out}})`$ and $`\alpha _\text{s}^n\mathrm{log}^n(Q/K_{\text{out}})`$ ($`n1`$). Throughout the analysis, we systematically neglected next-to-next-to-leading terms in the radiator (radiator series $`\alpha _\text{s}^n\mathrm{log}^{n1}\mathrm{ln}(Q/K_{\text{out}})`$, $`n1`$) as well as non-exponentiating corrections of the relative order $`𝒪\left(\alpha _\text{s}\right)`$. The latter belong to the coefficient function $`1+c\alpha _\text{s}+𝒪\left(\alpha _\text{s}^2\right)`$. Its first coefficient, $`c`$, is analysed in where the order $`\alpha _\text{s}^2`$ expansion of the approximate resummed cross is compared with numerical calculation based on the exact $`\alpha _\text{s}^2`$ matrix element. Matching the resummed logarithmic expressions with the exact result is necessary for justifying exponentiation of the next-to-leading SL terms $`\alpha _\text{s}\mathrm{ln}(Q/K_{\text{out}})`$. Accompanying gluon radiation pattern follows the colour topology of the underlying parton antenna. The corresponding patterns were known, to SL accuracy, for two-parton — $`q\overline{q}`$ and $`gg`$ — sources (quark and gluon form factors, 2-jet shapes). In the present paper we presented the first such analysis for 3-parton ensembles. The $`q\overline{q}g`$–initiated events possess a rich colour structure which determines secondary parton flows and makes them event-geometry-dependent. We have demonstrated that after taking into account SL effects due to inter-jet gluon flows, the result can be cast as a sum of quark-antiquark ($`C_F`$) and gluon ($`N_c`$) contributions with the proper hardness scales depending on $`T`$ and $`T_M`$. A significant dependence on event geometry is the key feature which singles out the $`K_{\text{out}}`$ distribution among other $`e^+e^{}`$ event shape observables. Soft gluon field components with small transverse momenta (“gluers” with $`k_t\begin{array}{c}>\hfill \\ \hfill \end{array}\mathrm{\Lambda }_{QCD}`$) are believed to be responsible for hadronization. This belief, known as hypothesis of the local parton-hadron duality, has been verified in a number of experimental studies of various features of multiple hadroproduction in hard processes (for review see ). Moreover, in recent years a new theoretical techniques have been developed for triggering and quantifying genuine confinement effects by studying gluer radiation. A rich dependence of gluon (and, therefore, gluer) radiation on the colour topology makes 3-jet observables, and $`K_{\text{out}}`$ in particular, an interesting field for the study of non-perturbative effects. A separate publication will be devoted to analysis of the leading non-perturbative power-suppressed corrections to 3-jet observables. The $`K_{\text{out}}`$ distribution, similar to the case of 2-jet Broadening observables , will be shown to exhibit $`\mathrm{log}K_{\text{out}}`$–enhanced $`1/Q`$ contribution with the magnitude depending on 3-jet geometry, that is on $`T`$ and $`T_M`$ values. To access this interesting physics experimental studies of $`K_{\text{out}}`$ should be carried out specifically for events with moderate values of $`1T`$, away from the 2-jet region. Two sorts of experimental studies of these predictions can be envisaged. The most straightforward comparison calls for experimental identification of the gluon jet. Gluon tagging, however, is unnecessary for the study of the total $`K_{\text{out}}`$-distribution for genuine 3-jet events: the prediction given by the sum of three jet configurations corresponding to given $`T`$ and $`T_M`$, weighted with the proper 3-jet Born cross section factors bears practically as much information. Tagging brings in more information when a single-jet (right-jet) $`K_{\text{out}}`$-distribution is studied. In this case essentially different $`K_{\text{out}}`$-spectra will be seen depending on whether the right-hemisphere parton is a quark or a gluon. ## Acknowledgements We are grateful to Gavin Salam for helpful discussions and suggestions. ## Appendix A Kinematics ### A.1 Kinematics of $`q\overline{q}g`$ jets We consider the quark, antiquark and gluon with skeleton momenta $`P_a`$ in the centre-of-mass frame, $`_{a=1}^3P_a=Q=(Q,0,0,0)`$, $$P_1=E(x_1,0,0,t_1),P_2=E(x_2,0,T_M,t_2),P_3=E(x_3,0,T_M,t_3).$$ (74) We have $`2E=Q`$, $`x_a=2(P_aQ)/Q^2`$, $`x_1+x_2+x_3=2`$ and $`t_1=t_2+t_3`$. Assuming $`x_2>x_3`$, we obtain $$\begin{array}{cc}& x_1=T,t_1=T,\hfill \\ & x_2=\frac{2T}{2}+\frac{T}{2}\rho ,t_2=\frac{T}{2}+\frac{2T}{2}\rho ,\hfill \\ & x_3=\frac{2T}{2}\frac{T}{2}\rho ,t_3=\frac{T}{2}\frac{2T}{2}\rho ,\hfill \end{array}$$ (75) where $$\rho \sqrt{1\frac{T_M^2}{1T}}<1.$$ (76) Thrust-major is confined to the kinematical region $$\frac{2(1T)}{T}\sqrt{2T1}<T_M<\sqrt{1T},$$ (77) where the upper limit comes from reality of $`\rho `$, and the lower limit comes from requiring $`x_1>x_2`$. It is straightforward to show that in this kinematical region $`t_3>0`$, i.e. that both $`P_2`$ and $`P_3`$ lie in the left hemisphere. We introduce Sudakov variables based on two light-light vectors aligned with the thrust axis, $$P=E(1,0,0,1),\overline{P}=E(1,0,0,1).$$ (78) The jet momenta have the following Sudakov decomposition: $$P_1=TP,P_2=A_2P+B_2\overline{P}+P_t,P_3=A_3P+B_3\overline{P}P_t,$$ (79) where the transverse momentum with respect to the thrust axis is $`P_t=(0,0,ET_M,0)`$. In terms of $`T`$ and $`T_M`$ the longitudinal Sudakov momentum components are $$\begin{array}{cc}& A_2=\frac{1T}{2}(1\rho ),B_2=\frac{1}{2}(1+\rho ),\hfill \\ & A_3=\frac{1T}{2}(1+\rho ),B_3=\frac{1}{2}(1\rho ).\hfill \end{array}$$ (80) Since $`P_a`$ with $`a=2,3`$ belong to the left hemisphere, we have $`A_a<B_a`$. Denoting by $`\mathrm{\Theta }_a`$ the angle of $`\stackrel{}{P}_a`$ with respect to the thrust axis we introduce the angular variables $`\tau _a`$ $$\begin{array}{cc}& \frac{\tau _2}{Q}\mathrm{tan}\frac{\mathrm{\Theta }_2}{2}=\frac{P_t}{QA_2}=\frac{1+\rho }{T_M},\hfill \\ & \frac{\tau _3}{Q}\mathrm{tan}\frac{\mathrm{\Theta }_3}{2}=\frac{P_t}{QA_3}=\frac{1\rho }{T_M}.\hfill \end{array}$$ (81) We remark that $`\tau _3`$ is negative since $`P_3`$ lies in the third quadrant (“down” hemisphere). Since $`P_2`$ and $`P_3`$ are in the left hemisphere and $`\left\|P_{2z}\right\|>\left\|P_{3z}\right\|`$, we have $`0<\pi \mathrm{\Theta }_2<\mathrm{\Theta }_3\pi <\frac{1}{2}\pi `$. Hence, $$\frac{\tau _2}{Q}>\frac{\tau _3}{Q}>1.$$ The following relations hold: $$\tau _2\left\|\tau _3\right\|=\frac{Q^2}{1T},\tau _2+\left\|\tau _3\right\|=\frac{2Q}{T_M}.$$ (82) The minimal value of thrust-major, $`T_M=2(1T)/(2T)`$ for a given $`T`$ in (77), corresponds to the configuration with the softest parton momentum orthogonal to the thrust axis, $$\rho \frac{T}{2T},\frac{\tau _2}{Q}\frac{1}{1T},\frac{\left\|\tau _3\right\|}{Q}\mathrm{\hspace{0.25em}1}.$$ The maximal value $`T_M=\sqrt{1T}`$ is achieved with a symmetric $`23`$ pair: $$\rho 0,\frac{\tau _2}{Q}\frac{\left\|\tau _3\right\|}{Q}\frac{1}{\sqrt{1T}}.$$ We also introduce two-dimensional vectors in the transverse plane $`\{x,y\}`$, $$\begin{array}{c}\hfill \stackrel{}{\tau }_2\frac{\stackrel{}{P}_t}{A_2}=(0,\tau _2),\stackrel{}{\tau }_3\frac{\stackrel{}{P}_t}{A_3}=(0,\tau _3).\end{array}$$ (83) Since $`P_1`$ is along the thrust axis we have $`\stackrel{}{\tau }_1=(0,0)`$. ### A.2 Born cross sections The squared Born matrix element $`M_0(𝒞_\delta )`$ has a well-known expression in terms of the variables $`x_a`$. For $`\delta =3`$ (with $`P_3`$ the gluon momentum) we have $$M_0^2(𝒞_3)=\frac{C_F}{2\pi }\alpha _\text{s}(Q)\frac{x_1^2+x_2^2}{(1x_1)(1x_2)}.$$ (84) The phase space factor $`\mathrm{\Gamma }_0`$ for the Born 3-parton system is $$\begin{array}{cc}\hfill \mathrm{\Gamma }_0& =\frac{1}{8(2\pi )^5}\underset{a=1}{\overset{3}{}}\frac{d^3p_a}{E_a}\delta ^4(Qp_a)\delta ^3(\stackrel{}{p}_1\stackrel{}{P}_1)\delta ^2(\stackrel{}{p}_{t2}\stackrel{}{P}_{t2})\hfill \end{array}$$ (85) The dipole invariant masses in terms of $`T,T_M`$ and variables $`\tau _i`$ are $$\begin{array}{cc}& Q_{12}^2=2P_1P_2=\frac{T}{2}(1+\rho )Q^2=\frac{1}{2}TT_M\tau _2Q,\hfill \\ & Q_{13}^2=2P_1P_3=\frac{T}{2}(1\rho )Q^2=\frac{1}{2}TT_M\|\tau _3\|Q,\hfill \\ & Q_{23}^2=2P_2P_3=(1T)Q^2=\frac{Q^4}{\tau _2\|\tau _3\|}.\hfill \end{array}$$ (86) The scales in (59), $$p_{ta}^2\frac{2(P_bP_a)(P_aP_c)}{(P_bP_c)},$$ — the invariant transverse momentum of the hard parton $`P_a`$ with respect to the $`bc`$-dipole — read $$p_{t1}^2=\frac{Q^2}{2}\left(\frac{TT_M}{1T}\right)^2,p_{t2}^2=Q^2(1T)\left(\frac{1+\rho }{T_M}\right)^2,p_{t3}^2=Q^2(1T)\left(\frac{1\rho }{T_M}\right)^2.$$ (87) ### A.3 Soft partons For secondary massless parton of momentum $`k`$ we write the Sudakov representation $$k=\alpha P+\beta \overline{P}+k_t,\alpha \beta =\frac{k_t^2}{Q^2},$$ (88) with $`\stackrel{}{k}_t`$ the transverse momentum with respect to the thrust axis. The right-hemisphere condition imposes the restriction upon longitudinal parton components, $$\alpha >\beta \alpha >\frac{k_t}{Q}.$$ Probability of soft gluon emission off the $`P_a,P_b`$ dipole is described by the squared matrix element $$\frac{P_aP_b}{2(P_ak)(kP_b)}=\frac{\alpha ^2\left(\stackrel{}{\tau }_a\stackrel{}{\tau }_b\right)^2}{\left(\stackrel{}{k}_t\alpha \stackrel{}{\tau }_a\right)^2\left(\stackrel{}{k}_t\alpha \stackrel{}{\tau }_b\right)^2}.$$ (89) ## Appendix B Radiators Here we compute the radiators for the total $`K_{\text{out}}`$ distribution, $`_\delta (\nu ,\beta ,\gamma )`$, to SL accuracy. The radiator is given by the combination of dipole contributions, $$r_{ab}=[dk]w_{ab}(k)\underset{\mathrm{}=1}{\overset{4}{}}[1u_{\mathrm{}}]\mathrm{\Theta }_{\mathrm{}}(k),$$ (90) with $`\mathrm{\Theta }_{\mathrm{}}(k)`$ the support function restricting gluon momentum $`k`$ to the quadrant $`C_{\mathrm{}}`$ and the source functions $`u_{\mathrm{}}`$ defined in (52), (53). Expressing the parton phase space $`[dk]`$ in terms of Sudakov variables and invoking the soft distribution $`w_{ab}`$ in (89) we have $$r_{ab}=_Q^Q𝑑k_x_{K_{ym}}^{K_{ym}}\frac{dk_y}{\pi }_0^{\alpha _m}\frac{d\alpha }{\alpha }\left\{\frac{\alpha _\text{s}(k_{t,ab})}{\pi }\frac{\alpha ^2(\stackrel{}{\tau }_a\stackrel{}{\tau }_b)^2}{(\stackrel{}{k}_t\alpha \stackrel{}{\tau }_a)^2(\stackrel{}{k}_t\alpha \stackrel{}{\tau }_b)^2}\right\}\underset{\mathrm{}}{}[1u_{\mathrm{}}]\mathrm{\Theta }_{\mathrm{}}(k).$$ (91) We recall that $`\stackrel{}{\tau }_a`$ are jet transverse direction vectors, $`\stackrel{}{\tau }_1=0`$ and $`\stackrel{}{\tau }_{2,3}`$ given in (83). We have three integrations to perform. The result of the first two integrals, in $`\alpha `$ and $`k_y`$, has the structure $$k_y𝑑\alpha \frac{\alpha _\text{s}}{\left\|k_x\right\|}\left(\mathrm{ln}\frac{Q^{}}{\left\|k_x\right\|}+𝒪\left(k_x\right)\right),\frac{Q^{}}{Q}=𝒪\left(1\right),$$ (92) where we have omitted $`𝒪\left(k_x\right)`$ terms as producing non-logarithmic $`𝒪\left(\alpha _\text{s}\right)`$ corrections upon $`k_x`$-integration. (For the same reason we have set the lower limit of $`\alpha `$-integration, $`\alpha \begin{array}{c}>\hfill \\ \hfill \end{array}(k_t/Q)^2`$, to zero, since the integral converges.) Similar correction originates from the precise limit of the $`k_x`$ integration, $`\left\|k_x\right\|\begin{array}{c}<\hfill \\ \hfill \end{array}Q`$. Therefore we were free to set it to $`Q`$ in (91). <sup>2</sup><sup>2</sup>2 We remind the reader that all these $`𝒪\left(\alpha _\text{s}\right)`$ uncertainties, as well as those we shall encounter below, are related to the coefficient function . At the same time, the exact limits of $`\alpha `$ and $`k_y`$ integrals do matter when they affect the scale of the logarithm, $`Q^{}`$, in (92) which are responsible for the SL correction to the radiator. The logarithmic terms in (92) originate from three “collinear regions” when the gluon transverse direction vector $`\stackrel{}{\tau }=\stackrel{}{k}_t/\alpha `$ is close to one of $`\stackrel{}{\tau }_a`$. In the direction $`P_1`$ it is the $`\alpha `$-integration that produces a collinear logarithm. Noticing that $`P_{1z}=TQ/2`$, we then set $`\alpha _m`$ to the maximal value of $`\alpha `$ kinematically allowed in the right-hemisphere jet, $$\alpha _m=T.$$ (93) In the regions collinear to the left-hemisphere jets $`a=2,3`$, $`\alpha `$ and $`k_y`$ are linked by the condition $`\stackrel{}{\tau }\stackrel{}{\tau }_a`$. In this case the kinematical limit can be expressed in terms of the $`y`$-component: $`P_{2y}=\left\|P_{3y}\right\|=T_MQ/2`$ translates into $$K_{ym}=\frac{T_MQ}{2}.$$ (94) In (91) we employed soft radiation probabilities $`w_{ab}`$. However, in the collinear regions hard parton splitting should be accounted for which produces another important SL corrections similar to those coming from the kinematical limits (93) and (94). These corrections will be taken care of in the end of Appendix B.2.2. We now proceed with successive integrations. Since the sources $`u`$ do not depends on $`\alpha `$ we first compute the $`\alpha `$-integral. Then we compute the $`k_y`$\- and $`k_x`$-integrals. ### B.1 Integrating over $`\alpha `$ We calculate separately contributions from the right ($`R`$) and left ($`L`$) hemispheres, $$R\frac{k_t}{Q}<\alpha <\alpha _m=T,L\alpha <\frac{k_t}{Q},$$ (95) and define $$\begin{array}{c}\hfill I_{ab}^R=_{k_t/Q}^{\alpha _m}\frac{\alpha d\alpha (\stackrel{}{\tau }_a\stackrel{}{\tau }_b)^2}{(\stackrel{}{k}_t\alpha \stackrel{}{\tau }_a)^2(\stackrel{}{k}_t\alpha \stackrel{}{\tau }_b)^2},I_{ab}^L=_0^{k_t/Q}\frac{\alpha d\alpha (\stackrel{}{\tau }_a\stackrel{}{\tau }_b)^2}{(\stackrel{}{k}_t\alpha \stackrel{}{\tau }_a)^2(\stackrel{}{k}_t\alpha \stackrel{}{\tau }_b)^2}.\end{array}$$ (96) The $`\alpha `$-integration yields $$\begin{array}{cc}& _0^\alpha \frac{\alpha d\alpha (\stackrel{}{\tau }_a\stackrel{}{\tau }_b)^2}{(\stackrel{}{k}\alpha \stackrel{}{\tau }_a)^2(\stackrel{}{k}\alpha \stackrel{}{\tau }_b)^2}\hfill \\ & =\frac{1}{k_y^2+h_{ab}^2k_x^2}\{\frac{h_{ab}}{2}\mathrm{ln}\frac{(\alpha \tau _ak_y)^2+k_x^2}{(\alpha \tau _bk_y)^2+k_x^2}+\frac{k_y}{\|k_x\|}\left(\mathrm{arctan}\frac{\alpha \tau _ak_y}{\|k_x\|}+\mathrm{arctan}\frac{\alpha \tau _bk_y}{\|k_x\|}\right)\}.\hfill \end{array}$$ (97) Here $`h_{ab}`$ is a function of $`\tau _a`$ and $`\tau _b`$ which has the following geometrical meaning: $$h_{ab}=\frac{\tau _a+\tau _b}{\tau _a\tau _b}=\frac{\mathrm{tan}\frac{1}{2}\theta _a+\mathrm{tan}\frac{1}{2}\theta _b}{\mathrm{tan}\frac{1}{2}\theta _a\mathrm{tan}\frac{1}{2}\theta _b}=\frac{\mathrm{sin}\frac{1}{2}(\theta _a+\theta _b)}{\mathrm{sin}\frac{1}{2}(\theta _a\theta _b)},$$ (98) with $`\theta _a`$ the angle between the momentum $`\stackrel{}{P}_a`$ and the thrust axis. We have $`h_{a1}=1`$ ($`a=2,3`$) and $`h_{23}=\rho `$. In the limit of small aplanarity, $`\left\|k_x\right\|Q`$, for the dipoles involving the right-hemisphere jet $`P_1`$ we derive $`I_{1a}^R`$ $`=`$ $`{\displaystyle \frac{1}{k_t^2}}\left\{{\displaystyle \frac{1}{2}}\mathrm{ln}{\displaystyle \frac{\alpha _m^2\tau _a^2}{\kappa _a^2+k_x^2}}+{\displaystyle \frac{k_y}{\left\|k_x\right\|}}\left({\displaystyle \frac{\pi }{2}}\left(12\vartheta (\tau _a)\right)\mathrm{arctan}{\displaystyle \frac{\kappa _a}{\left\|k_x\right\|}}\right)\right\},`$ (99) $`I_{1a}^L`$ $`=`$ $`{\displaystyle \frac{1}{k_t^2}}\left\{{\displaystyle \frac{1}{2}}\mathrm{ln}{\displaystyle \frac{\kappa _a^2+k_x^2}{k_t^2}}+{\displaystyle \frac{k_y}{\left\|k_x\right\|}}\left(\mathrm{arctan}{\displaystyle \frac{\kappa _a}{\left\|k_x\right\|}}+\mathrm{arctan}{\displaystyle \frac{k_y}{\left\|k_x\right\|}}\right)\right\}.`$ (100) For the left-hemisphere dipole we get $$\begin{array}{cc}\hfill I_{23}^R=& \frac{1}{k_y^2+\rho ^2k_x^2}\left\{\frac{\rho }{2}\mathrm{ln}\frac{\tau _2^2}{\tau _3^2}+\frac{\rho }{2}\mathrm{ln}\frac{\kappa _3^2+k_x^2}{\kappa _2^2+k_x^2}\frac{k_y}{\left\|k_x\right\|}\left(\mathrm{arctan}\frac{\kappa _2}{\left\|k_x\right\|}+\mathrm{arctan}\frac{\kappa _3}{\left\|k_x\right\|}\right)\right\},\hfill \\ \hfill I_{23}^L=& \frac{1}{k_y^2+\rho ^2k_x^2}\{\frac{\rho }{2}\mathrm{ln}\frac{\kappa _2^2+k_x^2}{\kappa _3^2+k_x^2}\hfill \\ & +\frac{k_y}{\left\|k_x\right\|}(\mathrm{arctan}\frac{\kappa _2}{\left\|k_x\right\|}+\mathrm{arctan}\frac{\kappa _3}{\left\|k_x\right\|}+2\mathrm{arctan}\frac{k_y}{\left\|k_x\right\|})\}.\hfill \end{array}$$ (101) Here $$\kappa _a\frac{k_t}{Q}\tau _ak_y,\kappa _2>0,\kappa _3<0.$$ ### B.2 Integrating over $`k_y`$ To obtain the $`k_x`$-integrand we need to perform the $`k_y`$-integration $$\begin{array}{c}\hfill r_{ab}=_Q^Q𝑑k_x\frac{\alpha _\text{s}(2\|k_x\|)}{\pi }B_{ab}(k_x).\end{array}$$ (102) The origin of the running coupling argument in (102) is explained below in Appendix B.5. To compute the functions $`B_{ab}(k_x)`$ we consider two terms $`B_{ab}^{(U)}`$ and $`B_{ab}^{(D)}`$ coming from the upper ($`U`$) and lower ($`D`$) hemispheres, each containing left- and right-hemisphere contributions, $$B_{ab}^{(U)}=_0^{K_{ym}}\frac{dk_y}{\pi }\left\{I_{ab}^L\left[1u_2\right]+I_{ab}^R\left[1u_{12}\right]\right\},$$ (103) and $$B_{ab}^{(D)}=_{K_{ym}}^0\frac{dk_y}{\pi }\left\{I_{ab}^L\left[1u_3\right]+I_{ab}^R\left[1u_{13}\right]\right\}.$$ (104) From (52) we have $$u_{12}=u(\beta _{12},\gamma ),u_{13}=u(\beta _{13},\gamma ),u_2=u(\beta _2,0),u_3=u(\beta _3,0),$$ (105) where $`\beta _{1a}=\beta _1+\beta _a`$ and $$u(\beta ,\gamma )=e^{\nu [\left\|k_x\right\|+i\beta k_x+i\gamma \left\|k_y\right\|]}.$$ (106) Only the right-hemisphere sources $`u_{1a}`$ depend on $`k_y`$. #### B.2.1 Dipole $`23`$ Consider first the $`23`$-dipole contributions $`B_{23}^{(U/D)}`$. The DL piece is obviously contained in the left-hemisphere piece $`I_{23}^L`$. At the same time, the contributions involving $`I_{23}^R`$ are subleading: their $`k_y`$-integrals are not enhanced by $`\mathrm{log}\left\|k_x\right\|`$. Therefore, with SL accuracy the accompanying $`[1u_{1a}]`$ factors in (103) and (104) can be simplified, set equal to $`[1u_a]`$ and factored out: $$I_{23}^L[1u_a]+I_{23}^R[1u_{1a}](I_{23}^L+I_{23}^R)[1u_a].$$ Then, the terms in $`I_{23}^R`$ and $`I_{23}^L`$ due to the boundary between the $`R`$\- and $`L`$-hemispheres cancel in the sum, and, using that the sources are now $`k_y`$-independent, we obtain $$\begin{array}{cc}\hfill B_{23}^{(U)}& =\left[1u_2\right],B_{23}^{(D)}=\left[1u_3\right];\hfill \\ \hfill & =_0^{K_{ym}}\frac{dk}{\pi (k^2+\rho ^2k_x^2)}\left(\frac{\rho }{2}\mathrm{ln}\frac{\tau _2^2}{\tau _3^2}+\frac{2k}{\left\|k_x\right\|}\mathrm{arctan}\frac{k}{\left\|k_x\right\|}\right),\hfill \end{array}$$ (107) where $`k=k_y`$ ($`k=k_y`$) in the up (down) hemisphere. We extract the logarithmic contribution coming from the $`\mathrm{arctan}`$ term in the region $`\left\|k_x\right\|kK_m=\frac{1}{2}T_MQ`$ (see (94)) and using $$_0^{\mathrm{}}\frac{2kdk}{\pi (k^2+\rho ^2k_x^2)}\mathrm{arctan}\frac{\left\|k_x\right\|}{k}=\mathrm{ln}\frac{1+\rho }{\rho },$$ (108) we finally arrive at $$=\frac{1}{2}\mathrm{ln}\frac{(1T)Q^2}{4k_x^2},$$ (109) where we have used the relations $$T_M^2=(1T)(1\rho ^2)\text{and}\frac{\tau _2}{\|\tau _3\|}=\frac{1+\rho }{1\rho },$$ following from (76) and (81), respectively. #### B.2.2 Dipoles $`12`$ and $`13`$ The contributions $`1a`$ can be simplified in a similar way. The left-hemisphere contribution (100) is accompanied by $`k_y`$-independent source $`u_a`$. The same source can be attributed, however, to the second term in (99) as well, since the $`k_y`$-integration is non-logarithmic and produces a subleading contribution. After this simplification the $`R/L`$-boundary terms cancel in the sum of $`I_{1a}^R`$ and $`I_{1a}^L`$, and we arrive at $`B_{1a}^{(U)}`$ $`=`$ $`{\displaystyle _0^{K_{ym}}}{\displaystyle \frac{dk_y}{\pi k_t^2}}({\displaystyle \frac{1}{2}}\mathrm{ln}{\displaystyle \frac{\alpha _m^2\tau _a^2}{k_t^2}}\left[1u_{12}\right]+{\displaystyle \frac{k_y}{\left\|k_x\right\|}}[{\displaystyle \frac{\pi }{2}}sgn(\tau _a)+\mathrm{arctan}{\displaystyle \frac{k_y}{\left\|k_x\right\|}}]\left[1u_2\right]),`$ (110) $`B_{1a}^{(D)}`$ $`=`$ $`{\displaystyle _{K_{ym}}^0}{\displaystyle \frac{dk_y}{\pi k_t^2}}({\displaystyle \frac{1}{2}}\mathrm{ln}{\displaystyle \frac{\alpha _m^2\tau _a^2}{k_t^2}}\left[1u_{13}\right]+{\displaystyle \frac{k_y}{\left\|k_x\right\|}}[{\displaystyle \frac{\pi }{2}}sgn(\tau _a)+\mathrm{arctan}{\displaystyle \frac{k_y}{\left\|k_x\right\|}}]\left[1u_3\right]).`$ (111) As before, it is important to keep the dependence of the sources on the finite parameters $`\beta `$ and $`\gamma `$ only in the contributions enhanced by $`\mathrm{ln}(Q/\left\|k_x\right\|)`$. Such factor (collinear logarithm due to the jet #1) is explicitly present in the first terms in (110), (111). Adding these contributions together and integrating over $`k_y`$ results in substituting the R-hemisphere sources by the average source, $$\frac{[1\overline{u}_1]}{2\left\|k_x\right\|}\mathrm{ln}\frac{Q^2}{k_x^2},$$ where $$1\overline{u}_1\frac{1}{\pi }_0^{\mathrm{}}\frac{dy}{1+y^2}\left[(1u_{12})+(1u_{13})\right],y\left\|k_y/k_x\right\|.$$ Collinear enhancement factor due to the jet #2 originates from logarithmic $`k_y`$-integration in the second term of (110) (second quadrant), due to the jet #3 — in the second term of (111) (third quadrant). Remaining finite pieces can be evaluated using (108). Setting the limits $`\alpha _m=T`$ and $`K_{ym}=T_MQ/2`$ according to (93) and (94) and using (81), we finally obtain $$B_{1a}=B_{1a}^{(U+D)}=\frac{(1\overline{u}_1)}{2\left\|k_x\right\|}\mathrm{ln}\frac{Q^2}{k_x^2}+\frac{(1u_a)}{2\left\|k_x\right\|}\mathrm{ln}\frac{Q^2}{k_x^2}+\frac{(1u_0)}{\left\|k_x\right\|}\mathrm{ln}\frac{TT_M\|\tau _a\|}{8Q},$$ (112) where $`u_0`$ in the last subleading term is a source whose $`\beta `$-, $`\gamma `$-dependence can be chosen arbitrary. Using this freedom we can absorb the last term in (112) into rescaling of the first two namely, $$\mathrm{ln}\frac{Q^2}{k_x^2}\mathrm{ln}\frac{TT_M\|\tau _a\|Q}{8k_x^2}.$$ (113) We observe that the hard scales in (109) and (113) have a simple geometrical interpretation. Indeed, $$(1T)Q^2=\mathrm{\hspace{0.25em}2}P_2P_3=Q_{23}^2,\frac{TT_M\|\tau _a\|Q}{2}=2P_1P_a=Q_{1a}^2,$$ (114) with $`Q_{ab}`$ the invariant dipole masses, see (86). Our analysis was based up to now on the soft radiation matrix element, (89). To fully take into account SL effects from the region of large secondary parton momenta, we have, in addition to fixing the upper limits $`\alpha _m`$ and $`k_{ym}`$, to consider also hard collinear parton splitting. Due to collinear factorization, these corrections are process-independent and can be easily taken into account by proper rescaling of jet hardness parameters. They amount to supplying the invariant dipole masses by the factors $$Q_{ab}^2Q_{ab}^2\mathrm{exp}\left\{\frac{1}{2}(g_a+g_b)\right\},g_a=\{\begin{array}{cc}\frac{3}{2}& \text{for a quark/antiquark},\hfill \\ \frac{\beta _0}{2N_c}& \text{for a gluon}.\hfill \end{array}$$ (115) We finally obtain $$\begin{array}{cc}\hfill B_{23}& =\frac{\left[(1u_2)+(1u_3)\right]}{2\left\|k_x\right\|}\mathrm{ln}\frac{Q_{23}^2e^{{\scriptscriptstyle \frac{1}{2}}(g_2+g_3)}}{4k_x^2},\hfill \\ \hfill B_{1a}& =\frac{\left[(1\overline{u}_1)+(1u_a)\right]}{2\left\|k_x\right\|}\mathrm{ln}\frac{Q_{1a}^2e^{{\scriptscriptstyle \frac{1}{2}}(g_1+g_a)}}{4k_x^2}.\hfill \end{array}$$ (116) ### B.3 Radiator by assembling bits and pieces Now that the $`\alpha `$\- and $`k_y`$-integrations have been performed, we are in a position to assemble the full radiators for three jet configurations: $`_3`$ $`=`$ $`{\displaystyle _Q^Q}𝑑k_x{\displaystyle \frac{\alpha _\text{s}(2\|k_x\|)}{\pi }}{\displaystyle \frac{N_c}{2}}\left(B_{13}+B_{23}{\displaystyle \frac{1}{N_c^2}}B_{12}\right),`$ (117) $`_2`$ $`=`$ $`{\displaystyle _Q^Q}𝑑k_x{\displaystyle \frac{\alpha _\text{s}(2\|k_x\|)}{\pi }}{\displaystyle \frac{N_c}{2}}\left(B_{12}+B_{23}{\displaystyle \frac{1}{N_c^2}}B_{13}\right),`$ (118) $`_1`$ $`=`$ $`{\displaystyle _Q^Q}𝑑k_x{\displaystyle \frac{\alpha _\text{s}(2\|k_x\|)}{\pi }}{\displaystyle \frac{N_c}{2}}\left(B_{12}+B_{13}{\displaystyle \frac{1}{N_c^2}}B_{23}\right).`$ (119) The answers assume a simple when expressed in terms of the invariant transverse momentum of the hard parton $`P_a`$ with respect to the $`bc`$-dipole $$p_{ta}^2\frac{Q_{ba}^2Q_{ac}^2}{Q_{bc}^2}.$$ It is straightforward to verify that, to SL accuracy, the answer can be represented in a symmetric form as $$_\delta =_Q^Q\frac{dk_x}{k_x}\frac{\alpha _\text{s}(2\|k_x\|)}{2\pi }\left(C_1^{(\delta )}\mathrm{ln}\frac{Q_1^2}{k_x^2}[1\overline{u}_1]+C_2^{(\delta )}\mathrm{ln}\frac{Q_2^2}{k_x^2}[1u_2]+C_3^{(\delta )}\mathrm{ln}\frac{Q_3^2}{k_x^2}[1u_3]\right),$$ (120) where the hard scales are given by: $$\begin{array}{c}\hfill Q_1^2=\frac{p_{t1}^2}{4}e^{g_1},Q_2^2=\frac{p_{t2}^2}{4}e^{g_2},Q_3^2=\frac{p_{t3}^2}{4}e^{g_3}.\end{array}$$ (121) The rule for the colour factors in (120) is simple: $$C_a^{(a)}=N_c;C_b^{(a)}=C_F,\text{for}ab,$$ (122) and the hard-splitting rescaling factors $`e^{g_a}`$ are defined in (115). The universal representations (117)–(119) do not have a smooth 2-jet limit. In subsection 4.1 alternative formulae are presented which are equivalent to the previous ones in the 3-jet kinematics, $`T_M\begin{array}{c}<\hfill \\ \hfill \end{array}T=𝒪\left(1\right)`$, but are better behaved when the system assumes a quasi-two-jet kinematics, $`T_MT`$. ### B.4 Radiator for the right distribution In this case the source does not depend on $`k_y`$, and we have $$\begin{array}{cc}& _{ab}^R(\nu ,\beta )=_Q^Q𝑑k_x\frac{\alpha _\text{s}(2\|k_x\|)}{\pi }B_{ab}^R(k_x),\hfill \\ & B_{ab}^R(k_x)=\left[1e^{\nu (\|k_x\|+i\beta k_x)}\right]_{\mathrm{}}^{\mathrm{}}\frac{dk_y}{\pi }I_{ab}^R,\hfill \end{array}$$ (123) where the functions $`I_{ab}^R`$ are given in (99) and (101). The DL contribution is contained in the first term of $`I_{1a}^R`$ in (99). Extracting the large logarithm $`\mathrm{ln}(Q^2/k_x)`$ which is embodied into the DL function $`r(\overline{\mu },Q^2)`$ in (41), we calculate the geometry-dependent SL correction factors denoted there by $`F_{ab}`$. The $`k_y`$ integrals are convergent so we have set the upper limit $`K_{ym}\mathrm{}`$. Introducing the ratio of momenta $`t=k_y/\left\|k_x\right\|`$ we obtain the following expressions: $$\begin{array}{cc}\hfill F_{1a}(\tau _a)& =2\mathrm{ln}\frac{\alpha _m\left\|\tau _a\right\|}{Q}+\mathrm{\hspace{0.25em}2}_{\mathrm{}}^{\mathrm{}}\frac{dt}{\pi (1+t^2)}\{\frac{1}{2}\mathrm{ln}(1+\kappa _a^{}_{}{}^{}2)\hfill \\ & +t(\frac{\pi }{2}sgn(\tau _a)\mathrm{arctan}(\kappa _a^{}))\},\hfill \\ \hfill F_{23}(\tau _2,\tau _3)& =2\mathrm{ln}\frac{\tau _2}{\left\|\tau _3\right\|}+\mathrm{\hspace{0.25em}2}_{\mathrm{}}^{\mathrm{}}\frac{dt}{\pi (\rho ^2+t^2)}\{\frac{\rho }{2}\mathrm{ln}\left(\frac{1+\kappa _3^{}_{}{}^{}2}{1+\kappa _2^{}_{}{}^{}2}\right)\hfill \\ & t(\mathrm{arctan}(\kappa _2^{})+\mathrm{arctan}(\kappa _3^{}))\}.\hfill \end{array}$$ (124) where $`\kappa _a^{}\frac{\kappa _a}{\|k_x\|}=\sqrt{1+t^2}\frac{\tau _a}{Q}t`$. We remind the reader that convergence of the integrals is assured by $`\tau _2>1`$, $`\tau _3<1`$. ### B.5 Running coupling The two-loop analysis entails that the argument of the running coupling for the dipole distribution $`r_{ab}`$ in (90) is given by $`k_{t,ab}`$, the invariant gluon transverse momentum with respect to the $`ab`$-dipole (see(24)). We show here that, to SL accuracy, we can the scale at the value $`2\|k_x\|`$ instead, as has been stated in (102). This effective value is obtained after integrating $`\alpha _\text{s}(k_{t,ab})`$ over $`\alpha `$ and $`k_y`$. To show this we first observe that we need to control the precise argument of $`\alpha _\text{s}`$ only in the DL contributions which originate from the phase space regions where the gluon is collinear to one of the three hard partons $`P_a`$. Consider first the case of the contribution of $`r_{1a}`$ from the right hemisphere. Here the soft gluon is close to the thrust axis, $`P_1`$, so that the invariant transverse momentum reduces to the usual 2-dimensional momentum, $`k_{t,1a}^2k_t^2=k_x^2+k_y^2`$, which is $`\alpha `$-independent. This allows us to perform the $`\alpha `$-integration and obtain $`I_{1a}^R`$ in (99). Now, to determine the effective scale of $`\alpha _\text{s}`$ it suffices to consider $`k_y`$-integral of the leading piece of $`I_{1a}^R`$ proportional to $`\mathrm{ln}(Q^2/k_x^2)`$ and integrate over $`k_y`$. We have an integral of the type $$A=_{K_{ym}}^{K_{ym}}\frac{dk_y}{\pi k_t^2}\alpha _\text{s}(k_t)=\frac{1}{k_x}_{y_m}^{y_m}\frac{dy}{\pi (1+y^2)}\alpha _\text{s}(k_x\sqrt{1+y^2}),$$ (125) with $`k_x`$ positive, $`yk_y/k_x`$. For small $`k_x`$ the upper limit of the $`y`$-integral is large, $`y_m=K_{ym}/k_x=T_MQ/2k_x1`$, and can set infinite since the integral converges. Expanding the coupling to the first order, we get $$\begin{array}{cc}\hfill k_xA& =_{\mathrm{}}^{\mathrm{}}\frac{dy}{\pi (1+y^2)}\left(\alpha _\text{s}(k_x)\frac{\beta _0}{4\pi }\alpha _\text{s}^2(k_x)\mathrm{ln}(1+y^2)\right)+𝒪\left(\alpha _\text{s}^3\right)\hfill \\ & \left(\alpha _\text{s}(k_x)\frac{\beta _0}{4\pi }\alpha _\text{s}^2(k_x)\mathrm{ln}4\right)\alpha _\text{s}(2k_x).\hfill \end{array}$$ (126) The regions collinear to $`P_2`$ or $`P_3`$ seems more complicated since here $`k_{t,ab}`$ depend both on $`\alpha `$ and $`k_y`$. However, a similar analysis can be carried out in terms of the Sudakov variables with $`P,\overline{P}`$ aligned with the emitting parton, we obtain the same result for the argument of $`\alpha _\text{s}`$. ### B.6 The sources and the $`k_x`$ integrals Here we prove the substitution rule (60). The general structure of the radiators is $$D=_0^Q\frac{dk_x}{k_x}\mathrm{ln}\frac{Q^{}}{k_x}[\mathrm{\hspace{0.17em}1}u],ue^{\nu k_x}\mathrm{cos}(\nu \beta k_x)e^{i\nu \gamma yk_x},$$ (127) with $`Q^{}=𝒪\left(Q\right)`$ a hard scale. This we write as $$D=_{1/\overline{\mu }}^Q\frac{dk_x}{k_x}\mathrm{ln}\frac{Q^{}}{k_x}+\mathrm{\Delta },\mathrm{\Delta }=\frac{}{ϵ}\left\{_0^{1/\overline{\mu }}\frac{dk_x}{k_x}\left(\frac{k_x}{Q^{}}\right)^ϵ_0^Q\frac{dk_x}{k_x}\left(\frac{k_x}{Q^{}}\right)^ϵu\right\}\|_{ϵ=0},$$ (128) and optimise the choice of $`\overline{\mu }`$ such that $`\mathrm{\Delta }=𝒪\left(1\right)`$, i.e. it does not contain a $`\mathrm{ln}\nu `$-enhancement. Since $`\nu Q1`$, the second integral containing the exponential source function, $`u\mathrm{exp}(\nu k_x)`$, can be safely extended to infinity. Evaluating the expression in the curly brackets up to $`𝒪\left(ϵ\right)`$, we obtain $$\begin{array}{cc}\hfill \left\{\right\}& =ϵ^1\left[\left(\overline{\mu }Q^{}\right)^ϵ\mathrm{\Gamma }(1+ϵ)\left(\nu Q^{}\right)^ϵ\frac{1}{2}\left((1i\gamma yi\beta )^ϵ+(1i\gamma y+i\beta )^ϵ\right)\right]\hfill \\ & =\left(\nu Q^{}\right)^ϵ\left(\mathrm{ln}\frac{\overline{\mu }}{\nu }\gamma _E\frac{1}{2}\mathrm{ln}\left[(1i\gamma y)^2+\beta ^2\right]+𝒪\left(ϵ\right)\right).\hfill \end{array}$$ (129) Taking the $`ϵ`$-derivative in (128) we obtain a large parameter $`\mathrm{ln}(\nu Q^{})`$ accompanied by the factor which we set equal to zero to optimize the choice of $`\overline{\mu }`$: $$\overline{\mu }=\nu e^{\gamma _E}\sqrt{(1i\gamma y)^2+\beta ^2}.$$ (130) This means that, within SL accuracy, the source factor $`[1u]`$ can be substituted by $$[1u]\vartheta (k_x\overline{\mu }^1).$$ (131) Performing this substitution in (120) we get the result reported in the text, see (63). ## Appendix C Evaluation of $``$ The expression for $`_\delta `$ is rather complicated. Invoking (54) we split the $`\gamma `$-integral into two pieces, $$=_r+_i,$$ (132) namely the principal value and the $`\delta (\gamma )`$ contributions, $$\frac{1}{\gamma iϵ}=\frac{\text{P}}{\gamma }\pm i\pi \delta (\gamma ).$$ (133) We have $$\begin{array}{c}\hfill =_{\mathrm{}}^{\mathrm{}}\frac{d\beta _2}{\pi (1+\beta _2^2)^{1+{\scriptscriptstyle \frac{1}{2}}C_2^{(\delta )}r^{}}}_{\mathrm{}}^{\mathrm{}}\frac{d\beta _3}{\pi (1+\beta _3^2)^{1+{\scriptscriptstyle \frac{1}{2}}C_3^{(\delta )}r^{}}}_{\mathrm{}}^{\mathrm{}}\frac{d\beta _1}{\pi }(_r+_i),\end{array}$$ (134) where the two integrands are $$\begin{array}{cc}& _i=\frac{1}{2}\left(\frac{1}{1+\beta _{12}^2}+\frac{1}{1+\beta _{13}^2}\right)\left(\sqrt{1+\beta _{12}^2}\sqrt{1+\beta _{13}^2}\right)^{\frac{C_1^{(\delta )}r^{}}{2}},\hfill \\ & _r=(\frac{1}{1+\beta _{12}^2}\frac{1}{1+\beta _{13}^2})_0^{\mathrm{}}\frac{d\gamma }{\pi }(\sqrt{(1+\gamma )^2+\beta _{12}^2}\sqrt{(1+\gamma )^2+\beta _{13}^2})^{\frac{C_1^{(\delta )}r^{}}{2}}\frac{\mathrm{sin}(C_1^{(\delta )}r^{}A_1)}{\gamma }.\hfill \end{array}$$ (135) Here $$A_1=\frac{1}{2\pi }_0^{\mathrm{}}\frac{dx}{1x^2}\left[\mathrm{ln}\frac{(1+\gamma x)^2+\beta _{12}^2}{(1+\gamma )^2+\beta _{12}^2}\mathrm{ln}\frac{(1+\gamma x)^2+\beta _{13}^2}{(1+\gamma )^2+\beta _{13}^2}\right]$$ (136) Where as before the colour factors are given by (122). ### C.1 $``$ in the first order In the $`𝒪\left(r^{}\right)`$ approximation we have, for $`_r`$ $$\begin{array}{c}\hfill _r(\overline{\nu })C_1^{(\delta )}r^{}_{\mathrm{}}^{\mathrm{}}\frac{d\beta _2}{\pi (1+\beta _2^2)}_{\mathrm{}}^{\mathrm{}}\frac{d\beta _3}{\pi (1+\beta _3^2)}_{\mathrm{}}^{\mathrm{}}\frac{d\beta _1}{\pi }\left(\frac{1}{1+\beta _{12}^2}\frac{1}{1+\beta _{13}^2}\right)\frac{1}{\pi }_0^{\mathrm{}}\frac{d\gamma }{\gamma }A_1.\end{array}$$ (137) Using $$\begin{array}{c}\hfill \frac{1}{\pi }_0^{\mathrm{}}\frac{d\gamma }{\gamma }A_1=\frac{1}{8}\mathrm{ln}\frac{1+\beta _{12}^2}{1+\beta _{13}^2},\end{array}$$ (138) we find $$\begin{array}{cc}\hfill _r(\overline{\nu })& \frac{r^{}C_1^{(\delta )}}{2}_{\mathrm{}}^{\mathrm{}}\frac{d\beta _2}{\pi (1+\beta _2^2)}_{\mathrm{}}^{\mathrm{}}\frac{d\beta _3}{\pi (1+\beta _3^2)}_{\mathrm{}}^{\mathrm{}}\frac{d\beta _1}{2\pi }\left(\frac{1}{1+\beta _{12}^2}\frac{1}{1+\beta _{13}^2}\right)\frac{1}{2}\mathrm{ln}\frac{1+\beta _{12}^2}{1+\beta _{13}^2}.\hfill \end{array}$$ (139) Using $$_{\mathrm{}}^{\mathrm{}}\frac{d\beta }{\pi }\frac{\mathrm{ln}(1+\beta ^2)}{1+\beta ^2}=\mathrm{ln}4.$$ (140) and $$\begin{array}{cc}& _{\mathrm{}}^{\mathrm{}}\frac{d\beta _{13}}{\pi }\mathrm{ln}(1+\beta _{13}^2)_{\mathrm{}}^{\mathrm{}}\frac{d\beta _{12}}{\pi (1+\beta _{12}^2)}_{\mathrm{}}^{\mathrm{}}\frac{d\beta _1}{2\pi }\frac{1}{1+(\beta _{12}\beta _1)^2}\frac{1}{1+(\beta _{13}\beta _1)^2}=\mathrm{ln}4.\hfill \end{array}$$ (141) we obtain $$_r(\overline{\nu })\frac{1}{2}\mathrm{ln}4\frac{C_1^{(\delta )}}{2}r^{}(\overline{\nu }).$$ (142) Expanding $`_i`$ in $`r^{}`$ we obtain in the first order: $$\begin{array}{cc}\hfill _i(\overline{\nu })& =1\frac{1}{2}r^{}_{\mathrm{}}^{\mathrm{}}\frac{d\beta _2}{\pi (1+\beta _2^2)}_{\mathrm{}}^{\mathrm{}}\frac{d\beta _3}{\pi (1+\beta _3^2)}_{\mathrm{}}^{\mathrm{}}\frac{d\beta _1}{2\pi }\left(\frac{1}{1+\beta _{12}^2}+\frac{1}{1+\beta _{13}^2}\right)\hfill \\ & \left(C_2^{(\delta )}\mathrm{ln}(1+\beta _2^2)+C_3^{(\delta )}\mathrm{ln}(1+\beta _3^2)+\frac{1}{2}C_1^{(\delta )}\left[\mathrm{ln}(1+\beta _{12}^2)+\mathrm{ln}(1+\beta _{13}^2)\right]\right)\hfill \end{array}$$ (143) Using (140) and (141) we get: $$_i(\overline{\nu })=1\frac{1}{2}\mathrm{ln}4\left(\frac{3}{2}C_1^{(\delta )}+C_2^{(\delta )}+C_1^{(\delta )}\right)r^{}(\overline{\nu }).$$ (144) So, in the first order, $$=_r(\overline{\nu })+_i(\overline{\nu })=1\frac{1}{2}\mathrm{ln}4\left(2C_1^{(\delta )}+C_2^{(\delta )}+C_3^{(\delta )}\right)r^{}(\overline{\nu })$$ (145)
warning/0004/quant-ph0004009.html	ar5iv	text	# Reversible combination of inequivalent kinds of multipartite entanglement. ## Abstract We present a family of tri-partite entangled states that, in an asymptotical sense, can be reversibly converted into EPR states shared by only two of parties (say B and C), and tripartite GHZ states. Thus we show that bipartite and genuine tripartite entanglement can be reversibly combined in several copies of a single tripartite state. For such states the corresponding fractions of GHZ and of EPR states represent a complete quantification of their (asymptotical) entanglement resources. More generally, we show that the three different kinds of bipartite entanglement (AB, AC and BC EPR states) and tripartite GHZ entanglement can be reversibly combined in a single state of three parties. Finally, we generalize this result to any number of parties. Understanding the inequivalent ways in which the parts of a composite system can be entangled is one of the central open questions of quantum information theory. When the system consists only of two parts, A and B, and it has been prepared in a pure state $`\|\psi _{AB}`$, then its entanglement properties are qualitatively equivalent to those of an EPR state , $$\frac{1}{\sqrt{2}}(\|00+\|11),$$ (1) in the following sense . If two parties, Alice and Bob, share $`N`$ copies of the system in state $`\|\psi _{AB}`$ and are allowed to perform local operations assisted with classical communication (LOCC), then they can transform, reversibly in the large $`N`$ limit, the state of their systems into $`NE(\psi _{AB})`$ copies of an EPR state (1), where $`E(\psi _{AB})`$ is the entropy of entanglement of state $`\|\psi _{AB}`$ —namely the von Neumann entropy of the reduced density matrix describing either part $`A`$ or $`B`$—. Thus, reversibility of asymptotical conversions justifies that we regard all bipartite pure-state entanglement as equivalent and quantify it by means of $`E(\psi _{AB})`$. It has been recently shown that a GHZ state $$\frac{1}{\sqrt{2}}(\|000+\|111),$$ (2) of a tripartite system is inequivalent, even in this asymptotic sense, to EPR states distributed among the parties. This indicates that there is genuine tripartite entanglement. Similarly (see also ), for any number $`n`$ of parties sharing a system, the $`n`$-partite GHZ state $$\frac{1}{\sqrt{2}}(\|0^n+\|1^n),$$ (3) can not be reversibly converted by means of LOCC into any distribution of entangled states each one involving less than $`n`$ parties. Here we will refer to entanglement of the form (3) as canonical. Thus, in the general case of a $`m`$-partite system, one can find at least $`2^mm1`$ kinds of entanglement that are asymptotically inequivalent. They correspond, for each $`n=2,\mathrm{},m`$, to all $`m!/(n!(mn)!)`$ kinds of $`n`$-canonical entanglement, that is involving different subsets of $`n`$ parties . In this Letter we show that inequivalent kinds of multipartite entanglement can be reversibly combined into a pure state by means of LOCC. More specifically, we show that all kinds of $`n`$-canonical entanglement, $`n=2,\mathrm{},m`$, can be combined in a $`m`$-partite state, and then again reextracted, with asymptotically vanishing losses. For instance, we will prove that $`N`$ copies of some tripartite state $$\|\psi c_0\|000+c_1\|1\frac{1}{\sqrt{2}}(\|11+\|22),$$ (4) can be reversibly transformed, in the limit $`N\mathrm{}`$, into copies of an EPR state shared by Bob and Claire, and copies of a GHZ state, that is, $$\|\psi ^N\|\text{EPR}_{BC}^{Nl_{BC}}\|\text{GHZ}^{Nl_{ABC}}.$$ (5) This means that the asymptotical entanglement properties of $`\|\psi `$ can be completely characterized by simply specifying the values of $`l_{BC}`$ and $`l_{ABC}`$. Some other states $`\|\psi ^{}`$ of three parties will be reversibly converted into the three inequivalent kinds of bipartite EPR states and GHZ states, so that their entanglement can be characterized by the multicomponent measure $`L_\psi ^{}(l_{AB},l_{AC},l_{BC};l_{ABC})`$. For an arbitrary number $`m`$ of parties, a $`(2^mm1)`$–component measure will also quantify the entanglement properties of some other states $`\|\psi ^{\prime \prime }`$, by providing the amount of all inequivalent kinds of $`n`$-canonical entanglement ($`n=2,\mathrm{},m`$) that can be reversibly extracted from them. While it is not yet clear how many asymptotically inequivalent kinds of entanglement exist, not even whether there is only a finite number of them for the simplest non-trivial case —i.e. for tripartite systems—, our results arguably help in the ongoing effort to understand and classify multipartite quantum correlations, as they show for the first time that it is possible to quantify the entanglement of a given state by relating it to several inequivalent forms of entanglement. We start by analyzing the asymptotic properties of the tripartite state $`\|\psi \mathrm{I}\mathrm{C}^2\mathrm{I}\mathrm{C}^3\mathrm{I}\mathrm{C}^3`$ given by equation (4), where $`c_0,c_10`$, $`c_0^2+c_1^2=1`$. First we will show how the parties can extract, from $`N`$ copies of it and by means of LOCC, up to $`Nc_1^2+g_2(N)`$ EPR pairs between Bob and Claire and $`NS(c_0^2,c_1^2)+g_1(N)`$ GHZ states, where $`S(\{x_i\})=_ix_i\mathrm{log}_2x_i`$ and both $`g_1(x)/x`$ and $`g_2(x)/x`$ vanish as $`x\mathrm{}`$. Then we will prove that the same amounts of canonical entanglement —up to corrections that vanish as $`N\mathrm{}`$— suffice to create the state $`\|\psi ^N`$ (actually a state as faithful to $`\|\psi ^N`$ as wished if $`N`$ can be made arbitrarily large). Therefore we will have $`l_{BC}=c_1^2`$ and $`l_{ABC}=S(c_0^2,c_1^2)`$. Finally, we will then describe generalizations of this result to all possible kinds of canonical entanglement and to an arbitrary number of parties. Let us consider, then, two copies of $`\|\psi `$, which after some convenient relabeling of the local orthonormal basis in $`^{N=2}=\mathrm{I}\mathrm{C}^4\mathrm{I}\mathrm{C}^9\mathrm{I}\mathrm{C}^9`$ can be written as $`\|\psi ^2=c_1^2`$ $`\|1^0{\displaystyle \frac{1}{\sqrt{4}}}(\|1_1^01_1^0+\|2_1^02_1^0+\|3_1^03_1^0+\|4_1^04_1^0)`$ (6) $`+c_0c_1`$ $`\|1^1{\displaystyle \frac{1}{\sqrt{2}}}(\|1_1^11_1^1+\|2_1^12_1^1)`$ (7) $`+c_0c_1`$ $`\|2^1{\displaystyle \frac{1}{\sqrt{2}}}(\|1_2^11_2^1+\|2_2^12_2^1)`$ (8) $`+c_0^2`$ $`\|1^2\|1_1^21_1^2.`$ (9) By means of a local measurement the parties can project the state (9) into one of the three subspaces characterized by a constant coefficient $`c_0^kc_1^{2k}`$ ($`k=0,1,2`$). We point out the relevant fact that either Alice, Bob or Claire could perform such a measurement locally because $`\|\psi `$ is a linear combination $`c_0\|\varphi _1+c_1\|\varphi _2`$ of two locally orthogonal states $`\|\varphi _i`$, i.e. $$\text{Tr}[\rho _i^\alpha \rho _j^\alpha ]=0ij,\alpha =A,B,C,$$ (10) where $`\rho _i^\alpha `$ is the reduced density matrix of $`\|\varphi _i`$ for subsystem $`\alpha `$, and this implies that the three subspaces of (9) are also locally orthogonal; in other words, the parties can manipulate locally each of these subspaces independently. If the result of the measurement corresponds to $`k=0`$, then Bob and Claire will be sharing a $`2^2`$-level maximally entangled state, that is $`2`$ EPR pairs; if the outcome corresponds to $`k=1`$ then the parties end up sharing an EPR<sub>BC</sub> state and a GHZ state, as can be straightforwardly checked by expanding $`\|EPR_{BC}\|GHZ`$ as a linear combination of product states; finally, an outcome related to the subspace $`k=2`$ leaves the parties with a product state $`\|000`$. This structure of outcomes easily generalizes to the case of $`N`$ copies. Let us call block $`(N,k)`$, denoted by $`\|B_{N,k}`$, the normalized projection of $`\|\psi ^N`$ into the subspace characterized by the coefficient $`c_0^kc_1^{Nk}`$, that is $$\|\psi ^N=\underset{k=0}{\overset{N}{}}c_0^kc_1^{Nk}\sqrt{b_{N,k}}\|B_{N,k},$$ (11) $`b_{N,k}N!/[k!(Nk)!]`$. A direct computation shows that $`\|B_{N,k}`$ is of the form, $`\|B_{N,k}={\displaystyle \frac{1}{\sqrt{rt}}}`$ $`\{`$ (12) $`\|1^k(\|1_1^k1_1^k+\|2_1^k2_1^k`$ $`+`$ $`\mathrm{}+\|r_1^kr_1^k)`$ (13) $`+\|2^k(\|1_2^k1_2^k+\|2_2^k2_2^k`$ $`+`$ $`\mathrm{}+\|r_2^kr_2^k)`$ (14) $`\mathrm{}`$ (15) $`+\|t^k(\|1_t^k1_t^k+\|2_t^k2_t^k`$ $`+`$ $`\mathrm{}+\|r_t^kr_t^k)\},`$ (16) where $`r2^{Nk}`$, $`tb_{N,k}`$ and the local states satisfy $`i_a^k\|j_b^k^{}=\delta _{i,j}\delta _{k,k^{}}\delta _{a,b}`$ in both Bob and Claire and $`i^k\|j^k^{}=\delta _{i,j}\delta _{k,k^{}}`$ in Alice. Notice that (16) is equivalent to the tensor product of a $`r`$-level EPR state and a $`t`$-level GHZ state, $`\|B_{N,k}=({\displaystyle \frac{1}{\sqrt{r}}}{\displaystyle \underset{i=1}{\overset{r}{}}}\|ii_{BC})({\displaystyle \frac{1}{\sqrt{t}}}{\displaystyle \underset{i=1}{\overset{t}{}}}\|iii)`$ (17) $`=\|rEPR_{BC}\|tGHZ.`$ (18) Thus, by means of a local measurement projecting onto these blocks, the parties will obtain state (18) with probability $`P_{N,k}b_{N,k}c_0^{2k}c_1^{2(Nk)}`$. The expectation values $`<{\displaystyle \frac{\mathrm{log}_2r}{N}}>`$ $`=`$ $`{\displaystyle \underset{k=0}{\overset{N}{}}}P_{N,k}{\displaystyle \frac{Nk}{N}},`$ (19) $`<{\displaystyle \frac{\mathrm{log}_2t}{N}}>`$ $`=`$ $`{\displaystyle \underset{k=0}{\overset{N}{}}}P_{N,k}{\displaystyle \frac{\mathrm{log}_2b_{N,k}}{N}}`$ (20) correspond, respectively, to the fraction $`l_{BC}`$ of EPR<sub>BC</sub> states and to the fraction $`l_{ABC}`$ of GHZ states that are obtained, on average and per copy of $`\|\psi `$, from $`\|\psi ^N`$ . Because of the smooth behavior of the functions $`(Nk)/N`$ and $`(\mathrm{log}_2b_{N,k})/N`$ compared to the binomial distribution $`P_{N,k}`$, we can calculate the expectation values (20) —up to corrections that vanish in the limit $`N\mathrm{}`$— by just evaluating the two functions at the peak of $`P_{N,k}`$, namely at $`k_{\mathrm{max}}Nc_0^2`$, which gives us the announced amount of entanglement for each of the two canonical forms. Let us look now at the inverse transformation. We start, for clearness sake, by showing that $`2`$ EPR<sub>BC</sub> and $`2`$ GHZ states suffice to create state (9) locally and with certainty. Let us expand $`\|EPR_{BC}^2\|GHZ^2`$ as $`{\displaystyle \frac{1}{\sqrt{4}}}\{`$ $`\|1^0{\displaystyle \frac{1}{\sqrt{4}}}(\|1_1^01_1^0+\|2_1^02_1^0+\|3_1^03_1^0+\|4_1^04_1^0)`$ (21) $`+`$ $`\|1^1{\displaystyle \frac{1}{\sqrt{4}}}(\|1_1^11_1^1+\|2_1^12_1^1+\|3_1^13_1^1+\|4_1^14_1^1)`$ (22) $`+`$ $`\|2^1{\displaystyle \frac{1}{\sqrt{4}}}(\|1_2^11_2^1+\|2_2^12_2^1+\|3_2^13_2^1+\|4_2^14_2^1)`$ (23) $`+`$ $`\|1^2{\displaystyle \frac{1}{\sqrt{4}}}(\|1_1^21_1^2+\|2_1^22_1^2+\|3_1^23_1^2+\|4_1^24_1^2)\}.`$ (24) In this expression (cf. (9)) we would like to give the first row a weight $`c_1^2`$; in both the second and third rows we should get rid of $`\|3_i^23_i^2`$ and $`\|4_i^24_i^2`$ and give each of the rows a weight $`c_0c_1`$; the fourth row should be reduced to a product state with weight $`c_0^2`$. After these changes we will have state (9). We first note that the parties can transform, with certainty, the $`2`$ GHZ states into a triorthogonal state with arbitrary relative weights, $$\frac{1}{2}\underset{i=1}{\overset{4}{}}\|iii\underset{i=1}{\overset{4}{}}\lambda _i\|iii,$$ (25) by means of a local POVM $`\{O_j\}`$, $`j=1,\mathrm{}4`$, on (any) one of the parties followed by an outcome dependent, local unitary $`U_j`$ applied once in each of the parties’ subsystems. Here $`O_j`$ is proportional to $`_i\lambda _i\|i_4ji_4j\|`$ and $`U_j`$ takes $`\|i_4j`$ into $`\|i`$ on each local subsystem. The tensor product of $`2`$ EPR<sub>BC</sub> states with the resulting state in (25) is equivalent to (24) but with row $`i`$ having weight $`\lambda _i`$. Bob and Claire can now address each of the $`4`$ rows locally and reduce their length at wish. Indeed, in order to shorten the fourth row into a product vector, one of them, say Bob, can perform a POVM with $`4`$ positive operators $$Q_i=\|i_1^2i_1^2\|+\frac{1}{2}\underset{j=2,3,4}{}P_j,i=1,\mathrm{}4,$$ (26) where $`P_j`$ is a projector onto the local subspace supporting row $`j`$, e.g. $`P_2=_i\|i_1^1i_1^1\|`$; then Bob and Claire need to relabel the surviving term $`\|i_1^2i_1^2`$ of the first row as $`\|1_1^21_1^2`$. By means of similar POVMs followed by outcome dependent local unitaries Bob and Claire can tailor also the second and third row so that they contain only $`2`$ product terms each. Explicitly, a $`2`$-outcome POVM that reduces the second row reads $`Q_1^{}=\|1_1^11_1^11_1^11_1^1\|+\|2_1^12_1^12_1^12_1^1\|+{\displaystyle \frac{1}{\sqrt{2}}}{\displaystyle \underset{j=1,3,4}{}}P_j,`$ (27) $`Q_2^{}=\|3_1^13_1^13_1^13_1^1\|+\|4_1^14_1^14_1^14_1^1\|+{\displaystyle \frac{1}{\sqrt{2}}}{\displaystyle \underset{j=1,3,4}{}}P_j.`$ (28) Notice that such measurements do not modify the relative weight of the rows. A proper choice of the coefficients $`\lambda _i`$ in the first step of the transformation, namely $`\lambda _1=c_1^2,\lambda _2=\lambda _3=c_1c_0`$ and $`\lambda _4=c_0^2`$, completes therefore the protocol for preparing $`2`$ copies of $`\|\psi `$. In the case of a large number $`N`$ of copies the parties start with several EPR<sub>BC</sub> and GHZ states and want to create a state $$\|N_k_{}^{k_+}K\underset{k=k_{}}{\overset{k_+}{}}c_0^{Nk}c_1^k\sqrt{b_{N,k}}\|B_{N,k}$$ (29) such that $`F\|N_k_{}^{k_+}\|\psi ^N\|^2=1h(N)`$, where $`h(x\mathrm{})=0`$, that is a state which asymptotically can not be distinguished from $`\|\psi ^N`$. First we note that an arbitrary faithfulness can be achieved, asymptotically, by considering only the blocks $`\|B_{N,k}`$ (cf. eq. (16)) that correspond to $`k`$’s around $`k_{\mathrm{max}}=Nc_0^2`$. Indeed, a straightforward computation of the fidelity shows that it suffices to take $`k_\pm =c_0^2N\pm \alpha N^\beta `$ for some $`\alpha >0`$ and $`1/2<\beta <1`$: using that a binomial distribution is asymptotically equivalent to a normal (Gaussian) distribution, the fidelity $`F`$ can be seen to be bounded from below by $`\mathrm{\Phi }(2\alpha N^{\beta 1/2})`$, where $`\mathrm{\Phi }(x)1/\sqrt{2\pi }_x^xe^{y^2/2}𝑑y`$. For our choice of $`\alpha ,\beta `$, we have that $`F1`$ when $`N\mathrm{}`$. As with the $`N=2`$ case, our plan is, starting from a reasonable amount of EPR<sub>BC</sub> and GHZ states —which can be expanded in the fashion of (24)—, ($`i`$) to modify conveniently the weight of each row in the expansion and ($`ii`$) to shorten the length of each row, to obtain the pattern of lengths given by the block structure of (29). Notice that a straightforward generalization of (25) provides row $`i`$ with an arbitrary weight $`\lambda _i`$ by locally manipulating the initial $`GHZ`$ states, and that we have also already seen how to shorten each row independently by means of a local POVM in either Bob’s or Claire’s side (see POVMs (26) and (28) as examples). Thus, the only question left concerns the amount of canonical entanglement required to produce all blocks $`\|B_{N,k}`$ in (29). Since we have a mechanism to shorten but not to lengthen the rows, the number of EPR<sub>BC</sub> states must allow to obtain the longest rows, which are those of the block $`\|N,k_{}`$ and contain $`2^{Nk_{}}`$ product terms each. That is, $`Nk_{}`$ EPR<sub>BC</sub> states will suffice. The total number of GHZ states required is the logarithm of the total number of rows in (29), and thus reads $`\mathrm{log}_2(_{k=k_{}}^{k_+}b_{N,k})`$. Let $`k_0[k_{},k_+]`$ be the value that maximizes $`b_{N,k}`$ in this interval. Then the amount of GHZ states required to prepared (29) is bounded from above by $`\mathrm{log}_2[(k_+k_{}+1)b_{N,k_o}]`$. Finally, substitution of $`k_\pm `$ and $`k_0`$ in this bound and the previous estimation for EPR<sub>BC</sub> states shows that both amounts of canonical entanglement needed to prepare (29) are the expected ones, apart from corrections which scale sublinearly in $`N`$ and that therefore become irrelevant for $`N`$ sufficiently large. This concludes the proof that state (4) is asymptotically equivalent to canonical entanglement. We can now generalize the previous example and reversibly combine the three kinds of bipartite entanglement and the canonical tripartite entanglement in a single state. Indeed, the tripartite state $`\|\psi ^{}`$ $`c_0\|000`$ (30) $`+`$ $`c_1\|1{\displaystyle \frac{1}{\sqrt{2}}}(\|11+\|22)`$ (31) $`+`$ $`c_2{\displaystyle \frac{1}{\sqrt{2}}}(\|233+\|334)`$ (32) $`+`$ $`c_3{\displaystyle \frac{1}{\sqrt{2}}}(\|44+\|55)\|5,`$ (33) can be transformed by means of LOCC, into EPR states shared by $`2`$ parties and into GHZ states, the asymptotic ratios being, $`l_{AB}=c_1^2`$, $`l_{AC}=c_2^2`$ and $`l_{BC}=c_3^2`$ for the bipartite entanglement and $`l_{ABC}=S(\{c_i^2\})`$ for the tripartite entanglement. This result follows from considering analogous transformations to the ones described above. The expansion of the state of $`N`$ copies of $`\|\psi ^{}`$ into locally orthogonal subspaces as in (11) depends now on $`3`$ independent indices, the weights defining the block structure being $`c_0^{k_0}c_1^{k_1}c_2^{k_2}c_3^{Nk_1k_2k_3}`$. The binomial probability distribution is replaced by a multinomial distribution —centered at $`k_i=Nc_i^2`$— and each one of the new blocks $`\|B_{N,k_0,k_1,k_2}`$ is equivalent to the tensor product of some number of GHZ, EPR<sub>AB</sub>, EPR<sub>AC</sub> and EPR<sub>BC</sub> states. A local measurement onto the blocks leads, for sufficiently large $`N`$, to the desired expectation values for the fractions of canonical entanglement. Conversely, these amounts of entanglement suffice to create all the relevant blocks $`\|B_{N,k_0,k_1,k_2}`$. This is done by introducing some weights $`\lambda _i`$ in the initial GHZ states and by locally tailoring the (now multidimensional) rows in the pertinent expansion, as we explained in the previous example. More generally, let the $`m`$-partite state $$\|\psi ^{\prime \prime }\underset{i=0}{\overset{l}{}}c_i\|\varphi _i$$ (34) be a linear combination of locally orthogonal states \[see equation (10)\] such that each one is the tensor product of a canonical state $`\|\tau _i`$ (3) for $`n_i`$ of the parties, and a product vector for the remaining $`mn_i`$ parties . Then $`N`$ copies of the state $`\|\psi ^{\prime \prime }`$ are asymptotically equivalent to $`Nc_0^2`$ copies of $`\|\tau _0`$, $`\mathrm{}`$, $`Nc_l^2`$ copies of $`\|\tau _l`$ and to $`NS(\{c_i^2\})`$ copies of a $`m`$-canonical state, i.e. $$\|\psi ^{\prime \prime }^N[\underset{i=0}{\overset{l}{}}\|\tau _i^{Nc_i^2}](\|0^m+\|1^m)^{NS(\{c_i^2\})}.$$ (35) In this Letter we have provided examples of multipartite states whose entanglement properties can be classified and quantified in relation to the set of canonical states (3). The criteria underlying such classification is the asymptotical equivalence of states under LOCC. We have shown that at least in some cases, namely for states of the form (34), this criteria brings a significant simplification in the general problem of classifying entanglement. Indeed, our results show that the states (34), which depend on the set of continuous non-local parameters $`\{c_i\}`$, contain only a finite set of inequivalent forms of entanglement. We have gone further and quantified the amount of each form of entanglement contained in state (34), which gives rise to a multicomponent measure. We do not know to what extend the coefficients, as well as the reference states of this measure are essentially unique. This work was supported by the Austrian Science Foundation under the SFB “control and measurement of coherent quantum systems” (Project 11), the European Community under the TMR network ERB–FMRX–CT96–0087, the European Science Foundation and the Institute for Quantum Information GmbH. G.V also acknowledges a Marie Curie Fellowship HPMF-CT-1999-00200 (European Community).
warning/0004/math0004153.html	ar5iv	text	# A functional expression for the curvature of hyper-dimensional Riemannian spaces ## 1. Introduction The well-known Riemannian mathematical model of defining the curvature of hyper-dimensional curvilinear metric spaces, via Gaussian concept of the two-dimensional surface curvature, ; and , is an imperfect in spite of that it functional related to internal geometry of a space: 1. Firstly, for a reason that the concept of the curvature of hyper-dimensional metric spaces reduced to the concept of Gaussian curvature of two-dimensional geodesic surface, it cannot be said that Riemannian concept of the curvature of hyper-dimensional spaces is general one because it is an inapplicable to one-dimensional space. 2. As secondly, since there are more than one geodesic surface in vector spaces of higher dimensions, it is obvious that Riemannian concept of the curvature of hyper-dimensional spaces, in the general case, is not uniquely defined one. 3. Finally, at any an individual, concrete case, from the practical point of view, it is not simple to come to the quantitatively usable functional expression for Riemannian curvature of an analyzed hyper-dimensional metric space. Hence, the one other mathematical model of defining the curvature of the hyper-dimensional metric spaces, which essentially differs from Riemannian model, is presented in this paper. Concretely, the mathematical model is being discussed, which can be said to be generalization of well-known model of defining the curvature of curve and surface in the differential geometry. ### 1.1. Basic characteristics of space continuum The concept of geometrical point is one of fundamental concepts. Closely related to the concept of geometrical points is the system of values $`(a_1,a_2,\mathrm{}a_N)`$ of some an arbitrary $`N`$ variables $`(x_1,x_2,\mathrm{}x_N)`$ such that a set of all geometrical points, and for all real values of the variables, is a real $`N`$ \- dimensional space of a space continuum, . The geometrical point $`𝐎`$, defined by system of zero values $`(0,0,\mathrm{}\mathrm{\hspace{0.17em}0})`$, is an origin of system of reference (of co-ordinate system) of the space. The vector $`\stackrel{}{r}\left(x^i\right)`$, defined with respect to the origin $`𝐎`$, is a position vector. Note that the concept of a vector, in the vector hyper-dimensional spaces $`\left(N>3\right)`$, should be conditionally comprehend in the sense of its geometrical presentation in a form of segments, hence it bears a name linear tensor, . Covariant vectors $`\stackrel{}{e}_i`$: $`\stackrel{}{e}_i=_{x_i}\stackrel{}{r}\left(x^j\right)`$, where $`_{x_i}=\frac{}{x^i}`$, form the covariant vector basis $`\left\{\stackrel{}{e}_i\right\}_{i=1}^N`$ of the tangent space of a space continuum. The vectors $`\stackrel{}{e}^i`$, such that at each point of a space: $`\stackrel{}{e}_i\stackrel{}{e}^k=\delta _i^k`$, where the second order system of the unit values $`\delta _i^k`$ \- is an unit $`N\times N`$ matrix (Kronecker’s delta-symbol), ; and , form a basis $`\left\{\stackrel{}{e}^i\right\}_{i=1}^N`$, which is called the dual basis of the covariant vector basis $`\left\{\stackrel{}{e}_i\right\}_{i=1}^N`$. This is so-called natural isomorphism from $`\left\{\stackrel{}{e}_i\right\}_{i=1}^N`$ onto $`\left\{\stackrel{}{e}^i\right\}_{i=1}^N`$. The infinitesimal value $`d\stackrel{}{r}`$ of the position vector $`\stackrel{}{r}`$ of representative point is defined by $`d\stackrel{}{r}=dx^i\stackrel{}{e}_i=dx_i\stackrel{}{e}^i`$, where the well-known Einstein’s convention is applied to summation with respect to the repetitive indexes (uppers and lowers), and , herein as well as in the further text of the paper. ## 2. The main results ### 2.1. A curvature of hyper-dimensional Riemannian spaces By the following transformation low: $`x^i=x^i\left(q^\alpha \right)`$; $`i=1,2,\mathrm{},N`$, $`\alpha =1,2,\mathrm{},MN`$, in the general case, an arbitrary $`M`$ \- dimensional metric space is defined, $$ds^2=_{x^i}\stackrel{}{r}_{q^\alpha }x^i_{x^j}\stackrel{}{r}_{q^\beta }x^jdq^\alpha dq^\beta =e_{ij}_{q^\alpha }x^i_{q^\beta }x^jdq^\alpha dq^\beta =g_{\alpha \beta }dq^\alpha dq^\beta ,$$ where the positive definite symetric square matrices $`e_{ij}`$: $`e_{ij}=\stackrel{}{e}_i\stackrel{}{e}_j`$ and $`g_{\alpha \beta }`$: $$g_{\alpha \beta }=_{q^\alpha }\stackrel{}{r}_{q^\beta }\stackrel{}{r}=\stackrel{}{g}_\alpha \stackrel{}{g}_\beta =_{x^i}\stackrel{}{r}_{q^\alpha }x^i_{x^j}\stackrel{}{r}_{q^\beta }x^j=\stackrel{}{e}_i\stackrel{}{e}_j_{q^\alpha }x^i_{q^\beta }x^j,$$ of degree: $`N`$ and $`M`$, respectively, are fundamental (metric) tensors of an ambient $`N`$ \- dimensional Euclidean (more exactly Cartesian) space $`x^i`$, as well as, in the general case of Riemannian covering map ($`M<N`$), of an internal $`M`$ \- dimensional Riemannian curvilinear space $`q^\alpha `$ immersed in it. As it is well-known if $`M=N`$ a smooth map $`x^iq^\alpha `$ is called an isometric immersion. The smallest possible dimensional difference $`C`$: $`C=NM`$, is said to define a class of the Riemannian vector space $`q^\alpha `$, . Hence, Riemannian space $`q^\alpha `$: $`q^\alpha =q`$, of class $`C`$: $`C=N1`$, is an arbitrary curve of the ambient $`N`$ \- dimensional Euclidean (Cartesian) space $`x^i`$. Since the vector $`d_q\stackrel{}{g}`$, where $`d_q=\frac{d}{dq}`$, as derivative of the fundamental vector $`\stackrel{}{g}`$: $`\stackrel{}{g}=d_q\stackrel{}{r}`$, lying in the tangent space of Riemannian space $`q`$ of class $`C`$: $`C=N1`$ (on the tangent of curve), along the curve, from the point of view of the interior of an ambient space $`x^i`$, is a vector of the ambient space $`x^i`$, then at each points of the curve there exist $`C`$: $`C=N1`$, the linearly independent and mutually orthogonal unit vectors: $`\stackrel{}{n}^\mathrm{\Lambda }`$; $`\mathrm{\Lambda }=1,2,\mathrm{},N1`$, being orthogonal on the unit vector of tangent $`\stackrel{}{t}`$: $`\stackrel{}{t}=d_s\stackrel{}{r}`$, of the curve, such that $$d_{qq}^2\stackrel{}{r}=d_q\left(d_qx^i\stackrel{}{e}_i\right)=_{x^j}\stackrel{}{e}_id_qx^jd_qx^i+d_{qq}^2x^i\stackrel{}{e}_i=$$ (2.1) $$=d_q\stackrel{}{g}=\left(d_q\stackrel{}{g}\stackrel{}{n}^\mathrm{\Lambda }\right)\stackrel{}{n}_\mathrm{\Lambda }+\left(d_q\stackrel{}{g}\stackrel{}{g}\right)\stackrel{}{g}.$$ The covariant vectors $`\stackrel{}{n}_\mathrm{\Sigma }`$, as well as the vectors $`\stackrel{}{n}^\mathrm{\Lambda }`$ satisfying the condition: $`\stackrel{}{n}_\mathrm{\Sigma }\stackrel{}{n}^\mathrm{\Lambda }=\delta _\mathrm{\Sigma }^\mathrm{\Lambda }`$, form covariant $`\left\{\stackrel{}{n}_\mathrm{\Sigma }\right\}_{\mathrm{\Sigma }=1}^{N1}`$, as well as dual vector basis $`\left\{\stackrel{}{n}^\mathrm{\Lambda }\right\}_{\mathrm{\Lambda }=1}^{N1}`$ of a normal vector space of Riemannian space $`q`$ of class $`C`$: $`C=N1`$, respectively. From the point of view of the internal geometry of Riemannian space $`q`$, the vector $`\stackrel{}{\kappa }`$: $$\stackrel{}{\kappa }=\left(d_q\stackrel{}{g}\stackrel{}{n}^\mathrm{\Lambda }\right)\stackrel{}{n}_\mathrm{\Lambda },$$ is a vector of curvature of curve, more exactly, from the point of view of the internal geometry of an ambient Euclidean space $`x^i`$, it is a vector of geodesic curvature of curve, . Namely, by projection of the vector $`\stackrel{}{\kappa }`$ onto the normal vector space of Riemannian space $`q`$ and onto the tangent vector space of ambient Euclidean space $`x^i`$, it is obtained that $$\stackrel{}{\kappa }\stackrel{}{n}^\mathrm{\Sigma }=d_q\stackrel{}{g}\stackrel{}{n}^\mathrm{\Sigma }=\left(_{x^j}\stackrel{}{e}_i\stackrel{}{e}^kd_qx^id_qx^j+d_{qq}^2x^k\right)\left(\stackrel{}{e}_k\stackrel{}{n}^\mathrm{\Sigma }\right)=$$ $$=\left(\mathrm{\Gamma }_{ij}^kd_qx^id_qx^j+d_{qq}^2x^k\right)n_k^\mathrm{\Sigma }$$ and $$\stackrel{}{\kappa }\stackrel{}{e}^l=\left(d_q\stackrel{}{g}\stackrel{}{n}^\mathrm{\Lambda }\right)\stackrel{}{n}_\mathrm{\Lambda }\stackrel{}{e}^l=\left(\mathrm{\Gamma }_{ij}^kd_qx^id_qx^j+d_{qq}^2x^k\right)\stackrel{}{e}_k\stackrel{}{n}^\mathrm{\Lambda }\left(\stackrel{}{n}_\mathrm{\Lambda }\stackrel{}{e}^l\right),$$ where the mixed systems $`\mathrm{\Gamma }_{ij}^k`$: $`\mathrm{\Gamma }_{ij}^k=_{x^j}\stackrel{}{e}_i\stackrel{}{e}^k`$, are coefficients of connection (Christoffel’s symbols) of the second type, and an intensity of vector $$\frac{1}{\stackrel{}{g}\stackrel{}{g}}\left(\mathrm{\Gamma }_{ij}^kd_qx^id_qx^j+d_{qq}^2x^k\right)\stackrel{}{e}_k=\left(d_s\stackrel{}{t}\stackrel{}{e}^k\right)\stackrel{}{e}_k=d_s\stackrel{}{t},$$ being in the direction of the vector $`\stackrel{}{\kappa }`$: $$\frac{1}{\stackrel{}{g}\stackrel{}{g}}\left(\mathrm{\Gamma }_{ij}^kd_qx^id_qx^j+d_{qq}^2x^k\right)\left(\stackrel{}{e}_k\stackrel{}{n}^\mathrm{\Lambda }\right)\stackrel{}{n}_\mathrm{\Lambda }=\frac{1}{\stackrel{}{g}\stackrel{}{g}}\left(\stackrel{}{\kappa }\stackrel{}{e}^l\right)\stackrel{}{e}_l=\frac{1}{\stackrel{}{g}\stackrel{}{g}}\stackrel{}{\kappa },$$ defines both the curvature $`\kappa `$ of the curve, from the point of view of the internal geometry of Riemannian space $`q`$ of class $`C`$: $`C=N1`$ and geodesic curvature $`\kappa _g`$ of the curve, from the point of view of the internal geometry of ambient Euclidean space $`x^i`$, and . Analogously to the vector of curvature: $`\stackrel{}{\kappa }=\left(d_q\stackrel{}{g}\stackrel{}{n}^\mathrm{\Lambda }\right)\stackrel{}{n}_\mathrm{\Lambda }`$, of a Riemannian space $`q`$ of class $`C`$: $`C=N1`$, in the general case of Riemannian space $`q^\alpha `$ of class $`C`$: $`C=NM`$, immersed in $`N`$ \- dimensional ambient Euclidean space $`x^i`$, from the viewpoint of its internal geometry, the matrix scheme of vectors $$\stackrel{}{K}_{\alpha \beta }=\left(_{q^\beta }\stackrel{}{g}_\alpha \stackrel{}{n}^\mathrm{\Lambda }\right)\stackrel{}{n}_\mathrm{\Lambda }=\tau _{\alpha \beta }^\mathrm{\Lambda }\stackrel{}{n}_\mathrm{\Lambda },$$ (2.2) is a matrix scheme of the sectional curvature vectors, of a Riemannian space $`q^\alpha `$ of class $`C`$: $`C=NM`$. The determinant $`\left\|\tau _{\alpha \beta }^\mathrm{\Lambda }\tau _{\gamma \delta }^\mathrm{\Sigma }g^{\delta \beta }n_{\mathrm{\Lambda }\mathrm{\Sigma }}\right\|`$ of the matrix obtained by the matrix multiplication (rows $`\times `$ columns) of the matrix schemes of curvature vectors: $`\tau _{\alpha \beta }^\mathrm{\Lambda }\stackrel{}{n}_\mathrm{\Lambda }`$ and $`\tau _{\gamma \delta }^\mathrm{\Sigma }g^{\delta \beta }\stackrel{}{n}_\mathrm{\Sigma }`$, is the proportionate to the square of the curvature $`\kappa `$ of Riemannian space $`q^\alpha `$ of class $`C`$: $`C=NM`$, and with the coefficient of proportionality which is equal to the determinant of matrix of fundamental tensor $`g_{\alpha \beta }`$ $$\left\|\stackrel{}{K}_{\alpha \beta }\stackrel{}{K}_\gamma ^\beta \right\|=\left\|\tau _{\alpha \beta }^\mathrm{\Lambda }\tau _{\gamma \delta }^\mathrm{\Sigma }g^{\delta \beta }n_{\mathrm{\Lambda }\mathrm{\Sigma }}\right\|=\left\|g_{\alpha \beta }\right\|\kappa ^2,$$ more exactly, $$\kappa ^2=\frac{\left\|\tau _{\alpha \beta }^\mathrm{\Lambda }\tau _{\gamma \delta }^\mathrm{\Sigma }g^{\delta \beta }n_{\mathrm{\Lambda }\mathrm{\Sigma }}\right\|}{\left\|g_{\alpha \beta }\right\|}.$$ (2.3) If $`N`$ \- dimensional ambient space $`x^i`$ is either Riemannian curvilinear space of class $`C`$ or Euclidean curvilinear space of class $`C`$: $`C1`$, both with normal vector space $`\stackrel{}{w}_P`$; $`P=1,2,\mathrm{},C`$, then the matrix scheme of vectors $$\stackrel{}{K}_{\alpha \beta }=\left(_{q^\beta }\stackrel{}{g}_\alpha \stackrel{}{n}^\mathrm{\Lambda }\right)\stackrel{}{n}_\mathrm{\Lambda }+\left(_{q^\beta }\stackrel{}{g}_\alpha \stackrel{}{w}^P\right)\stackrel{}{w}_P,$$ is a matrix scheme of vectors of curvature of Riemannian space $`q^\alpha `$ of class $`\widehat{C}`$: $`\widehat{C}=C+NM`$. From the point of view of the internal geometry of the ambient space $`x^i`$, the partial matrix schemes of vectors of curvature: $`\stackrel{}{N}_{\alpha \beta }=\left(_{q^\beta }\stackrel{}{g}_\alpha \stackrel{}{w}^P\right)\stackrel{}{w}_P=\varkappa _{\alpha \beta }^P\stackrel{}{w}_P`$ and $`\stackrel{}{G}_{\alpha \beta }=\left(_{q^\beta }\stackrel{}{g}_\alpha \stackrel{}{n}^\mathrm{\Lambda }\right)\stackrel{}{n}_\mathrm{\Lambda }=\tau _{\alpha \beta }^\mathrm{\Lambda }\stackrel{}{n}_\mathrm{\Lambda }`$, are matrix schemes of vectors of normal curvature, as well as of geodesic curvature, of Riemannian space $`q^\alpha `$ of class $`\widehat{C}`$, respectively. The square of intensity of curvature $`\kappa `$ of Riemannian space of class $`\widehat{C}`$, in this case is equal to the sum of square of the normal curvature $`\kappa _n`$: $$\kappa _n=\frac{\left\|\varkappa _{\alpha \beta }^P\varkappa _{\gamma \delta }^Lg^{\beta \delta }w_{PL}\right\|}{\left\|g^{\alpha \beta }\right\|}$$ and of the geodesic curvature $`\kappa _g`$: $$\kappa _g=\frac{\left\|\tau _{\alpha \beta }^\mathrm{\Lambda }\tau _{\gamma \delta }^\mathrm{\Sigma }g^{\beta \delta }n_{\mathrm{\Lambda }\mathrm{\Sigma }}\right\|}{\left\|g^{\alpha \beta }\right\|},$$ of Riemannian space $`q^\alpha `$ of class $`\widehat{C}`$: $$\kappa ^2=\kappa _n^2+\kappa _g^2.$$ The space $`q^\alpha `$ of class $`\widehat{C}`$, as well as of the geodesic curvature zero $`\kappa _g`$: $`\kappa _g=0`$, with respect to the ambient space $`x^i`$ of class $`C`$, is a geodesic sub-space of that ambient space.$`\mathrm{}`$ #### 2.1.1. Functional expression for the curvature of hyper-dimensional Riemannian spaces, from the point of view of the internal geometry of space. If the vector functional equality (2.1) $$_{q^\beta }\stackrel{}{g}_\alpha =\left(_{q^\beta }\stackrel{}{g}_\alpha \stackrel{}{g}^\gamma \right)\stackrel{}{g}_\gamma +\left(_{q^\beta }\stackrel{}{g}_\alpha \stackrel{}{n}^\mathrm{\Lambda }\right)\stackrel{}{n}_\mathrm{\Lambda },$$ in the general case of the ambient Riemannian space $`q^\alpha `$ of class $`C`$: $`C=NM`$, is partially differentiated with respect to dual co-ordinates $`q^\delta `$, it is obtained that $$_{q^\beta q^\delta }^2\stackrel{}{g}_\alpha =_{q^\delta }\mathrm{\Gamma }_{\alpha \beta }^\gamma \stackrel{}{g}_\gamma +\mathrm{\Gamma }_{\alpha \beta }^\gamma _{q^\delta }\stackrel{}{g}_\gamma +_{q^\delta }\tau _{\alpha \beta }^\mathrm{\Lambda }\stackrel{}{n}_\mathrm{\Lambda }+\tau _{\alpha \beta }^\mathrm{\Lambda }_{q^\delta }\stackrel{}{n}_\mathrm{\Lambda },$$ (2.4) that is, $$_{q^\delta q^\beta }^2\stackrel{}{g}_\alpha =_{q^\beta }\mathrm{\Gamma }_{\alpha \delta }^\gamma \stackrel{}{g}_\gamma +\mathrm{\Gamma }_{\alpha \delta }^\gamma _{q^\beta }\stackrel{}{g}_\gamma +_{q^\beta }\tau _{\alpha \delta }^\mathrm{\Lambda }\stackrel{}{n}_\mathrm{\Lambda }+\tau _{\alpha \delta }^\mathrm{\Lambda }_{q^\beta }\stackrel{}{n}_\mathrm{\Lambda }.$$ (2.5) Since differentials $`d\stackrel{}{g}_\alpha =_{q^\beta }\stackrel{}{g}_\alpha dq^\beta `$ are an absolute ones, more exactly, the covariant vector basis $`\left\{\stackrel{}{g}_\alpha \right\}`$ of a Riemannian space $`q^\alpha `$ of class $`C`$: $`C=NM`$, is uniquely defined at each point of a space, and in accordance with that the condition of integrability: $`_{q^\beta q^\delta }^2\stackrel{}{g}_\alpha _{q^\delta q^\beta }^2\stackrel{}{g}_\alpha =0`$, is satisfied, then by projection of difference of the respectable vector functional expressions on both sides of the previous equations: (2.4) and (2.5), onto covariant vector basis $`\left\{\stackrel{}{g}_\rho \right\}`$, the functional relation is obtained $$_{q^\delta }\mathrm{\Gamma }_{\alpha \beta }^\gamma g_{\gamma \rho }_{q^\beta }\mathrm{\Gamma }_{\alpha \delta }^\gamma g_{\gamma \rho }+\mathrm{\Gamma }_{\alpha \beta }^\lambda \mathrm{\Gamma }_{\lambda \delta }^\gamma g_{\gamma \rho }\mathrm{\Gamma }_{\alpha \delta }^\lambda \mathrm{\Gamma }_{\lambda \beta }^\gamma g_{\gamma \rho }=\tau _{\alpha \beta }^\mathrm{\Lambda }\tau _{\delta \rho ,\mathrm{\Lambda }}\tau _{\alpha \delta }^\mathrm{\Lambda }\tau _{\beta \rho ,\mathrm{\Lambda }},$$ having in mind both the equation (2.2) and orthogonality of vectors: $`\stackrel{}{g}_\alpha `$ and $`\stackrel{}{n}_\mathrm{\Lambda }`$. On account of the fact that the functional expression on the left hand side of preceding relation represents Riemann-Christoffel’s tensor of Riemannian space $`q^\alpha `$ of class $`C`$: $`C=NM`$, ; and , it finally follows that $$R_{\rho \alpha \delta \beta }=\tau _{\alpha \beta }^\mathrm{\Lambda }\tau _{\delta \rho ,\mathrm{\Lambda }}\tau _{\alpha \delta }^\mathrm{\Lambda }\tau _{\beta \rho ,\mathrm{\Lambda }}.$$ (2.6) By application of Gauss-Chiò’s procedure for condensation of determinants, , and on the basis of tensorial relation (2.6), as well as of functional relation (2.3), it is possible the curvature of Riemannian space $`q^\alpha `$ of class $`C`$: $`C=NM`$, to come into functional relation to internal geometry of a space, more exactly to the components of Rimann-Christoffel’s tensor of the curvature $`R_{\rho \alpha \delta \beta }`$. Namely, if in addition to the matrix scheme of vectors: $`\stackrel{}{K}_{\alpha \beta }=\tau _{\alpha \beta }^\mathrm{\Lambda }\stackrel{}{n}_\mathrm{\Lambda }`$, the matrix schemes of vectors: $$\stackrel{}{\mathrm{\Phi }}_{\alpha \beta }=\left[\begin{array}{ccccc}\varphi ^\mathrm{\Lambda }\stackrel{}{n}_\mathrm{\Lambda }\hfill & 0\hfill & \mathrm{}\hfill & 0\hfill & 0\hfill \\ \tau _{12}^\mathrm{\Lambda }\stackrel{}{n}_\mathrm{\Lambda }\hfill & \tau _{11}^\mathrm{\Lambda }\stackrel{}{n}_\mathrm{\Lambda }\hfill & \mathrm{}\hfill & 0\hfill & 0\hfill \\ \mathrm{}\hfill & \mathrm{}\hfill & \mathrm{}\hfill & \mathrm{}\hfill & \mathrm{}\hfill \\ \tau _{1\left(M1\right)}^\mathrm{\Lambda }\stackrel{}{n}_\mathrm{\Lambda }\hfill & 0\hfill & \mathrm{}\hfill & \tau _{11}^\mathrm{\Lambda }\stackrel{}{n}_\mathrm{\Lambda }\hfill & 0\hfill \\ \tau _{1M}^\mathrm{\Lambda }\stackrel{}{n}_\mathrm{\Lambda }\hfill & 0\hfill & \mathrm{}\hfill & 0\hfill & \tau _{11}^\mathrm{\Lambda }\stackrel{}{n}_\mathrm{\Lambda }\hfill \end{array}\right];$$ $$\stackrel{}{Z}_{\alpha \beta }=\left[\begin{array}{ccccc}\tau _{12}^\mathrm{\Lambda }\stackrel{}{n}_\mathrm{\Lambda }\hfill & 0\hfill & \mathrm{}\hfill & 0\hfill & 0\hfill \\ \tau _{11}^\mathrm{\Lambda }\stackrel{}{n}_\mathrm{\Lambda }\hfill & \zeta ^\mathrm{\Lambda }\stackrel{}{n}_\mathrm{\Lambda }\hfill & \mathrm{}\hfill & \tau _{1\left(M1\right)}^\mathrm{\Lambda }\stackrel{}{n}_\mathrm{\Lambda }\hfill & \tau _{1M}^\mathrm{\Lambda }\stackrel{}{n}_\mathrm{\Lambda }\hfill \\ \mathrm{}\hfill & \mathrm{}\hfill & \mathrm{}\hfill & \mathrm{}\hfill & \mathrm{}\hfill \\ 0\hfill & 0\hfill & \mathrm{}\hfill & \tau _{12}^\mathrm{\Lambda }\stackrel{}{n}_\mathrm{\Lambda }\hfill & 0\hfill \\ 0\hfill & 0\hfill & \mathrm{}\hfill & 0\hfill & \tau _{12}^\mathrm{\Lambda }\stackrel{}{n}_\mathrm{\Lambda }\hfill \end{array}\right];$$ $$\stackrel{}{\mathrm{\Xi }}_{\alpha \beta }=\left[\begin{array}{ccccc}\xi ^\mathrm{\Lambda }\stackrel{}{n}_\mathrm{\Lambda }\hfill & 0\hfill & \mathrm{}\hfill & 0\hfill & 0\hfill \\ \tau _{22}^\mathrm{\Lambda }\stackrel{}{n}_\mathrm{\Lambda }\hfill & \tau _{21}^\mathrm{\Lambda }\stackrel{}{n}_\mathrm{\Lambda }\hfill & \mathrm{}\hfill & 0\hfill & 0\hfill \\ \mathrm{}\hfill & \mathrm{}\hfill & \mathrm{}\hfill & \mathrm{}\hfill & \mathrm{}\hfill \\ \tau _{2\left(M1\right)}^\mathrm{\Lambda }\stackrel{}{n}_\mathrm{\Lambda }\hfill & 0\hfill & \mathrm{}\hfill & \tau _{21}^\mathrm{\Lambda }\stackrel{}{n}_\mathrm{\Lambda }\hfill & 0\hfill \\ \tau _{2M}^\mathrm{\Lambda }\stackrel{}{n}_\mathrm{\Lambda }\hfill & 0\hfill & \mathrm{}\hfill & 0\hfill & \tau _{21}^\mathrm{\Lambda }\stackrel{}{n}_\mathrm{\Lambda }\hfill \end{array}\right];$$ $$\stackrel{}{\mathrm{\Psi }}_{\alpha \beta }=\left[\begin{array}{ccccc}\tau _{22}^\mathrm{\Lambda }\stackrel{}{n}_\mathrm{\Lambda }\hfill & 0\hfill & \mathrm{}\hfill & 0\hfill & 0\hfill \\ \tau _{21}^\mathrm{\Lambda }\stackrel{}{n}_\mathrm{\Lambda }\hfill & \psi ^\mathrm{\Lambda }\stackrel{}{n}_\mathrm{\Lambda }\hfill & \mathrm{}\hfill & \tau _{2\left(M1\right)}^\mathrm{\Lambda }\stackrel{}{n}_\mathrm{\Lambda }\hfill & \tau _{2M}^\mathrm{\Lambda }\stackrel{}{n}_\mathrm{\Lambda }\hfill \\ \mathrm{}\hfill & \mathrm{}\hfill & \mathrm{}\hfill & \mathrm{}\hfill & \mathrm{}\hfill \\ 0\hfill & 0\hfill & \mathrm{}\hfill & \tau _{22}^\mathrm{\Lambda }\stackrel{}{n}_\mathrm{\Lambda }\hfill & 0\hfill \\ 0\hfill & 0\hfill & \mathrm{}\hfill & 0\hfill & \tau _{22}^\mathrm{\Lambda }\stackrel{}{n}_\mathrm{\Lambda }\hfill \end{array}\right],$$ are also introduced into analysis, where: $`\varphi ^\mathrm{\Lambda }\stackrel{}{n}_\mathrm{\Lambda }`$;$`\zeta ^\mathrm{\Lambda }\stackrel{}{n}_\mathrm{\Lambda }`$;$`\xi ^\mathrm{\Lambda }\stackrel{}{n}_\mathrm{\Lambda }`$ and $`\psi ^\mathrm{\Lambda }\stackrel{}{n}_\mathrm{\Lambda }`$, are arbitrary vector functions of the normal vector space $`\stackrel{}{n}_\mathrm{\Lambda }`$ of Riemannian space $`q^\alpha `$ of class $`C`$: $`C=NM`$, then having in view the fact that determinant of the product of matrices, as it is well-known, is equal to the product of determinants of any matrix separately, and , it follows that $$(\stackrel{}{\mathrm{\Phi }}_{\gamma \eta }\stackrel{}{K}_\delta ^\eta )(\stackrel{}{K}_{\alpha \lambda }\stackrel{}{\mathrm{\Psi }}_\beta ^\lambda )=\frac{\left(\tau _{11}^\mathrm{\Lambda }\varphi _\mathrm{\Lambda }\right)\left(\tau _{22}^\mathrm{\Lambda }\psi _\mathrm{\Lambda }\right)}{\left\|g_{\alpha \beta }\right\|^2}\times $$ (2.7) $$\times \left\|\begin{array}{cccc}R_{1212}\hfill & R_{1213}\hfill & \mathrm{}\hfill & R_{121M}\hfill \\ R_{1312}\hfill & R_{1313}\hfill & \mathrm{}\hfill & R_{131M}\hfill \\ \mathrm{}\hfill & \mathrm{}\hfill & \mathrm{}\hfill & \mathrm{}\hfill \\ R_{1M12}\hfill & R_{1M13}\hfill & \mathrm{}\hfill & R_{1M1M}\hfill \end{array}\right\|\left\|\begin{array}{cccc}R_{1212}\hfill & R_{1232}\hfill & \mathrm{}\hfill & R_{12M2}\hfill \\ R_{3212}\hfill & R_{3232}\hfill & \mathrm{}\hfill & R_{32M2}\hfill \\ \mathrm{}\hfill & \mathrm{}\hfill & \mathrm{}\hfill & \mathrm{}\hfill \\ R_{M212}\hfill & R_{M232}\hfill & \mathrm{}\hfill & R_{M2M2}\hfill \end{array}\right\|=$$ $$=\frac{\left(\tau _{11}^\mathrm{\Lambda }\varphi _\mathrm{\Lambda }\right)\left(\tau _{22}^\mathrm{\Lambda }\psi _\mathrm{\Lambda }\right)}{\left\|g_{\alpha \beta }\right\|^2}\mathrm{\Delta }_{\tau _{11}^\mathrm{\Lambda }}^R\mathrm{\Delta }_{\tau _{22}^\mathrm{\Lambda }}^R,$$ more exactly<sup>1</sup><sup>1</sup>1 $$\left\|\stackrel{}{\mathrm{\Phi }}_{\alpha \delta }\stackrel{}{\mathrm{\Psi }}_\beta ^\delta \right\|=\frac{1}{\left\|g_{\alpha \beta }\right\|}\left\|\begin{array}{cccc}\varphi ^\mathrm{\Lambda }\tau _{22,\mathrm{\Lambda }}\hfill & 0\hfill & \mathrm{}\hfill & 0\hfill \\ \left(\tau _{22}^\mathrm{\Lambda }\tau _{12,\mathrm{\Lambda }}+\tau _{11}^\mathrm{\Lambda }\tau _{21,\mathrm{\Lambda }}\right)\hfill & \psi ^\mathrm{\Lambda }\tau _{11,\mathrm{\Lambda }}\hfill & \mathrm{}\hfill & \tau _{11}^\mathrm{\Lambda }\tau _{2M,\mathrm{\Lambda }}\hfill \\ \mathrm{}\hfill & \mathrm{}\hfill & \mathrm{}\hfill & \mathrm{}\hfill \\ \tau _{22}^\mathrm{\Lambda }\tau _{1M,\mathrm{\Lambda }}\hfill & 0\hfill & \mathrm{}\hfill & \tau _{11}^\mathrm{\Lambda }\tau _{22,\mathrm{\Lambda }}\hfill \end{array}\right\|=$$ $$=\frac{1}{\left\|g_{\alpha \beta }\right\|}\left(\varphi ^\mathrm{\Lambda }\tau _{22,\mathrm{\Lambda }}\right)\left(\psi ^\mathrm{\Lambda }\tau _{11,\mathrm{\Lambda }}\right)\left(\tau _{11}^\mathrm{\Lambda }\tau _{22,\mathrm{\Lambda }}\right)^{M2}.$$ $$\|\stackrel{}{K}_{\alpha \delta }\stackrel{}{K}_\beta ^\delta \|\left(\tau _{11}^\mathrm{\Lambda }\tau _{22,\mathrm{\Lambda }}\right)^{M2}=\frac{1}{\left\|g_{\alpha \beta }\right\|}\left\|\begin{array}{cccc}R_{1212}\hfill & R_{1213}\hfill & \mathrm{}\hfill & R_{121M}\hfill \\ R_{1312}\hfill & R_{1313}\hfill & \mathrm{}\hfill & R_{131M}\hfill \\ \mathrm{}\hfill & \mathrm{}\hfill & \mathrm{}\hfill & \mathrm{}\hfill \\ R_{1M12}\hfill & R_{1M13}\hfill & \mathrm{}\hfill & R_{1M1M}\hfill \end{array}\right\|\times $$ $$\times \left\|\begin{array}{cccc}R_{1212}\hfill & R_{1232}\hfill & \mathrm{}\hfill & R_{12M2}\hfill \\ R_{3212}\hfill & R_{3232}\hfill & \mathrm{}\hfill & R_{32M2}\hfill \\ \mathrm{}\hfill & \mathrm{}\hfill & \mathrm{}\hfill & \mathrm{}\hfill \\ R_{M212}\hfill & R_{M232}\hfill & \mathrm{}\hfill & R_{M2M2}\hfill \end{array}\right\|=\frac{1}{\left\|g_{\alpha \beta }\right\|}\mathrm{\Delta }_{\tau _{11}^\mathrm{\Lambda }}^R\mathrm{\Delta }_{\tau _{22}^\mathrm{\Lambda }}^R,$$ considering the fact that the functional expression in the relation (2.7), which is an independent of the choice of arbitrary vector functions: $`\varphi ^\mathrm{\Lambda }\stackrel{}{n}_\mathrm{\Lambda }`$ and $`\psi ^\mathrm{\Lambda }\stackrel{}{n}_\mathrm{\Lambda }`$ (the specially interesting case is one in which $`\varphi ^\mathrm{\Lambda }=\psi ^\mathrm{\Lambda }=\tau _{12}^\mathrm{\Lambda }`$), defines the determinant $`\left\|\stackrel{}{K}_{\alpha \delta }\stackrel{}{K}_\beta ^\delta \right\|`$. Similarly $$\|\stackrel{}{K}_{\alpha \delta }\stackrel{}{K}_\beta ^\delta \|\left(\tau _{12}^\mathrm{\Lambda }\tau _{21,\mathrm{\Lambda }}\right)^{M2}=\frac{1}{\left\|g_{\alpha \beta }\right\|}\left\|\begin{array}{cccc}R_{1212}\hfill & R_{1232}\hfill & \mathrm{}\hfill & R_{12M2}\hfill \\ R_{1213}\hfill & R_{1233}\hfill & \mathrm{}\hfill & R_{12M3}\hfill \\ \mathrm{}\hfill & \mathrm{}\hfill & \mathrm{}\hfill & \mathrm{}\hfill \\ R_{121M}\hfill & R_{123M}\hfill & \mathrm{}\hfill & R_{12MM}\hfill \end{array}\right\|\times $$ $$\times \left\|\begin{array}{cccc}R_{2121}\hfill & R_{2123}\hfill & \mathrm{}\hfill & R_{212M}\hfill \\ R_{2131}\hfill & R_{2133}\hfill & \mathrm{}\hfill & R_{213M}\hfill \\ \mathrm{}\hfill & \mathrm{}\hfill & \mathrm{}\hfill & \mathrm{}\hfill \\ R_{21M1}\hfill & R_{21M3}\hfill & \mathrm{}\hfill & R_{21MM}\hfill \end{array}\right\|=\frac{1}{\left\|g_{\alpha \beta }\right\|}\mathrm{\Delta }_{\tau _{12}^\mathrm{\Lambda }}^R\mathrm{\Delta }_{\tau _{21}^\mathrm{\Lambda }}^R.$$ On the basis of the functional formulation of the curvature of Riemannian space $`q^\alpha `$ of class $`C`$: $`C=NM`$, the relation (2.3): $`\kappa ^2=\frac{\left\|\stackrel{}{K}_{\alpha \delta }\stackrel{}{K}_\beta ^\delta \right\|}{\left\|g_{\alpha \beta }\right\|}`$, as well as of the two previously derived relations, and for $`\left(M>2\right)`$, it finally follows that $$\kappa ^2=\frac{1}{\left\|g_{\alpha \beta }\right\|^2\left(R_{1212}\right)^{M2}}\left[\left(\mathrm{\Delta }_{\tau _{11}^\mathrm{\Lambda }}^R\mathrm{\Delta }_{\tau _{22}^\mathrm{\Lambda }}^R\right)^{\frac{1}{M2}}\left(\mathrm{\Delta }_{\tau _{12}^\mathrm{\Lambda }}^R\mathrm{\Delta }_{\tau _{21}^\mathrm{\Lambda }}^R\right)^{\frac{1}{M2}}\right]^{M2}.$$ (2.8) Clearly, in the case of Riemannian space $`q^\alpha `$ of class $`C`$: $`C=NM`$ ($`M=2`$), functional expression (2.3) for the curvature of surface is reduced to the well-known Gaussian curvature of surface in the theory of surfaces, and $$\kappa =\frac{R_{1212}}{\left\|g_{\alpha \beta }\right\|}.$$ In the case when the component of Riemann-Christoffel’s tensor is equal to zero, it is possible, in the functional expression (2.8) for the curvature of Riemannian spaces, to take any another combination of the components of the matrix scheme of the curvature vectors of space: $`\stackrel{}{K}_{\alpha \beta }=\tau _{\alpha \beta }^\mathrm{\Lambda }\stackrel{}{n}_\mathrm{\Lambda }`$, as the support elements of Gauss-Chiò’s procedure for condensation of determinants, and for which some of the components of Riemann-Christoffel’s tensor of curvature are not equal to zero. Clearly, if all components of Riemann-Christoffel’s tensor of the curvature of space are identically equal to zero, then the space is Riemannian space of curvature zero (Euclidean space).$`\mathrm{}`$ ### 2.2. Sub-spaces of curvature of hyper-dimensional spaces Let $`N`$ \- dimensional ambient space of space continuum $`x^i`$, just as in the first Comment of preceding Section 2.1 of this paper, be either Riemannian space of class $`C`$ or Euclidean curvilinear space of class $`C`$: $`C1`$, both with the normal vector space $`\stackrel{}{w}_P`$; $`P=1,2,\mathrm{},C`$. If and only if the vector components: $$\left(_{q^\delta }\stackrel{}{N}_{\alpha \beta }g^{\alpha \delta }\stackrel{}{e}_i\right)\stackrel{}{e}^i=\varkappa _{\alpha \beta }^Pg^{\alpha \delta }\left(_{q^\delta }\stackrel{}{w}_P\stackrel{}{e}_i\right)\stackrel{}{e}^i,$$ of vectors obtained by a partial differentiation of the matrix scheme of vector of normal curvature: $`\stackrel{}{N}_{\alpha \beta }=\left(_{q^\delta }\stackrel{}{g}_\alpha \stackrel{}{w}^P\right)\stackrel{}{w}_P=\varkappa _{\alpha \beta }^P\stackrel{}{w}_P`$, of Riemannian space $`q^\alpha `$ of class $`\widehat{C}`$: $`\widehat{C}=C+NM`$, immersed in ambient space $`x^i`$ of class $`C`$, are vectors of the tangent vector space $`\stackrel{}{g}_\alpha `$ of Riemannian space $`q^\alpha `$ of class $`\widehat{C}`$, in other words if and only if the mutually equivalent conditions: $$\left(_{q^\delta }\stackrel{}{N}_{\alpha \beta }g^{\alpha \delta }\stackrel{}{e}_i\right)\stackrel{}{e}^i=\widehat{\kappa }^2e_{ij}_{q^\beta }x^j\stackrel{}{e}^i=\widehat{\kappa }^2\left(\stackrel{}{g}_\beta \stackrel{}{e}_i\right)\stackrel{}{e}^i=\widehat{\kappa }^2\stackrel{}{g}_\beta $$ (2.9) and $$\varkappa _{\alpha \beta }^Pg^{\alpha \delta }\left(_{q^\delta }\stackrel{}{w}_P\stackrel{}{e}_i\right)=\varkappa _{\alpha \beta }^Pg^{\alpha \delta }t_{ij,P}_{q^\delta }x^j=\widehat{\kappa }^2e_{ij}_{q^\beta }x^j,$$ (2.10) are satisfied<sup>2</sup><sup>2</sup>2The second of the two relations of equality is reduced to: $`t_{k,P}^jt_{ij}^P_{q^\beta }x^k=\widehat{\kappa }^2e_{ij}_{q^\beta }x^j`$, clearly on the condition that all vectors of vector components $`\left(_{q^\delta }\stackrel{}{w}^P\stackrel{}{e}^i\right)\stackrel{}{e}_i`$ are also vectors of the tangent vector space $`\stackrel{}{g}_\alpha `$. Namely, since: $`t_{k,P}^j_{q^\beta }x^k=_{q^\beta }\stackrel{}{w}_P\stackrel{}{e}^i`$ and $`\varkappa _{\alpha \beta ,P}g^{\alpha \delta }_{q^\delta }x^i=\left(_{q^\beta }\stackrel{}{w}_P\stackrel{}{g}_\alpha \right)\stackrel{}{g}^\alpha \stackrel{}{g}^\delta \left(\stackrel{}{e}^i\stackrel{}{g}_\delta \right)`$, then $`\varkappa _{\alpha \beta ,P}g^{\alpha \delta }_{q^\delta }x^i=t_{k,P}^j_{q^\beta }x^k`$., then Riemannian sub-space $`q^\alpha `$ of class $`\widehat{C}`$, is a sub-space of curvature of the ambient space $`x^i`$ of class $`C`$. In the case in which a Riemannian space $`q^\alpha `$ of class $`\widehat{C}`$: $`\widehat{C}=C+NM`$, is an arbitrary curve $`\left(M=1\right)`$, the preceding conditions are reduced to the conditions: $`\left(_q\stackrel{}{N}\frac{1}{\stackrel{}{g}\stackrel{}{g}}\stackrel{}{e}_i\right)\stackrel{}{e}^i=\widehat{\kappa }^2e_{ij}_qx^j\stackrel{}{e}^i`$ and $`\left(d_q\stackrel{}{g}\frac{1}{\stackrel{}{g}\stackrel{}{g}}\stackrel{}{w}^P\right)t_{ij,P}_qx^j\stackrel{}{e}^i=\widehat{\kappa }^2e_{ij}_qx^j\stackrel{}{e}^i`$, as well as $`\left(d_s\stackrel{}{t}\stackrel{}{w}^P\right)t_{ij,P}_sx^j\stackrel{}{e}^i=k_{ij}_sx^j\stackrel{}{e}^i=\widehat{\kappa }^2e_{ij}_sx^j\stackrel{}{e}^i`$. The last of them is the well-known condition for the curve to be a line of the curvature of the ambient space $`x^i`$ of class $`C`$, . It is obvious from this that conditions: (2.9) and (2.10), are a generalization of the preceding conditions. By projection of the condition (2.9), onto covariant vector basis $`\left\{\stackrel{}{g}_\gamma \right\}`$, it is obtained that $$\varkappa _{\alpha \beta }^Pg^{\alpha \delta }t_{ij,P}_{q^\delta }x^j_{q^\gamma }x^i=\varkappa _{\alpha \beta }^P\varkappa _{\gamma ,P}^\alpha =\widehat{\kappa }^2e_{ij}_{q^\beta }x^j_{q^\gamma }x^i=\widehat{\kappa }^2g_{\beta \gamma },$$ more exactly, $$\left(\kappa _n\right)^2=\frac{\left\|\varkappa _{\alpha \beta }^P\varkappa _{\gamma ,P}^\alpha \right\|}{\left\|g_{\beta \gamma }\right\|}=\widehat{\kappa }^{2M},$$ with regard to the relation (2.3). Riemannian space of curvature: $`q^\alpha `$, of class $`\widehat{C}`$, as a sub-space of the ambient space $`x^i`$ of class $`C`$, is said to be the principal, if and only if $$t_{j,P}^kt_{ik}^P_{q^\beta }x^j=\widehat{\kappa }^2e_{ij}_{q^\beta }x^j,$$ (2.11) in other words, if all vectors of the vector components: $`\left(_{q^\delta }\stackrel{}{w}_P\stackrel{}{e}^i\right)\stackrel{}{e}_i`$, are vectors of the tangent vector space (see the Footnote 2 of the paper). In order to exist nontrivial solutions of the previous homogeneous linear system (2.11) of ordinary differential equations with respect to the unknowns $`_{q^\beta }x^j`$, the following condition $$\left\|t_{j,P}^kt_{ik}^P\widehat{\kappa }^2e_{ij}\right\|=0,$$ (2.12) must be satisfied, and . On the basis of the developed form of the preceding condition (2.12), expressed by polynomial of $`N`$-th degree with respect to the unknown $`\widehat{\kappa }^2`$, as well as of Viète’s formulas , and taking the relation (2.3) into consideration, it finally follows that $$k^2=\frac{\left\|t_{j,P}^kt_{ik}^P\right\|}{\left\|e_{ij}\right\|}=\underset{i=1}{\overset{𝑁}{}}\widehat{\kappa }_i^2.$$ If the varied of all eigevalues $`\widehat{\kappa }_i`$ of the matrix: $`t_{j,P}^kt_{ik}^P`$, is an unit, then the principal Riemannian sub-spaces of curvature: $`q^\alpha `$, of class $`\widehat{C}`$, are one-dimensional ones $`\left(M=1\right)`$, in other words there exist $`N`$ principal directions of curvature with the normal curvatures: $`\left(\kappa _n\right)_i=\widehat{\kappa }_i`$, of the ambient space $`x_i`$ of class $`C`$. On the condition that there exists at least one of all eigevalues $`\stackrel{~}{\kappa }`$ with the varied $`M`$, then there exists at least one principal $`M`$ \- dimensional Riemannian sub-space of curvature: $`q^\alpha `$, of class $`\widehat{C}`$, with the normal curvature: $`\kappa _n=\stackrel{~}{\kappa }^M`$, as well as $`NM`$ mutually orthogonal principal directions of curvature with the normal curvatures: $`\left(\kappa _n\right)_i=\widehat{\kappa }_i`$ , such that are orthogonal onto the principal $`M`$ \- dimensional Riemannian sub-spaces of curvature: $`q^\alpha `$, of class $`\widehat{C}`$. In this emphasized case $$k=\stackrel{~}{\kappa }^M\underset{i=1}{\overset{NM}{}}\widehat{\kappa }_i.$$ Furthermore, if the vector components: $$\left(_{q^\delta }\stackrel{}{N}_{\alpha \beta }\stackrel{}{e}_i\right)\stackrel{}{e}^i=\varkappa _{\alpha \beta }^P\left(_{q^\delta }\stackrel{}{w}_P\stackrel{}{e}_i\right)\stackrel{}{e}^i,$$ are also vectors of the tangent vector space $`\stackrel{}{g}_\delta `$ of Riemannian sub-space $`q^\alpha `$ of class $`\widehat{C}`$ $$\left(_{q^\delta }\stackrel{}{N}_{\alpha \beta }\stackrel{}{e}_i\right)\stackrel{}{e}^i=\widehat{\kappa }^2g_{\alpha \beta }e_{ki}_{q^\delta }x^k\stackrel{}{e}^i=\widehat{\kappa }^2g_{\alpha \beta }\left(\stackrel{}{g}_\delta \stackrel{}{e}_i\right)\stackrel{}{e}^i=\widehat{\kappa }^2g_{\alpha \beta }\stackrel{}{g}_\delta ,$$ more exactly, $$\varkappa _{\alpha \beta }^Pt_{ij,P}_{q^\gamma }x^i_{q^\delta }x^j=\varkappa _{\alpha \beta }^P\varkappa _{\gamma \delta ,P}=\widehat{\kappa }^2g_{\alpha \beta }g_{\gamma \delta },$$ then, the normal curvature $`\kappa _n`$ of sub-space of curvature: $`q^\alpha `$, of class $`\widehat{C}`$, of the ambient space $`x^i`$ of class $`C`$, is determined by functional form $$\left(\kappa _n\right)^2=\frac{\left\|\varkappa _{\alpha \beta }^P\varkappa _{\delta ,P}^\alpha \right\|}{\left\|g_{\alpha \beta }\right\|}=\widehat{\kappa }^{2M},$$ more exactly, $${}_{}{}^{M}\sqrt{\left(\kappa _n\right)^2}=\widehat{\kappa }^2=\frac{\left(t_{ij}^Pt_{kl,P}t_{il}^Pt_{kj,P}\right)_{q^\alpha }x^i_{q^\beta }x^j_{q^\gamma }x^l_{q^\delta }x^k}{\left(e_{ij}e_{kl}e_{il}e_{kj}\right)_{q^\alpha }x^i_{q^\beta }x^j_{q^\gamma }x^l_{q^\delta }x^k};$$ $${}_{}{}^{M}\sqrt{\left(\kappa _n\right)^2}=\widehat{\kappa }^2=\frac{R_{ikjl}_{q^\alpha }x^i_{q^\beta }x^j_{q^\gamma }x^l_{q^\delta }x^k}{\left(e_{ij}e_{kl}e_{il}e_{kj}\right)_{q^\alpha }x^i_{q^\beta }x^j_{q^\gamma }x^l_{q^\delta }x^k}=$$ $$=\frac{R_{\alpha \delta \beta \gamma }}{g_{\alpha \beta }g_{\gamma \delta }g_{\alpha \gamma }g_{\beta \delta }}=\frac{R}{M\left(M1\right)},$$ where the scalar invariant $`R`$: $`R=g^{\alpha \gamma }g^{\beta \delta }R_{\alpha \delta \beta \gamma }`$, is so-called invariant of curvature (scalar curvature), . It can be proved by an application of the result of Schur Theorem, , that Riemannian sub-spaces of curvature: $`q^\alpha `$, of class $`\widehat{C}`$, of the ambient space $`x^i`$ of class $`C`$, whether they are principals or not, in this emphasized case are isotropic spaces of the constant normal curvature $`\kappa _n`$.$`\mathrm{}`$ #### 2.2.1. Example The curvature of Riemannian spaces: $$ds^2=e^{\mu \left(\rho \right)}d\rho ^2+\rho ^2d\theta ^2+\mathrm{sin}^2\theta d\phi ^2+e^{\nu \left(\rho \right)}d\tau ^2;$$ $$\mu \left(\rho \right)=\nu \left(\rho \right)=\mathrm{ln}\left(1\frac{2m}{\rho }\right);m=const.$$ and $$d\overline{s}^2=e^{\mu \left(\rho \right)}\left(d\rho ^2+\rho ^2d\theta ^2+\mathrm{sin}^2\theta d\phi ^2+d\tau ^2\right);$$ $$\mu \left(\rho \right)=\frac{2m}{\rho };m=const.,$$ with spherical symmetry. Components of Christoffel’s symbols: $$\mathrm{\Gamma }_{\alpha \beta ,\gamma }=\frac{1}{2}\left(_{q^\alpha }g_{\beta \gamma }+_{q^\beta }g_{\alpha \gamma }_{q^\gamma }g_{\alpha \beta }\right),$$ as well as of Reimann-Christoffel’s tensor of curvature: $$R_{\alpha \beta \delta \gamma }=_{q^\delta }\mathrm{\Gamma }_{\beta \gamma ,\alpha }_{q^\gamma }\mathrm{\Gamma }_{\beta \delta ,\alpha }+g^{\lambda \sigma }\left(\mathrm{\Gamma }_{\beta \delta ,\lambda }\mathrm{\Gamma }_{\alpha \gamma ,\sigma }\mathrm{\Gamma }_{\beta \gamma ,\lambda }\mathrm{\Gamma }_{\alpha \delta ,\sigma }\right),$$ which are not identically zeros, for these sub-classes of general class of Riemannian spaces with spherical symmetry, are the following forms: $$\mathrm{\Gamma }_{11,1}=\overline{\mathrm{\Gamma }}_{11,1}=\frac{1}{2}e^\mu _\rho \mu ;\mathrm{\Gamma }_{22,1}=\mathrm{\Gamma }_{12,2}=\rho ;\overline{\mathrm{\Gamma }}_{22,1}=\overline{\mathrm{\Gamma }}_{12,2}=\frac{\rho e^\mu }{2}\left(2+\rho _\rho \mu \right);$$ $$\mathrm{\Gamma }_{33,1}=\mathrm{\Gamma }_{13,3}=\rho \mathrm{sin}^2\theta ;\overline{\mathrm{\Gamma }}_{33,1}=\overline{\mathrm{\Gamma }}_{13,3}=\frac{\rho \mathrm{sin}^2\theta e^\mu }{2}\left(2+\rho _\rho \mu \right);$$ $$\mathrm{\Gamma }_{33,2}=\mathrm{\Gamma }_{23,3}=\rho \mathrm{sin}\theta \mathrm{cos}\theta ;\overline{\mathrm{\Gamma }}_{33,2}=\overline{\mathrm{\Gamma }}_{23,3}=\rho ^2e^\mu \mathrm{sin}\theta \mathrm{cos}\theta ;$$ $$\mathrm{\Gamma }_{44,1}=\mathrm{\Gamma }_{14,4}=\frac{1}{2}e^\nu _\rho \nu ;\overline{\mathrm{\Gamma }}_{44,1}=\overline{\mathrm{\Gamma }}_{14,4}=\frac{1}{2}e^\mu _\rho \mu ,$$ as well as $$R_{1212}=\frac{1}{2}\rho _\rho \mu ;\overline{R}_{1212}=\frac{1}{2}\rho e^\mu \left(_\rho \mu +_{\rho \rho }^2\mu \right);$$ $$R_{1313}=\frac{1}{2}\rho \mathrm{sin}^2\theta _\rho \mu ;\overline{R}_{1313}=\frac{1}{2}\rho e^\mu \mathrm{sin}^2\theta \left(_\rho \mu +\rho _{\rho \rho }^2\mu \right);$$ $$R_{2323}=\rho ^2\mathrm{sin}^2\theta (1e^\mu );\overline{R}_{2323}=\rho ^3e^\mu _\rho \mu \mathrm{sin}^2\theta \left(1+\frac{\rho }{4}_\rho \mu \right);$$ $$R_{1414}=\frac{1}{2}e^\nu \left\{\frac{1}{2}\left[_\rho \mu _\rho \nu \left(_\rho \nu \right)^2\right]_{\rho \rho }^2\nu \right\};\overline{R}_{1414}=\frac{1}{2}e^\mu _{\rho \rho }^2\mu ;$$ $$R_{2424}=\frac{1}{2}\rho e^{\nu \mu }_\rho \nu ;\overline{R}_{2424}=\frac{1}{2}\rho e^\mu _\rho \mu \left(1+\frac{\rho }{2}_\rho \mu \right).$$ By the functional relation (2.8), it follows that: $$\kappa =\frac{1}{R_{1212}g}\sqrt{\left(R_{1212}\right)^2R_{1313}R_{1414}R_{2323}R_{2424}}=\frac{1}{g}\sqrt{R_{1313}R_{1414}R_{2323}R_{2424}},$$ $$\kappa =\frac{\left\{\frac{1}{8}\rho ^4(1e^\mu )e^{2\nu \mu }\mathrm{sin}^4\theta _\rho \mu _\rho \nu \left[_{\rho \rho }^2\nu +\frac{1}{2}\left(_\rho \nu \right)^2\frac{1}{2}_\rho \mu _\rho \nu \right]\right\}^{\frac{1}{2}}}{\rho ^4e^{\mu +\nu }\mathrm{sin}^2\theta }=2\frac{m^2}{\rho ^6}$$ and $$\overline{\kappa }=\frac{1}{\overline{R}_{1212}\overline{g}}\sqrt{\left(\overline{R}_{1212}\right)^2\overline{R}_{1313}\overline{R}_{1414}\overline{R}_{2323}\overline{R}_{2424}}=\frac{1}{\overline{g}}\sqrt{\overline{R}_{1313}\overline{R}_{1414}\overline{R}_{2323}\overline{R}_{2424}},$$ $$\overline{\kappa }=\frac{\left[\frac{1}{8}\rho ^5e^{4\mu }\mathrm{sin}^4\theta \left(_\rho \mu \right)^2\left(_\rho \mu +\rho _{\rho \rho }^2\mu \right)_{\rho \rho }^2\mu \left(1+\frac{\rho }{4}_\rho \mu \right)\left(1+\frac{\rho }{2}_\rho \mu \right)\right]^{\frac{1}{2}}}{\rho ^4e^{4\mu }\mathrm{sin}^2\theta }=$$ $$=2\frac{m^2}{\rho ^6}e^{\frac{4m}{\rho }}\sqrt{\left(1+\frac{m}{2\rho }\right)\left(1+\frac{m}{\rho }\right)}.\mathrm{}$$ ## 3. CONCLUSION By functional form (2.8), derived from general functional form (2.3) generalizing the concept of the curvature of Riemannian both one and two dimensional spaces to the general concept of the curvature of Riemannian vector spaces $`q^\alpha `$ of class $`C`$: $`C=NM`$, the concept itself of the curvature of Riemannian vector spaces of higher dimensions $`\left(M2\right)`$, directly related to internal geometry of space, more exactly, to components of Reimann-Christoffel’s tensor of curvature $`R_{\alpha \delta \beta \gamma }`$. The process of generalization of fundamental concepts of the differential geometry, presented in this paper, gives the solid base to further generalization other, whether they are fundamentals or not, concepts and theorems of the differential geometry of a surface, and what may be the subject of separated analysis. The one of such concepts is that of Codazzi’s equations (formula) of a surface, and , which can be, in the general case of Riemannian space $`q^\alpha `$ of class $`C`$: $`C=NM`$, obtained by projection of vector condition of integrability: $$_{q^\delta q^\beta }^2\stackrel{}{g}_\alpha _{q^\beta q^\delta }^2\stackrel{}{g}_\alpha =0,$$ of absolute differentials of the fundamental vectors $`\stackrel{}{g}_\alpha `$:$`d\stackrel{}{g}_\alpha =_{q^\beta }\stackrel{}{g}_\alpha dq^\beta `$, onto the normal vector space of Riemannian space $`q^\alpha `$ of class $`C`$.
warning/0004/quant-ph0004033.html	ar5iv	text	# A new conception experimental test of Bell inequalities using non-maximally entangled states ## 0.1 Introduction In 1935 Einstein-Podolsky-Rosen , analysing the measurement on a entangled state, proposed that Quantum Mechanics (QM) could be an incomplete theory, representing a statistical approximation of a complete deterministic theory: this was the born of Local Hidden Variable Theories, where observables values are fixed by some hidden variable and probabilistic predictions become epistemic, being due to our ignorance of hidden variables values. A subsequent fundamental progress in discussing Local Hidden Variable (LHV) was Bell’s discovery that any theory of this kind must satisfy certain inequalities which can be violated in QM leading in principle to a possible experimental test of the validity of the standard interpretation of QM compared to LHV. Since then, many interesting experiments have been devoted to a test of Bell inequalities, the most interesting of them using photon pairs , leading to a substantial agreement with quantum mechanics and disfavouring LHV theories. But, up to now, no experiment has yet been able to exclude definitively such theories. In fact, so far, one has always been forced to introduce, at least, a further additional hypothesis , due to the low total detection efficiency, stating that the observed sample of particle pairs is a faithful subsample of the initial set of pairs. This problem is known as detection or efficiency loophole. The research for new experimental configurations able to overcome the detection loophole is, of course, of the greatest interest. In the 90’s a big progress in the direction of eliminating this loophole has been obtained by using parametric down conversion (PDC) processes. This technique has been largely employed for generating ”entangled” photon pairs, i.e. pairs of photons described by a common wave function which cannot be factorised into the product of two distinct wave functions pertaining to separated photons. The generation of entangled states by parametric down conversion (PDC) has replaced other techniques, such as the radiative decay of excited atomic states, as it was in the celebrated experiment of A. Aspect et al. , for it overcomes some former limitations. In particular, it overcomes the poor angular correlation of atomic cascade photons, that is at the origin of the small total efficiency of this type of experiments in which one is forced to select a small subsample of the produced photons, leading inevitably to the detection loophole, for PDC presents angular correlations better than 1 mrad The first experiments using this technique, where performed with type I PDC, which gives phase and momentum entanglement and can be used for a test of Bell inequalities using two spatially separated interferometers , as realised by Ref.. The use of beam splitters, however, strongly reduces the total quantum efficiency. In alternative, a polarisation entangled state can be generated . However, in generating this state, in most of the used configurations, half of the initial photon flux is lost and the efficiency loophole cannot be eliminated . Recently, an experiment where a polarisation entangled state is directly generated, has been realised using Type II PDC . This scheme has permitted, at the price of delicate compensations for having identical arrival time of the ordinary and extraordinary photon, a much higher total efficiency than the previous ones, which is, however, still far from the value of $`0.81`$ required for eliminating the detection loophole for a maximally entangled state (see later). Also, some recent experiments studying equalities among correlations functions rather than Bell inequalities are very far from giving a loophole free test of local realism . A large interest remains therefore for new experiments increasing total quantum efficiency in order to reduce and finally overcome the efficiency loophole. For this purpose, we have considered the possibility of generating a polarisation entangled state via the superposition of the spontaneous fluorescence emitted by two non-linear crystals (rotated for having orthogonal polarisation) driven by the same pumping laser . The crystals are put in cascade along the propagation direction of the pumping laser and the superposition is obtained by using an appropriate optics. If the path between the two crystal is smaller than the coherence length of the laser, the two photons path are indistinguishable and a polarisation entangled state is created. In fact, applying the evolution operator given by the PDC Hamiltonian one has, to the first order of the perturbative expansion: $$\|\mathrm{\Psi }=\|vacuum+f_1V_1\|H\|H+f_2V_2\|V\|V$$ (1) where $`f_i`$ keeps into account the properties of crystal $`i`$ ($`\|f_i\|^2`$ is the fraction of incident light down converted by the non-linear crystal) and $`V_i`$ the pump intensity at the crystal $`i`$. The possibility of easily obtaining a non maximally entangled state (where $`V_1f_1`$ and $`V_2f_2`$ are different) is very important, for it has been recognised that for non maximally entangled pairs the lower limit on the total detection efficiency for eliminating the detection loophole is as small as 0.67 (in the case of no background). This has to be compared with the maximally entangled pairs case, where a total efficiency larger than 0.81 is required. However, it must be noticed that, for non-maximally entangled states, the largest discrepancy between quantum mechanics and local hidden variable theories is reduced: thus a compromise between a lower total efficiency and a still sufficiently large value of this difference will be necessary when realising of an experiment addressed to overcome the detection loophole. ## 0.2 Description of the experiment The general scheme of the experiment is shown in fig. 1: two crystals of $`LiIO_3`$ (10x10x10 mm) <sup>1</sup><sup>1</sup>1Which we have measured to have $`d_{31}=3.5\pm 0.4`$ . are placed along the pump laser propagation, 250 mm apart, a distance smaller than the coherence length of the pumping laser. This guarantees indistinguishability in the creation of a couple of photons in the first or in the second crystal. A couple of planoconvex lenses of 120 mm focal length centred in between focalises the spontaneous emission from the first crystal into the second one maintaining exactly, in principle, the angular spread. A hole of 4 mm diameter is drilled into the centre of the lenses to allow transmission of the pump radiation without absorption and, even more important, without adding stray-light, because of fluorescence and diffusion of the UV radiation. This configuration, which realises a so-called ”optical condenser”, has been chosen among others, using an optical simulation program, as a compromise between minimisation of aberrations (mainly spherical and chromatic) and losses due to the number of optical components. The pumping beam at the exit of the first crystal is displaced from its input direction by birefringence: the small quartz plate (5 x5 x5 mm) in front of the first lens of the condensers compensates this displacement, so that the input conditions are prepared to be the same for the two crystals, apart from alignment errors. Finally, a half-wavelength plate immediately after the condenser rotates the polarisation of the Argon beam and excites in the second crystal a spontaneous emission cross-polarised with respect to the first one. With a phase matching angle of $`51^o`$ the spontaneous emissions at 633 and 789 nm (which are the wave lengths used for the test) are emitted at $`3.5^o`$ and $`4^o`$ respectively. The dimensions and positions of both plates are carefully chosen in order not to intersect this two stimulated emissions. We have used as photo-detectors two avalanche photodiodes with active quenching (EG&G SPCM-AQ) with a sensitive area of 0.025 $`mm^2`$ and dark count below 50 counts/s. PDC signal was coupled to an optical fiber (carrying the light on the detectors) by means of a microscope objective with magnification 20, preceded by a polariser (with extinction ratio $`10^6`$). The output signals from the detectors are routed to a two channel counter, in order to have the number of events on single channel, and to a Time to Amplitude Converter circuit, followed by a single channel analyser, for selecting and counting coincidence events. A very interesting degree of freedom of this configuration is given by the fact that by tuning the pump intensity between the two crystals, one can easily tune the value of $`f=(f_2V_2)/(f_1V_1)`$, which determines how far from a maximally entangled state ($`f=1`$) the produced state is. This is a fundamental property, which permits to select the most appropriate state for the experiment The main difficulty of this configuration is in the alignment, which is of fundamental importance for having a high visibility. This problem has been solved using a technique, that had been already applied in our laboratory for metrological studies, namely the use of an optical amplifier scheme, where a solid state laser is injected into the crystals together with the pumping laser, an argon laser at 351 nm wavelength (see fig.1). If the angle of injection is selected appropriately, a stimulated emission along the correlated direction appears, permitting to identify quite easily the two correlated directions. Then, stopping the stimulated emission of the first crystal, and rotating the polarisation of the diode laser one obtains the stimulated emission in the second crystal and can check the superposition with the former. We think that the proposed scheme is well suited for leading to a further step toward a conclusive experimental test of non-locality in quantum mechanics. The main advantage of the proposed configuration with respect to most of the previous experimental set-ups is that all the entangled pairs are selected (and not only $`<50\%`$ as with beams splitters), furthermore it does not require delicate compensations for the optical paths of the ordinary and extraordinary rays after the crystal. At the moment, the results which we are going to present are still far from a definitive solution of the detection loophole; nevertheless, being the first test of Bell inequalities using a non-maximally entangled state, they represents an important step in this direction. Furthermore, this configuration permits to use any pair of correlated frequencies and not only the degenerate ones. We have thus realised this test using for a first time two different wave lengths (at $`633`$ and $`789`$ nm). It must be acknowledged that a set-up for generating polarisation entangled pairs of photons, which presents analogies with our, has been realised recently in Ref. The main difference between the two experiments is that in the two crystals are very thin and in contact with orthogonal optical axes: this permits a ”partial” superposition of the two emissions with opposite polarisation. This overlapping is mainly due to the finite dimension of the pump laser beam, which reflects into a finite width of each wavelength emission. A much better superposition can be obtained with the present scheme, by fine tuning the crystals’ and optics’ positions and using the parametric amplifier trick. Furthermore, in the experiment of Ref. the value of $`f`$ is in principle tunable by rotating the polarisation of the pump laser, however this reduces the power of the pump producing PDC already in the first crystal, while in our case the whole pump power can always be used in the first crystal, tuning the PDC produced in the second one. More recently , they have also performed a test about local realism using a particular kind of Hardy equalities . Their result is in agreement with quantum mechanics modulo the detection loophole, no discussion concerning the elimination of loopholes for this equality is presented. As a first check of our apparatus, we have measured the interference fringes, varying the setting of one of the polarisers, leaving the other fixed. We have found a high visibility, $`V=0.973\pm 0.038`$. Our results are summarised by the value obtained for the Clauser-Horne sum $$CH=N(\theta _1,\theta _2)N(\theta _1,\theta _2^{})+N(\theta _1^{},\theta _2)+N(\theta _1^{},\theta _2^{})N(\theta _1^{},\mathrm{})N(\mathrm{},\theta _2)$$ (2) which is strictly negative for local realistic theory. In (2), $`N(\theta _1,\theta _2)`$ is the number of coincidences between channels 1 and 2 when the two polarisers are rotated to an angle $`\theta _1`$ and $`\theta _2`$ respectively ($`\mathrm{}`$ denotes the absence of selection of polarisation for that channel). On the other hand, quantum mechanics predictions for $`CH`$ can be larger than zero: for a maximally entangled state the largest value is obtained for $`\theta _1=67^o.5`$ , $`\theta _2=45^o`$, $`\theta _1^{}=22^o.5`$ , $`\theta _2^{}=0^o`$ and corresponds to a ratio $$R=[N(\theta _1,\theta _2)N(\theta _1,\theta _2^{})+N(\theta _1^{},\theta _2)+N(\theta _1^{},\theta _2^{})]/[N(\theta _1^{},\mathrm{})+N(\mathrm{},\theta _2)]$$ (3) equal to 1.207. For non-maximally entangled states the angles for which CH is maximal are somehow different and the maximum is reduced to a smaller value. The angles corresponding to the maximum can be evaluated maximising Eq. 2 with $`\begin{array}{c}N[\theta _1,\theta _2]=[ϵ_1^{\|\|}ϵ_2^{\|\|}(Sin[\theta _1]^2Sin[\theta _2]^2)+\hfill \\ ϵ_1^{}ϵ_2^{}(Cos[\theta _1]^2Cos[\theta _2]^2)\hfill \\ (ϵ_1^{}ϵ_2^{\|\|}Sin[\theta _1]^2Cos[\theta _2]^2+ϵ_1^{\|\|}ϵ_2^{}Cos[\theta _1]^2Sin[\theta _2]^2)\hfill \\ +\|f\|^2(ϵ_1^{}ϵ_2^{}(Sin[\theta _1]^2Sin[\theta _2]^2)+ϵ_1^{\|\|}ϵ_2^{\|\|}(Cos[\theta _1]^2Cos[\theta _2]^2)+\hfill \\ (ϵ_1^{\|\|}ϵ_2^{}Sin[\theta _1]^2Cos[\theta _2]^2+\hfill \\ ϵ_1^{}ϵ_2^{\|\|}Cos[\theta _1]^2Sin[\theta _2]^2)\hfill \\ +(f+f^{})((ϵ_1^{\|\|}ϵ_2^{\|\|}+ϵ_1^{}ϵ_2^{}ϵ_1^{\|\|}ϵ_2^{}ϵ_1^{}ϵ_2^{\|\|})(Sin[\theta _1]Sin[\theta _2]Cos[\theta _1]Cos[\theta _2])]/(1+\|f\|^2)\hfill \end{array}.`$ (11) where (for the case of non-ideal polariser) $`ϵ_i^{\|\|}`$ and $`ϵ_i^{}`$ correspond to the transmission when the polariser (on the branch $`i`$) axis is aligned or normal to the polarisation axis respectively. The limits for obtaining a detection loophole free experiment using non-maximally entangled states with non-ideal polarisers are summarised in fig.2, where one has maximised the Clauser-Horne inequality violation using the correlation function (11) for the non-maximally entangled case. From the figure one derives that for the ideal polariser case one can lower the limit on the total efficiency up to 0.67 for a opportune non-maximally entangled state. The use of non ideal polarisers slightly enhances this limit, however even for a extinction ratio ($`ϵ^{\|\|}ϵ^{}`$) of $`10^4`$ the effect is small (see fig. 2). Considering the extinction ratios of commercial polarisers, this simulation shows that this does not represent a significant problem. Anyway, having $`ϵ^{\|\|}<1`$ contributes to lower the detection efficiency. An important factor is the background on the single channel detection: its effect is exactly the same of a lower quantum efficiency of photodetection, for it increases the denominator in Eq. 3. The phase of $`f`$ must be kept next to zero. Any relative phase between the two components of the entangled state reflects into a reduction of Clauser-Horne inequality violation, up to reaching no violation at all for a phase difference of $`\pi /2`$. In our case we have generated a state with $`f0.4`$: in this case the largest violation of the inequality is reached for $`\theta _1=72^o.24`$ , $`\theta _2=45^o`$, $`\theta _1^{}=17^o.76`$ and $`\theta _2^{}=0^o`$, to $`R=1.16`$. Our experimental result is $`CH=512\pm 135`$ coincidences per second, which is almost four standard deviations different from zero and compatible with the theoretical value predicted by quantum mechanics. In terms of the ratio (3), our result is $`1.082\pm 0.031`$. For the sake of comparison, one can consider the value obtained with the angles which optimise Bell inequalities violation for a maximally entangled state. The result is $`CH=92\pm 89`$, which, as expected, shows a smaller violation than the value obtained with the correct angles setting. ## 0.3 Conclusions We have presented the first measurement of the violation of Clauser-Horne inequality (or of other Bell inequalities) using a non-maximally entangled state. This represents a relevant step in the direction of eliminating the detection loophole. Further developments in this sense are the purpose of this collaboration. Besides, this scheme gives a beautiful example of how entanglement and quantum interference derive from absence of welcher weg (which path) information , namely from the impossibility of distinguishing in which crystal the pair is produced. In fact, this scheme works substantially as a quantum eraser: when no polariser is inserted after the crystals, interference is cancelled because of different polarisations produced in the two crystals, when a polariser, rotated respect to vertical (horizontal) axis, is inserted, due to (at least partial) cancellation of the information about where PDC has happened interference is restored (and could be modulated changing the phase between the pair produced in the first and second crystal). ## 0.4 Acknowledgments We would like to acknowledge the support of the Italian Space Agency under contract LONO 500172 and of MURST via special programs ”giovani ricercatori”, Dip. Fisica Teorica Univ. Torino. Figures captions \- Fig. 1 Sketch of the source of polarisation entangled photons. CR1 and CR2 are two $`LiIO_3`$ crystals cut at the phase-matching angle of $`51^o`$. L1 and L2 are two identical piano-convex lenses with a hole of 4 mm in the centre. P is a 5 x 5 x 5 mm quartz plate for birefringence compensation and $`\lambda /2`$ is a first order half wave-length plate at 351 nm. U.V. identifies the pumping radiation at 351 nm. The infrared beam (I.R.) at 789 nm is generated by a diode laser and is used for system alignment only. The parametric amplifier scheme, described in the text, is shown as well. The dashed line identifies the idler radiation at 633 nm. A second half-wave plate on the I.R. beam (not shown in the figure) allows amplified idler emission from the second crystals too. The figure is not in scale. \- Fig. 2 a) Contour plot of the quantity $`CH/N`$ (see Eq. 2. N is the total number of detections) in the plane with $`f`$ (non maximally entanglement parameter, see the text for the definition) as y-axis and $`\eta `$ (total detection efficiency) as x-axis. The polarisers are supposed to have $`ϵ_i^{\|\|}=1`$. The leftmost region corresponds to the region where no detection loophole free test of Bell inequalities can be performed. The contour lines are at 0, 0.01, 0.1, 0.15, 0.2. b) the same as (a), but with an extinction ratio $`10^4`$.
warning/0004/cond-mat0004270.html	ar5iv	text	# Lattice models and Landau theory for type II incommensurate crystals ## I Introduction Like other transitions, the phase transition to an incommensurate (INC) phase (reviews are given by Bak, Selke and Janssen and Janner ) can be described on the phenomenological level within the frame of extended Landau theory. The necessary extension consists essentially in accounting for the expansion of the free energy density as a function not only of the components of the order parameter, but also of their spatial derivatives. Therefore the global free energy becomes a functional of spatially dependent components of the order parameter and the equilibrium configuration for given values of temperature and external parameters is found as a solution of a variational problem. The continuum Landau theory allows a natural classification of the possible forms of the free-energy functional for an INC transition into two classes, according to whether the driving term in the free energy expansion responsible for the appearance of the incommensurate state is linear (type I, Lifshitz invariant present) or quadratic (type II, no Lifshitz invariant) in the gradient of the order parameter. The properties of those two kinds of INC phases are different: for type I INC phases the lock-in transition is either continuous, or only slightly discontinuous, and approaching the lock-in temperature $`T_c`$ these phases exhibit the structuration of the modulated phase into discommensurations or solitons. On the other hand, the modulation of the type II INC phase remains practically sinusoidal down to the temperature $`T_c`$ and the lock-in transition is always of first order. Although the above statements can be considered as a rule of thumb, there are cases known where there is coexistence of solitonic and sinusoidal structural modulation. See Aramburu for details. In the following we will only be concerned with models describing type II INC phases. Landau theory has been rather successful in describing basic properties of these phases, but if one wants to have a better understanding of the true microscopic origin of the INC phase one has to go beyond this phenomenological approach. One possibility would be to study full microscopic models with realistic interactions. Another approach, to which the main part of this paper will be devoted, is to study semi-microscopic models which take into account the discrete nature of the systems and discuss properties in terms of (effective) inter-atomic interactions. The discreteness of a lattice leads to a number of important physical consequences such as pinning of solitons in the anisotropic next-nearest-neighbor Ising (ANNNI) model (see Yeomans for a review) and in the Frenkel-Kontorova model, and the existence of a devil’s staircase (infinite number of commensurate and incommensurate phases), for example found in betaine calcium chloride dihydrate (BCCD). Within Landau theory it is difficult to explain the occurrence of a specific sequence of transitions: several lock-in terms are needed then. For example, Ribeiro et al. chose the magnitude of four distinct lock-in contributions to the free energy in such a way as to stabilize the four most prominent commensurate phases in BCCD. Furthermore, in discrete models chaotic states are possible, which may provide an alternative description of phenomena observed in for example spin glasses, superionic conductors, the magnetic system CeSb and systems with pinning of charge density waves. See Bak for details. In the past few years different lattice models have been constructed, for example to describe the phase transitions in the A<sub>2</sub>BX<sub>4</sub> family, in BCCD and, more general, in crystals with $`Pcmn`$ symmetry. These models are 2-dimensional, with only nearest-neighbor interactions. Hlinka et al. studied a 3-dimensional nearest-neighbor model, applicable to BCCD. All these models have in common that the frustrated interaction, needed for having an incommensurate phase, comes from a nearest-neighbor mixing interaction. Recent X-ray, neutron and Raman experiments on the Sn<sub>2</sub>P<sub>2</sub>(S<sub>1-x</sub>Se<sub>x</sub>)<sub>6</sub> crystal family of uniaxial ferroelectrics motivated us to study lattice models. In the composition-temperature phase diagram of this crystal family a Lifshitz point is present. At this point the paraelectric phase, the ferroelectric phase and the INC phase become equal, and the boundaries separating these phases have equal derivative. Such a special point was, except for some ferroelectric liquid crystals, only found in the temperature-applied magnetic field phase diagram of the magnetic compound MnP. From both experimental and theoretical point of view this uniaxial Lifshitz point is interesting because critical exponents deviate from those found for ordinary critical points. This crystal family furthermore displays an interesting modulation wave vector behavior, shows cross-over effects from order-disorder to a displacive type of phase transition, and the ratio between the slopes of the soft phonon mode in the ferroelectric and paraelectric phase deviates substantially from the standard value $`R=2`$. These features are much easier to understand in a lattice model than in Landau theory. The paper is arranged as follows: in Section II we present a one-dimensional model; in many anisotropic systems, like Sn<sub>2</sub>P<sub>2</sub>(S<sub>1-x</sub>Se<sub>x</sub>)<sub>6</sub>, the modulation wave vector is in one specific direction. The incommensurability may arise if there is frustrating interaction between nearest-neighbor and next-nearest-neighbor couplings. We discuss general features of the model. In Section III we give some exact results regarding ground state properties. In Section IV we discuss the dynamics and the stability of the various phases. In Section V the phase diagrams, calculated partly analytically, partly numerically, are presented. Temperature effects are treated in Section VI. In Section VII we discuss the continuum limit of the model which gives the connection with the Landau theory. We conclude and give an outlook for further research in Section VIII. In Appendix A we give some exact results for phase boundaries and in Appendix B we present the next-nearest-neighbor extension of the Effective Potential Method for the determination of the ground state. This method was used to calculate some of the phase diagrams. ## II An extension of the DIFFOUR model In the following we will be concerned with an extension of the so-called DIFFOUR model (discrete frustrated $`\varphi ^4`$ model), for which the potential energy can be written as $`V={\displaystyle \underset{n}{}}\{`$ $`{\displaystyle \frac{A}{2}}x_n^2+{\displaystyle \frac{B}{4}}x_n^4+{\displaystyle \frac{C}{2}}\left(x_nx_{n1}\right)^2+{\displaystyle \frac{D}{2}}\left(x_nx_{n2}\right)^2`$ (2) $`+{\displaystyle \frac{E}{2}}[x_n^2(x_nx_{n1})^2+x_n^2(x_nx_{n+1})^2]\}.`$ The original DIFFOUR model, or EHM (elastically hinged molecule) model, has $`E=0`$. Although in principle this model gives incommensurate ground states, the behavior of the modulation wave vector as found in experiments can not be reproduced satisfactorily by the model. In order to account for this shortcoming we supply the DIFFOUR model with a non-linear coupling to neighbors. There are several possibilities: if we restrict ourselves to fourth-order terms we can consider a term $`(x_nx_{n1})^4`$. The resulting model has been studied by Lamb, who showed that the origin of this term is related to strain terms in the so-called magnetoelastic DIFFOUR model. Another possibility would be to consider a term of the form $`(x_nx_{n2})^4`$. Without the term mentioned above this would however be rather unphysical. Instead we will choose a term of the form $`x_n^2(x_nx_{n\pm 1})^2`$. This is the lowest-order dispersive fourth-order term, as will be shown in Section VII, and can, for example, be obtained from strain terms. The order parameter $`x_n`$ can be, for example, a displacement, a component of the polarization $`P`$ for ferroelectric systems, a component of the magnetization $`M`$ for magnetic systems, a rotation angle or a strain component. In this article we use for convenience terms like paraelectric, ferroelectric and antiferroelectric to distinguish between different ground states. The origin of incommensurability in this model is essentially competition between interactions with nearest- and next-nearest neighbors which may lead to frustration. Higher order terms are needed for stabilization. We expect that the extra fourth-order term $`(E0)`$ has a large effect on the phase diagrams for $`E=0`$. To actually determine phase diagrams it is not necessary to vary all 5 parameters $`A,B,C,D,E`$, which can be seen as follows: by taking $`x_n^{}=\sqrt{B/\|D\|}x_n`$ and $`V^{}=B/\|D\|^2V`$ we get the following renormalized parameters: $`A^{}=A/\|D\|`$, $`B^{}=1`$, $`C^{}=C/\|D\|`$, $`D^{}=D/\|D\|=\pm 1`$, and $`E^{}=E/B`$. For some purposes it is convenient to rewrite the potential in the following form: $`V={\displaystyle \underset{n}{}}\{`$ $`{\displaystyle \frac{\stackrel{~}{A}}{2}}x_n^2+{\displaystyle \frac{\stackrel{~}{B}}{4}}x_n^4+\stackrel{~}{C}x_nx_{n1}+\stackrel{~}{D}x_nx_{n2}`$ (4) $`+{\displaystyle \frac{\stackrel{~}{E}}{2}}[x_n^2(x_nx_{n1})^2+x_n^2(x_nx_{n+1})^2]\},`$ with $`\stackrel{~}{A}=A+2C+2D`$, $`\stackrel{~}{B}=B`$, $`\stackrel{~}{C}=C`$, $`\stackrel{~}{D}=D`$ and $`\stackrel{~}{E}=E`$. From this the connection with the ANNNI model can easily be made. Let us put $`\stackrel{~}{E}=0`$ and $`\stackrel{~}{B}=\stackrel{~}{A}`$. If we now take the limit $`\stackrel{~}{A}\mathrm{}`$ we end up with a model with two infinitely deep wells. The $`x_n`$ can only take on values $`\pm 1`$ and can thus be seen as spins. These spins are coupled to nearest-neighbors and next-nearest-neighbors via the $`\stackrel{~}{C}`$ and $`\stackrel{~}{D}`$ terms. So by increasing the depth of the double-well potential there is a crossover from displacive behavior to order-disorder behavior in the transition from the normal to the incommensurate phase. Inserting $`x_n^{}=(1)^nx_n`$ in the above potential leads to $`V={\displaystyle \underset{n}{}}\{`$ $`{\displaystyle \frac{\stackrel{~}{A}}{2}}x_{n}^{}{}_{}{}^{2}+{\displaystyle \frac{\stackrel{~}{B}}{4}}x_{n}^{}{}_{}{}^{4}\stackrel{~}{C}x_n^{}x_{n1}^{}+\stackrel{~}{D}x_n^{}x_{n2}^{}`$ (6) $`+{\displaystyle \frac{\stackrel{~}{E}}{2}}[x_{n}^{}{}_{}{}^{2}(x_n^{}+x_{n1}^{})^2+x_{n}^{}{}_{}{}^{2}(x_n^{}+x_{n+1}^{})^2]\}.`$ In the DIFFOUR model ($`\stackrel{~}{E}=0`$) this leads to the following symmetry: if $`\{x_n\}`$ is a state for $`\stackrel{~}{C}=X`$, then $`\{(1)^nx_n\}`$ is a state for $`\stackrel{~}{C}=X`$ with the same energy. This property can for example be seen in the ferroelectric-antiferroelectric phases. However, for $`\stackrel{~}{E}0`$ this symmetry $`\stackrel{~}{C}\stackrel{~}{C}`$ is no longer present. ## III Ground state properties Different ground states are possible, depending on the values of the parameters. The stationary states are solutions of $`V/x_n=0`$, giving $`Ax_n+Bx_n^3+C\left(2x_nx_{n1}x_{n+1}\right)+D\left(2x_nx_{n2}x_{n+2}\right)`$ (7) $`+E\left[4x_n^33x_n^2\left(x_{n1}+x_{n+1}\right)+2x_n\left(x_{n1}^2+x_{n+1}^2\right)\left(x_{n1}^3+x_{n+1}^3\right)\right]=0.`$ (8) If we impose periodic boundary conditions $`x_{N+n}=x_n`$ we arrive at a set of $`N`$ coupled non-linear equations. To find the lowest-energy state for each solution and for each value of $`N`$ the potential energy has to be evaluated. For low-period commensurate states (small values of $`N`$) analytic solutions of the above equation can be found. In the following, we study (for fixed $`N`$) periodic solutions of (8). For them we give the equilibrium values $`\{x_n\}`$ and the corresponding energy per particle $`v=V/N`$. In the paraelectric state $`(N=1)`$ all particles are in the equilibrium positions $$x_n=0,v=0.$$ (9) In the ferroelectric state $`(N=1)`$ all particles are uniformly displaced from their equilibrium positions $$x_n=\sqrt{\frac{A}{B}},v=\frac{A^2}{4B}.$$ (10) Note that $`B`$ always has to be positive, for the potential to be bounded from below. This implies that the ferroelectric state only can exist for $`A<0`$. For $`A>0`$ the ground state may be paraelectric. In the following we give some analytic results, based on numerical calculations of the shape of the solution $`\{x_n\}`$. In the antiferroelectric state $`(N=2)`$ particle positions alternate along the chain $$x_n=(1)^n\sqrt{\frac{A+4C}{B+16E}},v=\frac{\left(A+4C\right)^2}{4(B+16E)}.$$ (11) The potential is unbounded from below for $`EB/16`$ and stable solutions exist only for $`E>B/16`$ and $`A+4C<0`$. Both conditions have to be satisfied. For $`N=3`$ we determined the solution with lowest energy to be of the form $`(x_1,x_2,x_3)=(k\xi ,\xi ,k\xi )`$, with $`x_1/x_2<0`$, and $$\xi ^2=\frac{A+2(C+D)(1k)}{B+2E(23k+2k^2k^3)}=\frac{Ak+(C+D)(k1)}{Bk^3+E(2k^33k^2+2k1)}.$$ (12) The factor $`k`$ is determined by $`2\left[B(C+D)+E(A+C+D)\right]k^4\left[B(A+2C+2D)+2E(A+2C+2D)\right]k^3`$ (13) $`3AEk^2+\left[B(A+C+D)+2E(A+2C+2D)\right]k\left[B(C+D)+E(A+2C+2D)\right]=0.`$ (14) The energy per particle is given in terms of $`\xi `$ and $`k`$ as $$v=\frac{\xi ^2}{6}\left[A(1+2k^2)+2(C+D)(1k)^2\right]+\frac{\xi ^4}{12}\left[B(1+2k^4)+4E(1+k^2)(1k)^2\right].$$ (15) Note that the quartic equation (13) can be written as $`(1k)\{`$ $`2\left[B(C+D)+E(A+C+D)\right]k^3+\left[BA+2E(C+D)\right]k^2`$ (17) $`+[BA+E(3A+2C+2D)]k[B(C+D)+E(A+2C+2D)]\}=0,`$ where the special solution $`k=1`$ gives a ferroelectric state. The solution of the remaining cubic equation, which can be solved exactly for given parameters, gives a $`k`$ such that $`x_1/x_2<0`$, a true $`N=3`$ state. The lowest energy state for $`N=4`$ has $`x_1=x_2=\rho ,x_3=x_4=\rho `$ with $$\rho =\sqrt{\frac{A+2C+4D}{B+8E}},v=\frac{(A+2C+4D)^2}{4(B+8E)}.$$ (18) We have to keep in mind that we must satisfy $`E>B/16`$, which is not obvious from the above expression, but comes from the analysis of the antiferroelectric state. The lowest energy solution for $`N=6`$ can analytically be obtained, in the same manner as for $`N=3`$. It has the form $`(x_1,x_2,x_3,x_4,x_5,x_6)=(k\xi ,\xi ,k\xi ,k\xi ,\xi ,k\xi )`$. The lowest energy solution for $`N=8`$ reads $`(x_1,x_2,x_3,x_4,x_5,x_6,x_7,x_8)=(k\xi ,\xi ,\xi ,k\xi ,k\xi ,\xi ,\xi ,k\xi )`$. For $`N=5`$ one needs 2 different values of $`k`$: $`(x_1,x_2,x_3,x_4,x_5)=(k\xi ,k^{}\xi ,\xi ,k^{}\xi ,k\xi )`$ and for $`N=7`$ one needs three different values: $`k,k^{},k^{\prime \prime }`$, and the above method no longer works for $`N=5`$ and $`N=7`$. Therefore, to find states with larger periods, or even incommensurate periods, we rely on numerical calculations, for which true incommensurate states of course never can be found. However, the idea is that such a state can always be arbitrary well approximated by a commensurate state with wavelength $$\lambda =\frac{N}{s},s=1,2,\mathrm{},N2s,$$ (19) where $`N,s`$ are coprime numbers. Such a solution has a period $`N`$ and, in general, $`2s=`$ (number of local minima $`+`$ number of local maxima). In the special case where the $`\{x_n\}`$ take on positive and negative values, $`2s=`$ (number of sign changes within the period $`N`$). The bigger $`N`$ and $`s`$, the better the approximation. As an example of a numerical calculation we consider the ground state for $`A=2.24999,B=1,C=1,D=1,E=1`$. It is known to be incommensurate (see Sections IV and V) with a wavelength $`\mathrm{arccos}(\frac{1}{4})4.76679213`$. By the Farey construction we find that $`\frac{62}{13}4.76923077`$ should be a reasonable commensurate approximation. We numerically determined the ground state in terms of the $`\{x_n\}`$. The result is shown in Figure 1. After 62 particles the sequence repeats itself, and the solution passes through zero 26 times. We can label this state by its modulation wave vector $`\frac{13}{62}`$, measured in units of $`2\pi `$. For certain regimes in the parameter space the ground state can be determined analytically. The entire phase boundary of the paraelectric state and a part of the phase boundary of the ferroelectric state can be calculated. For the nearest-neighbor case in the DIFFOUR model proofs are given by Janssen and Tjon. We extend their proofs to the case in which we also have next-nearest-neighbor interaction. As the proof is rather lengthy it will be given in Appendix A. ## IV Phonon dispersion curves and stability limits To decide whether a solution of the equilibrium conditions is locally stable or not, one considers small displacements $`ϵ_n`$ from the positions given by a static solution $`\{x_n\}`$ satisfying (8): $$u_n=x_n+ϵ_n.$$ (20) The phonon frequencies are given by the square roots of the eigenvalues of the dynamical matrix and for stability all eigenvalues have to be non-negative. The dynamical matrix for a period-$`N`$ solution $`(N5)`$ has elements $`D_{n,n}`$ $`=`$ $`A+3Bx_n^2+2C+2D+2E\left[6x_n^23x_n(x_{n1}+x_{n+1})+x_{n1}^2+x_{n+1}^2\right],`$ (21) $`D_{n,n\pm 1}`$ $`=`$ $`C+E\left[3x_n^2+4x_nx_{n\pm 1}3x_{n\pm 1}^2\right],`$ (22) $`D_{n,n\pm 2}`$ $`=`$ $`D,`$ (23) with $`x_{N+n}=x_n`$ and the special cases $`D_{1,N}`$ $`=`$ $`\left[C+E\left(3x_1^2+4x_1x_N3x_N^2\right)\right]e^{iq},`$ (24) $`D_{1,N1}`$ $`=`$ $`D_{2,N}=De^{iq}.`$ (25) In the above expressions the $`\{x_n\}`$ are solutions of (8). Furthermore $`D_{n,m}=D_{m,n}^{}`$ and all other matrix elements are zero. In the following the phonon branches for certain low-period states will be examined. This will be done in terms of $`A,B,C,D,E`$. ### A Paraelectric state For the paraelectric state (9) the dynamical matrix is given by $$𝖣=A+2C(1\mathrm{cos}(q))+2D(1\mathrm{cos}(2q))=m\omega ^2.$$ (26) Rewriting this equation gives $$m\omega ^2=A+(4C+16D)\mathrm{sin}^2(\frac{q}{2})16D\mathrm{sin}^4(\frac{q}{2}).$$ (27) Note that this expression does not contain $`E`$. This means that the stability limits for the paraelectric phase in this model are the same as those for the paraelectric state in the DIFFOUR model. Now, we are looking for the minimum of this phonon branch. We distinguish the cases $`D>0`$ and $`D<0`$. The results are summarized in Table I. For $`D>0`$ and $`C=0`$, the branch has two minima, at $`\mathrm{sin}^2(\frac{q}{2})=0`$ and $`\mathrm{sin}^2(\frac{q}{2})=1`$. For $`A<0`$ the ferroelectric state and the antiferroelectric state are degenerate, for $`E=0`$ only. Comparison with calculations in Appendix A shows that as long as the paraelectric state is stable, it is the ground state. Destabilization is the condensation of a soft phonon. ### B Ferroelectric state For the ferroelectric state (10) one has $`m\omega ^2`$ $`=`$ $`A+3Bx^2+2C(1\mathrm{cos}(q))+2D(1\mathrm{cos}(2q))+4Ex^2(1\mathrm{cos}(q))`$ (28) $`=`$ $`2A+(4C+16D{\displaystyle \frac{8AE}{B}})\mathrm{sin}^2({\displaystyle \frac{q}{2}})16D\mathrm{sin}^4({\displaystyle \frac{q}{2}}).`$ (29) See Table II for the analysis. Note that for $`D>0`$ there is degeneracy for $`E=0`$. ### C Antiferroelectric state Finally, the phonon branches of the antiferroelectric state (11) will be investigated. The $`2\times 2`$ dynamical matrix is given by $`D_{1,1}=D_{2,2}`$ $`=`$ $`A+2C+2D(1\mathrm{cos}(q))(3B+28E){\displaystyle \frac{A+4C}{B+16E}},`$ (30) $`D_{1,2}=D_{2,1}^{}`$ $`=`$ $`(e^{iq}+1)(C+10E{\displaystyle \frac{A+4C}{B+16E}}).`$ (31) The eigenvalue equation reads $`m\omega ^2`$ $`=`$ $`A+2C+4D(3B+28E){\displaystyle \frac{A+4C}{B+16E}}4D\mathrm{cos}^2({\displaystyle \frac{q}{2}})`$ (33) $`\pm 2(C+10E{\displaystyle \frac{A+4C}{B+16E}})\mathrm{cos}({\displaystyle \frac{q}{2}}).`$ Results are summarized in Table III. Note that for $`C=10E\frac{A+4C}{B+16E}`$ the two branches coincide. ### D States with period $`N3`$ For the $`N=3`$ solution exact phonon frequencies can in principle be found. The elements of the dynamical matrix are given in terms of $`\xi `$ and $`k`$, defined in (12) and (13): $`D_{1,1}`$ $`=`$ $`D_{3,3}=A+2(C+D)+\left[3Bk^2+2E(4k^23k+1)\right]\xi ^2,`$ (34) $`D_{2,2}`$ $`=`$ $`A+2(C+D)+\left[3B+4E(k^23k+3)\right]\xi ^2,`$ (35) $`D_{1,2}`$ $`=`$ $`D_{2,3}=CE(3k^24k+3)\xi ^2De^{iq},`$ (36) $`D_{1,3}`$ $`=`$ $`D(C+2Ek^2\xi ^2)e^{iq}.`$ (37) The eigenvalues are then found as the solution (Cardano’s formula) of a cubic equation. For $`N=4`$ the dynamical matrix has elements in terms of $`\rho `$, defined in (18): $`D_{n,n}`$ $`=`$ $`A+2(C+D)+(3B+16E)\rho ^2,`$ (38) $`D_{1,2}`$ $`=`$ $`D_{3,4}=C2E\rho ^2,`$ (39) $`D_{1,3}`$ $`=`$ $`D_{2,4}=D(1+e^{iq}),`$ (40) $`D_{1,4}`$ $`=`$ $`(C10E\rho ^2)e^{iq},`$ (41) $`D_{2,3}`$ $`=`$ $`C10E\rho ^2.`$ (42) The resulting secular equation is a quartic one and exact solutions for the eigenvalues can be found using Ferrari’s formula. For solutions with larger periods we have to rely on numerical calculations. As an example we again consider the ground state for $`A=2.24999,B=1,C=1,D=1,E=1`$. See also the end of Section III. The calculated phonon dispersion curves in the commensurate approximation $`\lambda =\frac{62}{13}`$ are given in Figure 2. ## V Calculation of phase diagrams In this section we present some phase diagrams, calculated partly analytically, partly numerically. The traditional method to find the ground state numerically is to solve the equations for equilibrium (8). However, these equations also hold for metastable states, maxima and saddlepoints and it may happen that one finds a metastable state instead of the true ground state. This problem is not present for the so-called effective potential method (EPM), introduced by Griffiths and Chou. This method in principle always gives the ground state. Originally it was used to study Frenkel-Kontorova and similar one-dimensional models with only nearest-neighbor interaction. As an interesting application of this method, we mention a study of the ground state of the chiral XY model in a field. Below we give a brief outline of the method. Consider a one-dimensional system with only nearest-neighbor interaction in its ground state. If one atom is displaced from its equilibrium position (we assume that $`x_n`$ denotes the displacement), the surrounding atoms will change their positions in order to minimize the total energy. This deformation will in general cost some energy. A function, called the ‘effective potential’, describes the net energy cost as a function of the positions of the atoms. This effective potential achieves its minimum on points of the ground state and rigorous mathematical statements can be made. Numerical procedures to find solutions are based on discretization of the range $`x_n`$ of the atomic positions. The $`x_n`$ can now only adopt a finite number of values. For models with interactions up to next-nearest neighbors, as in the case of the extended DIFFOUR model, the EPM can be adapted, which will be discussed in Appendix B. The proofs of the existence of solutions for models with next-nearest-neighbor interactions, both in the continuous and discretized version, are rather long and will be given in a separate paper. Using the EPM and equation (8) we calculated various phase diagrams. First we varied both $`A`$ and $`E`$ with the other parameters fixed: $`B=1`$, $`C=1`$ and $`D=1`$. The resulting phase diagram is given in Figure 3. From the analysis in Appendix A we know that the phase boundary for the paraelectric state for $`\|C\|<4`$ is given by $`A=\frac{1}{4}C^22C+4`$. For $`C=1`$ and $`D=1`$ we find $`A=2\frac{1}{4}`$. At this boundary we have a transition to an incommensurate state with wave vector $`q=\mathrm{arccos}(\frac{C}{4D})=\mathrm{arccos}(\frac{1}{4})4.76679213`$. This is the state we discussed at the end of Section III. We can clearly see the effect of the $`E`$-term: for $`E=0`$ further decreasing of $`A`$ leads to a transition to a commensurate state with period 4. For $`E<0`$ this transition can be followed by transitions to $`N=3`$ or $`N=2`$ commensurate states. For $`E`$ sufficiently positive, the wavelength of the ground state increases for decreasing $`A`$. Between the commensurate states incommensurate ones can be found. By increasing $`E`$, the region between paraelectric and ferroelectric phases shrinks. A positive $`E`$ term favors long-wavelength solutions, with the ferroelectric state being the extreme limit $`(q=0)`$. Figure 4 gives the phase diagram found for $`B=1`$, $`D=1`$, $`E=0`$ and varying both $`A`$ and $`C`$. This is the phase diagram for the original DIFFOUR model. We have seen that the phase boundary for the paraelectric phase for $`\|C\|<4`$ is a parabola symmetric around $`C=0`$. For $`\|C\|4`$ this boundary is given by $`A+2C2=2+2\|C\|`$, two straight lines. The parabola and the lines meet at $`\|C\|=4`$, and have equal derivative at this point. Note the symmetry $`CC`$, which implies that the modulation wave vectors for the system with $`+C`$ and $`C`$ are related by $$q_C+q_C=\frac{1}{2}$$ (43) in units of $`2\pi `$. At $`(C=4,A=0)`$ the paraelectric phase, the ferroelectric phase and the incommensurate phase become equal. The lines separating the paraelectric phase from the incommensurate phase and the incommensurate phase from the ferroelectric phase have equal derivative at this point. In Landau theory (see Section VII) such a point would be called a Lifshitz point. From the symmetry $`CC`$ it is obvious that there is also a Lifshitz point at $`(C=4,A=16)`$. At this point the paraelectric phase, the antiferroelectric phase and the incommensurate phase become equal. Figure 5 gives the phase diagram for $`B=1`$, $`D=1`$, $`E=1`$ in terms of $`A`$ and $`C`$. The symmetry $`CC`$ is no longer present. However, the phase boundary of the paraelectric phase is independent of $`E`$. Also the wavelength of the phase emanating from this boundary is the same. In particular the positions of the Lifshitz points and the derivates at these points do not change. Note the boundary of the antiferroelectric phase: starting from the Lifshitz point and going down in the phase diagram, it initially bends to the right and then returns to lower values of $`C`$. Figures 4 and 5 have been obtained by solving equation (8) and comparing the energies of the solutions. We investigated the nearby surroundings of the Lifshitz point at $`(C=4,A=0)`$ to look how the transition line from the ferroelectric phase to the incommensurate phase changes by increasing $`E`$. See figure 6 for the results. One notices a tendency towards a smaller wedge $`W`$ (the vertical distance between the paraelectric-incommensurate phase boundary and the incommensurate-ferroelectric phase boundary) by increasing $`E`$, as was to be expected. ## VI Temperature dependent behavior As the model under consideration is one-dimensional with short-range interactions, there is no phase transition possible at $`T0`$. If we however consider weakly interacting linear chains in a three dimensional system, this system can be described by (2) as well if we interpret the variables $`x_n`$ as averages over planes perpendicular to a fixed direction (the $`c`$-axis). Phase transitions become possible due to inter-chain couplings. To study the temperature dependence of the parameters we take the thermal average of the conditions for equilibrium (8), resulting in $`Ax_n+Bx_n^3+C\left(2x_nx_{n1}x_{n+1}\right)+D\left(2x_nx_{n2}x_{n+2}\right)`$ (44) $`+E\left[4x_n^33x_n^2x_{n1}3x_n^2x_{n+1}+2x_nx_{n1}^2+2x_nx_{n+1}^2x_{n1}^3x_{n+1}^3\right]=0`$ (45) We have to distinguish between ground states with $`\{\overline{x}_n\}=0`$ and ground states with $`\{\overline{x}_n\}0`$, where the $`\{\overline{x}_n\}`$ are solutions of (8). In the former case we assume that the thermal fluctuations of the displacement $`x_n`$ do not depend on the lattice site, $`x_n^2x_n^2x_m^2x_m^2`$, and if we furthermore approximate the correlations by $`x_n^2x_mx_n^2x_m`$, the following holds: $`E\left[3x_n^2x_{n1}3x_n^2x_{n+1}+2x_nx_{n1}^2+2x_nx_{n+1}^2x_{n1}^3x_{n+1}^3\right]`$ (47) $`+(B+4E)x_n^3`$ $``$ $`E\left[3x_n^2\left(x_{n1}+x_{n+1}\right)+2x_n\left(x_{n1}^2+x_{n+1}^2\right)\left(x_{n1}^3+x_{n+1}^3\right)\right]`$ (50) $`+4E\left(x_n^2x_n^2\right)\left(2x_nx_{n1}x_{n+1}\right)`$ $`+(B+4E)x_n^3+B\left(x_n^2x_n^2\right)x_n.`$ Inserting the last expression in (45) we see that the conditions for equilibrium for the thermal average of the displacement, $`x_n`$, have the same form as those for the displacements $`x_n`$ themselves; the only difference being the replacement of the parameters $`A`$ and $`C`$ by temperature dependent ones: $`A`$ $``$ $`A+BT,`$ (51) $`C`$ $``$ $`C+4ET,`$ (52) where $`T=x_n^2x_n^2`$ is a measure of the thermal fluctuations. So a change in temperature will renormalize both parameters $`A`$ and $`C`$ (unlike in the DIFFOUR model with $`E=0`$). For all other ground state solutions ($`\overline{x}_n0`$) we calculate the thermal averages around $`\overline{x}_n`$, where the $`\{\overline{x}_n\}`$ satisfy (8), $$x_n^px_m^q=\frac{x_n^px_m^qe^{\beta (x_n\overline{x}_n)^2/2}e^{\beta (x_m\overline{x}_m)^2/2}\text{d}x_n\text{d}x_m}{e^{\beta (x_n\overline{x}_n)^2/2}e^{\beta (x_m\overline{x}_m)^2/2}\text{d}x_n\text{d}x_m}=x_n^px_m^q,$$ (53) where $`\beta =1/T`$. Three different integrals have to be calculated, yielding $`x_n^3`$ $`=`$ $`\overline{x}_n^3+3\overline{x}_nT`$ (54) $`x_n^2`$ $`=`$ $`\overline{x}_n^2+T`$ (55) $`x_n`$ $`=`$ $`\overline{x}_n`$ (56) Substitution in (45) and comparison with (8) then leads to $`A`$ $``$ $`A+(3B+4E)T,`$ (57) $`C`$ $``$ $`C+6ET.`$ (58) Note that the parameter $`E`$ now also enters in the temperature dependence of $`A`$. This linear behavior in $`T`$, with a kink at the temperature where the transition from the paraelectric phase to the incommensurate or the ferroelectric phase takes place, is corroborated by Monte Carlo calculations. Some of the results are shown in Figure 7. The results (52) and (58) are in sharp contrast with the assumptions made in standard Landau theory, to be discussed in Section VII, that there is only one temperature dependent parameter, and that its behavior above and below the transition temperature is the same. It is now straightforward to calculate the temperature dependent ground states and stability limits by making substitutions (52) and (58). We especially would like to focus on the phonon branches in the paraelectric and ferroelectric phases. Of experimental interest is the ratio between the slopes of the soft phonon mode in the ferroelectric and paraelectric phase, the so-called $`R`$-parameter: $$R=\frac{\mathrm{d}\omega ^2/\mathrm{d}T\|_{\mathrm{ferro}}}{\mathrm{d}\omega ^2/\mathrm{d}T\|_{\mathrm{para}}}.$$ (59) Self-consistent renormalized phonon theory gives the result $`R=2`$. In experiments however very often $`R2`$ is found. Taking into account the temperature dependence given in (52) and (58), we find using (27) and (29) $$R=\frac{2(3B+4E)32E^2/B\mathrm{sin}^2(q/2)}{B+16E\mathrm{sin}^2(q/2)},$$ (60) which, for $`E=0`$, gives (at the center of the Brillouin zone) $`R=6`$ instead of $`R=2`$, and for $`E0`$ can take on any value as long as $`B+16E>0`$ is satisfied. Molecular dynamics simulations on a 3-dimensional $`\varphi ^4`$ lattice by Padlewski et al. show that $`R=2`$ holds only for systems with long-range couplings being in the displacive limit. However, Sollich et al. showed that this double limit of displaciveness and long range interaction is not necessary if the system is displacive enough: they found $`R=2`$ for a system with only nearest-neighbor interactions, thereby questioning Padlewski’s claim of having studied a system in the displacive limit. ## VII Comparison with the continuum theory The continuum limit of the extended DIFFOUR model leads to a well-known expansion. By replacing the differences in the general order parameter $`x_n`$ by differentials in the principal order parameter for ferroelectrics, the polarization $`P(z)`$, $`(x_nx_{n1})^2`$ $``$ $`\left({\displaystyle \frac{\mathrm{d}P}{\mathrm{d}z}}\right)^2,`$ (61) $`(x_{n1}+x_{n+1}2x_n)^2`$ $``$ $`\left({\displaystyle \frac{\mathrm{d}^2P}{\mathrm{d}z^2}}\right)^2,`$ (62) and rearranging terms we arrive at the free energy density $$f=\frac{\alpha }{2}P^2+\frac{\beta }{4}P^4+\frac{\kappa }{2}\left(\frac{\mathrm{d}P}{\mathrm{d}z}\right)^2+\frac{\lambda }{2}\left(\frac{\mathrm{d}^2P}{\mathrm{d}z^2}\right)^2+\frac{\eta }{2}P^2\left(\frac{\mathrm{d}P}{\mathrm{d}z}\right)^2,$$ (63) with $`\alpha =A`$, $`\beta =B`$, $`\kappa =C+4D`$, $`\lambda =D`$ and $`\eta =2E`$. The above free energy density was used by Ishibashi and Shiba to study phase transitions in NaNO<sub>2</sub> and SC(NH<sub>2</sub>)<sub>2</sub> (thiourea), proper ferroelectrics in which the polarization component of interest transforms according to a one-dimensional irreducible representation. The $`\eta `$-term is allowed by symmetry because it is the product of the two invariants $`P^2`$ and $`(\mathrm{d}^2P/\mathrm{d}z^2)^2`$. Alternatively, as both sodium nitrite and thiourea admit an interaction of $`P`$ with another mode $`u`$ (strain for example), Dvořák showed that the $`\eta `$-term accounts in an effective manner for this interaction, thereby reducing $`g(P,u)f(P)`$. By taking the Fourier transform of the above free energy density we find $$\stackrel{~}{f}=\left(\frac{\alpha }{2}+\frac{\kappa }{2}q^2+\frac{\lambda }{2}q^4\right)P_q^2+\left(\frac{\beta }{4}+\frac{\eta }{2}q^2\right)P_q^4.$$ (64) This justifies the choice of the $`E`$-term in the extended version of the DIFFOUR model, discussed in Section II. A term $`(\mathrm{d}P/\mathrm{d}z)^4`$ has been included in the free energy expansion by Jacobs et al., but by taking the Fourier transform one finds $`q^4P_q^4`$ which is of higher order than the $`\eta `$-term used here. In a seminal paper Hornreich et al. discussed a multicritical point of a new type, which they called a Lifshitz point. In the spherical model limit they were able to calculate critical exponents and the shape of the phase boundaries of $`2^{\mathrm{nd}}`$ order and $`1^{\mathrm{st}}`$ order transitions in the vicinity of the Lifshitz point. Let us return to the above free energy density to give a definition of the Lifshitz point. At an ordinary paraelectric to ferroelectric phase transition the coefficient $`\alpha `$ changes sign. If we have an additional incommensurate phase we need the $`\kappa `$ and $`\lambda `$-terms, and at the Lifshitz point $`\kappa =0`$. Higher-order terms in the expansion are needed for stabilization. Converting $`\alpha =0,\kappa =0`$ to variables in the extended DIFFOUR model we find $`A=0,C=4`$ (for $`D=1`$). This is exactly the position of the Lifshitz point found in Section V. There is another analogy between Landau theory and the DIFFOUR model: Michelson showed that for systems with uniaxial polarization the phase transition lines separating the paraelectric phase from the incommensurate phase and the incommensurate phase from the ferroelectric phase are tangent at the Lifshitz point. This feature is also present in figures 4 and 5. Let us now discuss some properties of the solutions found in Landau theory. Ground states minimize the total free energy $$F=\frac{1}{d}_0^df(z)dz,$$ (65) and can be found by solving the Euler-Lagrange equation $$\lambda \frac{\mathrm{d}^4P}{\mathrm{d}z^4}\kappa \frac{\mathrm{d}^2P}{\mathrm{d}z^2}\eta \left[P\left(\frac{\mathrm{d}P}{\mathrm{d}z}\right)^2+P^2\frac{\mathrm{d}^2P}{\mathrm{d}z^2}\right]+\alpha P+\beta P^3=0.$$ (66) Golovko was able to obtain exact solutions for some special values of the parameters in a slightly more general free energy density (he added a term $`\frac{\gamma }{6}P^6`$ to the expansion (63)). However, his method is not general and we will not discuss it further. Instead we follow a different approach: numerically solving equation (66) shows that the solutions contain practically only one harmonic, the amplitude of higher harmonics is at most 3.5% of the former. As usual in Landau theory only the coefficient $`\alpha `$ is temperature dependent: $`\alpha =\alpha _0(TT_c)`$. It is found that between the high-temperature paraelectric solution $$P(z)=0,F=0,$$ (67) and the low-temperature ferroelectric solution $$P(z)=\sqrt{\frac{\alpha }{\beta }},F=\frac{\alpha ^2}{4\beta },$$ (68) an incommensurate solution exists. Just below the paraelectric-incommensurate transition at $`\alpha _i=\alpha _0(T_iT_c)=\frac{\kappa ^2}{4\lambda }`$ it has the form $$P(z)=\rho _0\mathrm{cos}(qz),F=\frac{(\alpha _0\alpha )^2}{2\left(3\beta +2\eta q_0^2\right)}.$$ (69) The amplitude $`\rho _0`$ and wave vector $`q`$ are given by $`\rho _0^2`$ $`=`$ $`{\displaystyle \frac{4(\alpha _0\alpha )}{3\beta +2\eta q_0^2}},`$ (70) $`q`$ $`=`$ $`q_0\left(1+{\displaystyle \frac{\eta }{8\kappa }}\rho _0^2\right),`$ (71) $`q_0^2`$ $`=`$ $`{\displaystyle \frac{\kappa }{2\lambda }}.`$ (72) The $`\eta `$-term makes the incommensurate phase less stable when $`\eta `$ is positive, implying that the transition temperature from the incommensurate state to the ferroelectric state increases as $`\eta `$ increases. See also the discussion by Tolédano. In the discrete model a positive $`E`$-term is responsible for this effect. ## VIII Conclusions and outlook In this paper we have calculated various properties of an extension of the DIFFOUR model. For this purpose a next-nearest-neighbor generalization of the Effective Potential Method was developed. The shape of the paraelectric phase boundary was proven rigorously, elaborating on a former proof which only included nearest-neighbor interactions. We found that the phase diagram changes considerably due to the extra $`E`$-term, but the transition at the paraelectric phase boundary does not depend on $`E`$. Positive $`E`$ favors longer-period solutions. By taking thermal fluctuations in two different regimes into account the parameters $`A`$ and $`C`$ can be considered as effectively temperature dependent. For $`C`$ this holds only for nonzero $`E`$, which explains the relevance of this extra term. This has strong consequences for two experimentally easy accessible quantities: the temperature dependence of the modulation wave vector, and the ratio between the slopes of the soft phonon mode in the ferroelectric and paraelectric phases ($`R`$-parameter). Although lattice and continuum models have some features in common, the differences are more striking. A lattice model would be a more natural choice than the phenomenological Landau treatment of incommensurate phases. Discrete models do not need ad hoc lock-in terms to explain different commensurate and incommensurate phases. Complex phase diagrams can in principle be obtained using a simple Hamiltonian which takes into account the discreteness of a lattice. The Sn<sub>2</sub>P<sub>2</sub>(S<sub>1-x</sub>Se<sub>x</sub>)<sub>6</sub> crystal family seems to be an excellent system for our future research: it is uniaxial, has an exceptional Lifshitz point in the composition-temperature phase diagram, shows cross-over effects from order-disorder to a displacive type of phase transition, and displays an interesting modulation wave vector behavior. All these phenomena can in principle be explained by the extended DIFFOUR model. ###### Acknowledgements. We thank Stephan Eijt for drawing our attention to this problem and Alexey Rubtsov for providing us with the results of the Monte Carlo calculations. ## A Exact results for paraelectric and ferroelectric phases In this Appendix explicitly calculated phase boundaries are given for the extended DIFFOUR model. We first consider $`E=0`$ and then discuss the effect of $`E0`$. Let us start with writing $`V`$ in the form $$V=\underset{n}{}\left\{\frac{a}{2}x_n^2+\frac{1}{4}x_n^4+cx_nx_{n1}+dx_nx_{n2}\right\},$$ (A1) with $`d=\pm 1`$. The remaining parameters $`a,c,d`$ are the tilde parameters defined in (4) after normalization of $`\stackrel{~}{B}`$ and $`\stackrel{~}{D}`$. Let us now try to write this as $$V=\underset{n}{}\left\{p(x_nqx_{n1}rx_{n2})^2+\frac{1}{4}x_n^4\right\}.$$ (A2) Comparison of the two expressions yields $`p(1+q^2+r^2)`$ $`=`$ $`{\displaystyle \frac{a}{2}},`$ (A3) $`p(2q+2qr)`$ $`=`$ $`c,`$ (A4) $`2pr`$ $`=`$ $`d.`$ (A5) From this one can see that if it is possible to write the potential in this form and $`a`$ is positive, then $`p`$ is positive. Eliminating $`q`$ and $`r`$ from the above equations yields the following fourth order polynomial equation (assuming non-zero $`a,c`$ and $`d`$): $$16p^4+(16d8a)p^3+(8d^2+4c^28ad)p^2+(4d^32ad^2)p+d^4=0.$$ (A6) First consider the $`d=+1`$ case. Equation (A6) then has the following complex solutions $`p`$ $`=`$ $`{\displaystyle \frac{1}{8}}(2+a+\sqrt{(2+a)^24c^2}`$ (A8) $`\pm \sqrt{2}\sqrt{4+a^22c^2+(2+a)\sqrt{(2+a)^24c^2}}),`$ $`p`$ $`=`$ $`{\displaystyle \frac{1}{8}}(2+a\sqrt{(2+a)^24c^2}`$ (A10) $`\pm \sqrt{2}\sqrt{4+a^22c^2(2+a)\sqrt{(2+a)^24c^2}}).`$ The first requirement for having a real solution is that $`(2+a)^24c^20`$, i.e. $`a2+2\|c\|`$. Consider the first two solutions (A8). First look at $`\|c\|4`$. Then for $`a2+2\|c\|`$ we find $`(2+a)\sqrt{(2+a)^24c^2}>0`$. And for the first term in the root we find $`4+a^22c^2>2c^28\|c\|0`$. So for $`\|c\|4`$ the only requirement for having a real (positive, $`a>0`$) solution is $`a2+2\|c\|`$. For $`\|c\|<4`$ we have $`4+a^22c^2+(2+a)\sqrt{(2+a)^24c^2}=0`$ for $`a=2+\frac{1}{4}c^2`$. It can be seen that the argument of the root is positive for $`a>2+\frac{1}{4}c^2`$. So, for $`\|c\|<4`$ the requirement for having a real (positive, $`a>0`$) solution is: $`a2+\frac{1}{4}c^2`$ (then $`a>2+2\|c\|`$ automatically holds too). The requirements for $`\|c\|4`$ and $`\|c\|<4`$ form a continuous line in the $`ac`$-parameter space. Above this line $`V`$ can be written as (A2) with $`p`$ positive. Therefore $`V0`$. The lower bound is reached by the trivial solution which always exists, so the paraelectric phase is the ground state above this line. In Sec. IV it is shown that the above line corresponds exactly with the stability lines of the trivial solution, showing that the modulated phases arise from the destabilization of the trivial solution due to the condensation of a soft phonon mode. In case $`d=1`$ the solutions of the fourth order polynomial equation are $`p`$ $`=`$ $`{\displaystyle \frac{1}{8}}(2+a+\sqrt{(2+a)^24c^2}`$ (A12) $`\pm \sqrt{2}\sqrt{4+a^22c^2+(2+a)\sqrt{(2+a)^24c^2}}),`$ $`p`$ $`=`$ $`{\displaystyle \frac{1}{8}}(2+a\sqrt{(2+a)^24c^2}`$ (A14) $`\pm \sqrt{2}\sqrt{4+a^22c^2(2+a)\sqrt{(2+a)^24c^2}}).`$ Look at the two solutions (A12). The first condition is $`a2+2\|c\|`$. $`(2+a)`$ is positive if this requirement is fulfilled. Further $`4+a^22c^22c^2+8\|c\|0`$. So, here we have only one requirement for all $`c`$, namely $`a2+2\|c\|`$. Above this line the paraelectric phase is the ground state. With the same sort of reasoning we can also try to prove that for certain parameter values the ground state is ferroelectric Here we work with the following form of $`V`$ (with $`D=\pm 1`$): $$V=\underset{n}{}\left\{\frac{A}{2}x_n^2+\frac{1}{4}x_n^4+\frac{C}{2}(x_nx_{n1})^2+\frac{D}{2}(x_nx_{n2})^2\right\}.$$ (A15) We try to write this as $$V=\underset{n}{}\left\{\frac{A}{2}x_n^2+\frac{1}{4}x_n^4+P(x_nQx_{n1}Rx_{n2})^2\right\}.$$ (A16) Comparison of the two expressions yields $`P(1+Q^2+R^2)`$ $`=`$ $`C+D,`$ (A17) $`P(2Q+2QR)`$ $`=`$ $`C,`$ (A18) $`2PR`$ $`=`$ $`D.`$ (A19) If there exists a solution and $`C+D`$ is positive, then $`P`$ is positive. Rewriting the above equations yields the following fourth order polynomial equation (assuming nonzero $`C`$ and $`D`$): $$16P^4+(32D16C)P^3+(24D^2+16CD+4C^2)P^2+(8D^34CD^2)P+D^4=0.$$ (A20) For the case $`D=1`$ the complex solutions are $$P=\frac{1}{4}\left(2+C\pm \sqrt{C}\sqrt{4+C}\right),$$ (A21) both having multiplicity 2. For $`0<C<4`$ the solution is not real. For $`C4`$ the solution is real. In order to have a positive solution we must have $`C+D>0`$, so $`C>1`$. So, for $`C>4`$ the potential can be written as (A16) with $`P`$ positive. So $$V\underset{n}{}\left\{\frac{A}{2}x_n^2+\frac{1}{4}x_n^4\right\}.$$ (A22) The ferroelectric phase, which exists if $`A<0`$, reaches this lower bound. So for $`B=1,D=1,A<0`$ the ground state is ferroelectric for $`C>4`$. In terms of the tilde parameters: for $`\stackrel{~}{B}=1,\stackrel{~}{D}=1,\stackrel{~}{A}<22\stackrel{~}{C}`$, the ferroelectric phase is the ground state for $`\stackrel{~}{C}<4`$. For $`D=+1`$ the complex solutions are $$P=\frac{1}{4}\left(2+C\pm \sqrt{C}\sqrt{4+C}\right).$$ (A23) For $`4<C<0`$ the solution is not real. For $`C0`$ the solution is real. In that case $`P`$ is positive. For $`B=1,D=1,A<0`$ the ferroelectric phase is the ground state for $`C0`$. In other words: for $`\stackrel{~}{B}=1,\stackrel{~}{D}=1,\stackrel{~}{A}<22\stackrel{~}{C}`$ the ferroelectric phase is the ground state for $`\stackrel{~}{C}0`$. In terms of the tilde parameters analogous statements about the anti-ferroelectric phase can easily be made ($`\stackrel{~}{C}\stackrel{~}{C}`$). In case $`E0`$ the following holds: the parts of the phase diagram where the paraelectric phase is the ground state in the DIFFOUR model $`(E=0)`$ also belong to the paraelectric phase for this extended model for all allowed values of $`E`$. In this extended model there are no other parts of the phase diagram where the trivial solution is the ground state, because the stability conditions for this solution are the same as in the DIFFOUR model (see Section IV). For the ferroelectric phase the statements are less rigorous: if the ferroelectric phase is the ground state in the DIFFOUR model $`(E=0)`$, then it is also the ground state in the extended model for $`E>0`$. Again we can prove that $`V_n\left\{\frac{A}{2}x_n^2+\frac{B}{4}x_n^4\right\}`$. However, for positive $`E`$ the part of the phase diagram where the ferroelectric phase is the ground state becomes bigger. ## B Effective Potential Method for next-nearest neighbors In this Appendix, which is based on the account given by Griffiths, we discuss how the EPM can be generalized to be applicable to systems in which there is next-nearest neighbor interaction. Consider a classical one-dimensional chain of atoms. The total potential energy of the system is given by $$H=\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}\left\{V(x_n)+W(x_{n+1},x_n)+D(x_{n+2},x_n)\right\}.$$ (B1) So, interactions up to next-nearest-neighbors are included. The effective potential method that is used to find ground states for systems with interaction with first neighbors, can also be used for systems where second neighbor interaction is included. Instead of a scalar variable at site $`n`$, one now has to deal with a vector consisting of the values for $`x`$ for two adjacent atoms. Writing $`𝐱_n=(x_{2n},x_{2n1})`$ the above potential energy is of the form $$H=\underset{n}{}K(𝐱_n,𝐱_{n1}),$$ (B2) with $`K(𝐱_n,𝐱_{n1})`$ $`=`$ $`V(x_{2n})+V(x_{2n1})+W(x_{2n},x_{2n1})+W(x_{2n1},x_{2n2})`$ (B4) $`+D(x_{2n},x_{2n2})+D(x_{2n1},x_{2n3}).`$ The fact that here a vector at site $`n`$ is considered does not change the EPM and the proofs given for this method. The solution can be found by solving $$\eta +R(𝐱_n)=\underset{𝐱_{n1}}{\mathrm{min}}\left\{R(𝐱_{n1})+K(𝐱_n,𝐱_{n1})\right\}.$$ (B5) Here $`\eta `$ is 2 times the ground state energy per particle. The vector consists of two components with respect to which one has to minimize. Because of this minimization over two components, which has to be performed frequently, the numerical procedures based on discretization of the range of possible $`x_n`$-values will take a very long time. The method we used is slightly different. Consider the $`n^{\mathrm{th}}`$ couple of adjacent atoms. Couple $`(n1)`$ does not consist of two other atoms as is the case above. Instead, the right atom of couple $`(n1)`$ is the same as the left atom of couple $`n`$. In this case only a minimization over one atomic degree of freedom (the ‘position’) is left. This leads to more reasonable computation times. The fact that couple $`(n1)`$ is not independent of couple $`n`$ requires an adaptation of the proof of the existence of a solution both in the continuous case and in the discretized case used for numerical procedures. We only give an outline for the method, proofs of the existence of solutions and generalization to systems with interactions up to $`s^{\mathrm{th}}`$ neighbors are to be given in a separate paper. The method for deriving the equations and the numerical procedures for solving the equations remain essentially the same. Numerical calculations suggest that the error in $`\eta `$ has a cubic dependence on the grid size rather than a quadratic dependence for the Frenkel-Kontorova model. Let us give the following explanation for the method: Imagine that a system described by (B1) is in its ground state. If we now change the positions of two adjacent atoms, the surrounding atoms will in general also change their positions in order to minimize the total energy. This net energy change caused by the deformation of one couple will be called the effective two-particle-potential. This will describe the energy cost as a function of the positions of two adjacent atoms. At site $`n`$, the effective two-particle-potential $`R(x_{n+1},x_n)`$, due to the presence of the atoms $`i<n`$, can be formally written as $$R(x_{n+1},x_n)\underset{i<n}{\mathrm{min}}\left\{\underset{in+1}{}\left[V(x_i)+W(x_i,x_{i1})+D(x_i,x_{i2})\eta \right]\right\},$$ (B6) where the minimum is taken over all atomic positions $`x_i`$ with $`i<n`$ and $`\eta `$ is the (unknown) ground state energy per particle. By rewriting this equation, one obtains $`R(x_{n+1},x_n)=\underset{x_{n1}}{\mathrm{min}}\underset{i<n1}{\mathrm{min}}\{`$ $`{\displaystyle \underset{in}{}}\left[V(x_i)+W(x_i,x_{i1})+D(x_i,x_{i2})\eta \right]`$ (B8) $`+V(x_{n+1})+W(x_{n+1},x_n)+D(x_{n+1},x_{n1})\eta \},`$ which gives $$\eta +R(x_{n+1},x_n)=V(x_{n+1})+W(x_{n+1},x_n)+\underset{x_{n1}}{\mathrm{min}}\left[R(x_n,x_{n1})+D(x_{n+1},x_{n1})\right].$$ (B9) This is the minimization eigenvalue equation for $`R`$. The same procedure can be followed for the effect of the atoms $`i>n+1`$. The effective two-particle-potential due to these atoms is called $`S(x_{n+1},x_n)`$, which gives $$\eta +S(x_{n+1},x_n)=V(x_n)+W(x_{n+1},x_n)+\underset{x_{n+2}}{\mathrm{min}}\left[S(x_{n+2},x_{n+1})+D(x_{n+2},x_n)\right].$$ (B10) The total effective two-particle-potential $`F(x_{n+1},x_n)`$ of a couple of adjacent atoms in a double infinite chain, is given by $$F(x_{n+1},x_n)=R(x_{n+1},x_n)+S(x_{n+1},x_n)V(x_{n+1})V(x_n)W(x_{n+1},x_n),$$ (B11) where the last three terms are subtracted on the right side to avoid double counting. The equations (B9) and (B10) can also be obtained in another way. Let $`\stackrel{~}{R}_N(x_{n+1},x_n)`$ be the minimal energy of a chain of $`N`$ atoms with the constraint that the atoms $`N`$ and $`N1`$ are at fixed positions $`x_{n+1}`$ and $`x_n`$ respectively, while the other atoms are free to rearrange themselves in an optimal way so as to minimize the total energy. This leads to $`\stackrel{~}{R}_2(x_{n+1},x_n)`$ $`=`$ $`V(x_{n+1})+V(x_n)+W(x_{n+1},x_n),`$ (B12) $`\stackrel{~}{R}_3(x_{n+1},x_n)`$ $`=`$ $`V(x_{n+1})+V(x_n)+W(x_{n+1},x_n)`$ (B14) $`+\underset{x_{n1}}{\mathrm{min}}\left[V(x_{n1})+W(x_n,x_{n1})+D(x_{n+1},x_{n1})\right]`$ $`=`$ $`V(x_{n+1})+W(x_{n+1},x_n)+\underset{x_{n1}}{\mathrm{min}}\left[\stackrel{~}{R}_2(x_n,x_{n1})+D(x_{n+1},x_{n1})\right],`$ (B15) $`\stackrel{~}{R}_{N+1}(x_{n+1},x_n)`$ $`=`$ $`V(x_{n+1})+W(x_{n+1},x_n)+\underset{x_{n1}}{\mathrm{min}}\left[\stackrel{~}{R}_N(x_n,x_{n1})+D(x_{n+1},x_{n1})\right].`$ (B16) Now assume that for $`N\mathrm{}`$, $`\stackrel{~}{R}_N(x_{n+1},x_n)`$ approaches some function $`R(x_{n+1},x_n)`$ plus a constant proportional to $`N2`$: $$\stackrel{~}{R}_N(x_{n+1},x_n)R(x_{n+1},x_n)+(N2)\eta .$$ (B17) In that case equation (B9) follows. However, it is not clear that (B17) will always be satisfied. But by imposing a special boundary condition, namely $$\stackrel{~}{R}_2(x_{n+1},x_n)=R(x_{n+1},x_n),$$ (B18) (B17) will be satisfied exactly. The previous boundary condition is the same as saying that the left-most couple experiences the effective two-particle-potential instead of the true two-particle-potential. The minimum energy of this system as a function of the positions of the two right-most atoms is given by $`R+N\eta `$. Assuming that $`R`$ is a bounded function, the energy per particle of such a system will tend to $`\eta `$ as $`N\mathrm{}`$. $`\eta `$ is thus the average energy per particle in any ground state, since the extra boundary condition only changes the total energy by a term of order 1. So $`R`$ is the effective two-particle-potential for the right-most couple of a semi-infinite chain. The same is true for $`S`$ for the left-most couple of a semi-infinite chain extending to the right. $`F`$ is the total effective two-particle-potential for a couple in a double-infinite chain. $`R`$,$`S`$ and $`F`$ can of course only be defined up to an additive constant. In the above derivations the problems arising from the summation of an infinite number of terms in (B6) are neglected. In fact, one considers local deformations of length $`M`$, with the limit $`M\mathrm{}`$. This will be explained below. Define the effective two-particle-potential due to the local deformation of length $`M`$ as $`R^{(M)}(x_{n+1},x_n)\underset{n+1M<i<n}{\mathrm{min}}\{`$ $`{\displaystyle \underset{n+3M<in+1}{}}\left[K(x_i,x_{i1},x_{i2})\eta \right]`$ (B21) $`+[K(x_{n+3M},x_{n+2M},u_{n+1M})`$ $`+K(x_{n+2M},u_{n+1M},u_{nM})2\eta ]\},`$ where $`u_i`$ refers to the ground state value for atom $`i`$, and where we have introduced $$K(x_{n+1},x_n,x_{n1})V(x_{n+1})+W(x_{n+1},x_n)+D(x_{n+1},x_{n1}).$$ (B22) The right hand site of equation (B21) can be rewritten as $$R^{(M)}(x_{n+1},x_n)=\underset{x_{n1}}{\mathrm{min}}\left[R^{(M1)}(x_n,x_{n1})+K(x_{n+1},x_n,x_{n1})\eta \right].$$ (B23) It is reasonable to assume that in the limit $`M\mathrm{}:R^{(M)}(x_{n+1},x_n)R^{(M1)}(x_{n+1},x_n)`$ (because $`x_{n+2M}u_{n+2M}`$). Writing $`R(x_{n+1},x_n)=lim_M\mathrm{}R^{(M)}(x_{n+1},x_n)`$ the minimization eigenvalue equation results. In the second version of obtaining the equations it is clear that it is in fact the limit of local deformations, however with the boundary condition that the left most couple of atoms experiences the effective two-particle-potential. Here, one should take the length of the chain going to infinity in order to let $`\eta `$ go to the ground state energy per particle. The above explanation also holds for $`S`$. It is best to picture the situation as a local deformation of the ground state. Now, the nonlinear minimization eigenvalue equations for $`R`$ and $`S`$ are rewritten. The eigenvalue equation for $`R`$ now becomes $$\eta +R(x_{n+1},x_n)=\underset{x_{n1}}{\mathrm{min}}\left[R(x_n,x_{n1})+K(x_{n+1},x_n,x_{n1})\right].$$ (B24) Let the function $`L`$ be defined by $$L(x_{n+1},x_n)=S(x_{n+1},x_n)V(x_{n+1})V(x_n)W(x_{n+1},x_n).$$ (B25) The minimization eigenvalue equation for $`S`$ can now be rewritten as $$\eta +L(x_{n+1},x_n)=\underset{x_{n+2}}{\mathrm{min}}\left[L(x_{n+2},x_{n+1})+K(x_{n+2},x_{n+1},x_n)\right].$$ (B26) In terms of $`R`$ and $`L`$ one has $$F(x_{n+1},x_n)=R(x_{n+1},x_n)+L(x_{n+1},x_n).$$ (B27) In fact there may be multiple solutions of the eigenvalue equation, not only differing by a trivial constant. The existence of different solutions is related to the existence of different degenerate ground states. The general solution is given by $$R(x_{n+1},x_n)=\underset{\alpha }{\mathrm{min}}\left[R_\alpha (x_{n+1},x_n)+K_\alpha \right].$$ (B28) The $`R_\alpha `$ correspond to the pure phases and the $`K_\alpha `$ are arbitrary constants. For each solution of the minimization eigenvalue equation for $`R`$ (B24), a $`\tau `$ map can be defined, where $`\tau (x_{n+1},x_n)=\{(x_n,x_{n1})\}`$ with $`x_{n1}`$ one of the values for which the minimum on the right hand side of (B24) is achieved. An $`R`$-orbit is defined as $$n:(x_n,x_{n1})\tau (x_{n+1},x_n)\eta +R(x_{n+1},x_n)=R(x_n,x_{n1})+K(x_{n+1},x_n,x_{n1}).$$ (B29) Similarly, for the minimization eigenvalue equation for $`L`$ (B26), a $`\sigma `$ map can be defined, where $`\sigma (x_{n+1},x_n)=\{(x_{n+2},x_{n+1})\}`$ with $`x_{n+2}`$ one of the values for which the minimum on the right hand side of (B26) is achieved. An $`L`$-orbit is defined as $$n:(x_{n+2},x_{n+1})\sigma (x_{n+1},x_n)\eta +L(x_{n+1},x_n)=L(x_{n+2},x_{n+1})+K(x_{n+2},x_{n+1},x_n).$$ (B30) A ground state is both an $`R`$-orbit and an $`L`$-orbit. Therefore, it can be proven that for a ground state $$F(x_{n+1},x_n)=R(x_{n+1},x_n)+L(x_{n+1},x_n)=F(x_n,x_{n1})$$ (B31) So, $`F`$ is constant on the positions of two adjacent atoms in a ground state, which is logical since it is the effective two-particle-potential. Numerical procedures are based on a discretized version of the system. In that case for each ground state there is a solution for the eigenvalue equation for which there is a path from each point to the ground state in the corresponding $`\tau `$ graph. So, the situation is as follows. There is a local deformation of length $`M`$ (with $`M\mathrm{}`$), in a chain coinciding with a particular ground state for $`\pm \mathrm{}`$. This ground state corresponds to a certain solution of the eigenvalue equation. The deformation is such that atoms $`n`$ and $`n+1`$ have values $`x_n`$ and $`x_{n+1}`$. By applying the corresponding $`\tau `$ map one can obtain the positions of the atoms left from $`n`$. For the discretized system one will finally reach the ground state in this way (In the continuous case it is supposed to converge to the ground state). The same can be done for the atoms right from $`n+1`$ by applying the $`\sigma `$ map. From this picture it is clear that the ground state is both an $`R`$\- orbit and an $`L`$-orbit. In fact the $`\tau `$ map and the $`\sigma `$ map may be multi-valued. So, the atomic positions of the atoms (say) left from $`n`$ do not have to be unique. The deformation can have parts consisting of minimizing cycles (cycles of minimal energy) different from the ground state configuration at $`\pm \mathrm{}`$. In the $`\tau `$ graph one can go directly to the minimizing cycle corresponding to the ground state configuration at $`\mathrm{}`$, or one can first stay for some time in another minimizing cycle if this exists. If there are several solutions for $`R`$ and $`S`$ (with several corresponding $`\tau `$ and $`\sigma `$ maps), there are several possibilities to construct $`F`$. It will often be logical to take the ground states, toward which the chain converges at $`\mathrm{}`$ and $`+\mathrm{}`$, the same. The numerical algorith we used is an adapted version of the one discussed by Floria and Griffiths. Here an example will be given to show that it is important that in fact limits of finite deformations are considered. Suppose that the ground state is ferroelectric, with two degenerate ground states: $`x_n=x=\pm l`$ where $`l0`$. If the deformation is just infinite as suggested in (B6) the value of $`R(l,l)`$ should be the same as the value for $`R(l,l)`$. However, something else is seen. Two solutions can be found corresponding to the two ground states. In the solution corresponding to the solution $`x_n=l`$, $`R(l,l)`$ has a higher value than $`R(l,l)`$. The difference is the defect energy, the energy cost for going from the $`+l`$ phase to the $`l`$ phase. (The defect energy (and the defect configuration) can also be calculated using the $`\tau `$ map.) From this it can be seen that the deformation is in fact embedded in the ground state $`u_i=+l`$ at $`\mathrm{}`$ (In the limit $`M\mathrm{}`$: $`x_{n+2M}u_{n+2M}`$ where $`u_{n+2M}=+l`$, or for the second version: the left most couple of atoms in the finite chain experiences the effective two-particle-potential corresponding to the ground state $`u_i=+l`$). It can also be expected that $`F`$ has local minima at the positions of two adjacent atoms in metastable states. However, since only two atoms are at a fixed position, while the other atoms are free to rearrange themselves in an optimal way, this may not be the case. When changing the positions of two adjacent atoms by an infinitesimal amount, the changes of the other atomic positions in the metastable state does not have to be infinitesimal. Therefore, it is not necessarily true that there is a local minimum in $`F`$ for positions of two adjacent atoms in a metastable state. When there are no other atomic positions in the metastable state (a period 1 solution), $`F`$ does have a local minimum. When the lowest metastable state has positions of two adjacent atoms which are not seen in a ground state (which will often be the case), there will be a local minimum in $`F`$ for these two positions. In that case the energy cannot be lowered by changing the other atoms by any amount, since the only states which have lower energy are ground states and these cannot be reached since the positions of the two adjacent atoms in consideration are not in a ground state (and the changes of them should be infinitesimal). By following the development of the shape of $`F`$ one may also investigate the kind of phase transitions that are involved. For example, a discontinuous change in the set of points where $`F`$ achieves its global minimum, indicates a first order transition.
warning/0004/cond-mat0004363.html	ar5iv	text	# X-ray scattering study of two length scales in the critical fluctuations of CuGeO3 ## I Introduction High resolution x-ray and neutron scattering studies of the critical fluctuations associated with structural and magnetic phase transitions typically reveal “two length scales”, that is, two distinctive scattering lineshapes superimposed upon each other in the critical scattering profile. Since the existence of the second length scale seems to contradict the fundamental assumption of modern critical phenomena theory that there exists only one characteristic length in the critical fluctuations, extensive experimental and theoretical efforts have been devoted to elucidating the exact origin of this phenomenon. However, in spite of a significant amount of work dedicated to this problem, a consensus still has not been reached. Presently, there exist two main approaches:(1) models based on intrinsic near-surface effects and (2) explanations involving near-surface random defects. The accumulating experimental evidence seems to favor the randomness interpretation although there is still no definitive experiment to pinpoint the exact origin of the second length scale fluctuations. In the present paper, we present a high-resolution synchrotron x-ray scattering study of the critical fluctuations associated with the spin-Peierls structual phase transition in CuGeO<sub>3</sub>. Not only do we clearly observe two lineshapes in the critical scattering profile, but we also observe a dramatic change of the anisotropy ratio of the correlation length divergence along the three primary crystal axes. The existence of the modified anisotropy ratio provides substantial evidence that near-surface dislocations are the origin of the second length scale fluctuations. Our paper is organized as follows: In Section II we provide details of the sample preparation and experimental measurements. In Section III we present our experimental results. A discussion of the results and conclusions is given in Section IV. ## II Experimental Procedures The experiment was carried out at MIT-IBM beamline X20A at the National Synchrotron Light Source. The x-ray beam was focused by a mirror, monochromatized by a pair of Ge (111) crystals, scattered from the sample, and analyzed by a Si (111) analyzer. The x-ray energy was 8.5 keV. High quality pure CuGeO<sub>3</sub> and Cu<sub>0.99</sub>Zn<sub>0.01</sub>GeO<sub>3</sub> single crystals grown by the travelling solvent floating zone method were used. Carefully cleaved samples were placed inside a Be can filled with helium heat-exchange gas and mounted on the cold finger of a 4K closed cycle cryostat. The experiment was carried out around the (1.5, 1, 1.5) SP dimerization peak position with the (H K H) zone in the scattering plane. ## III Experimental Results Pretransitional lattice fluctuations along the H, K and L directions have been measured in pure CuGeO<sub>3</sub> by x-ray and neutron scattering. All of the experiments have shown rapid and anisotropic broadening of the scattering peaks when the sample was heated across T<sub>SP</sub>, which was clear evidence of anisotropy in the magnetic interaction. Close to T<sub>SP</sub>, however, Schoeffel et al. observed a crossover temperature T<sub>CO</sub> where the ratio of the correlation lengths along the three crystal axis directions appeared to change abruptly: $`\xi _c/\xi _a4`$ and $`\xi _c/\xi _b1`$ below $`T_{CO}`$ and $`\xi _c/\xi _b1.6`$ above. This was used as evidence of a crossover to a 2D lattice fluctuation regime above $`T_{SP}`$. Later, both experimental and theoretical efforts were devoted to elucidating the exact nature of the 2D crossover. Harris et al., on the other hand, studied the critical behavior in the immediate vicinity above T<sub>SP</sub> and reported a different anisotropy ratio. However, the critical fluctuations reported by Harris et al. have length scales which are about an order of magnitude larger than those reported by Schoeffel et al.. The discrepancies in these two experiments demonstrate that one must treat the data near the transition more cautiously. In extracting the correlation length just above T<sub>SP</sub>, it is necessary to take into account explicitly that there exist two distinct scattering length scales. Distinguishing and separating their individual contributions to the total cross-section will be of primary importance. This comprises a principal motivation of this experiment. To reconcile the results of previous critical scattering studies of CuGeO<sub>3</sub> and to obtain some insight into the physical origin of the second length scale, we carefully studied the pretransitional critical behavior just above T<sub>SP</sub>. Though the exact origin of the long length scale fluctuations has not been determined, it has long been speculated that they originate from random surface stresses caused by defects. Hence, we prepared our samples by cleaving them several times until no observable cracks could be seen by visual inspection. In doing so, we took advantage of the fact that CuGeO<sub>3</sub> crystals are inclined to self-cleave along the $`a`$ crystal plane. Thus, no additional grinding or polishing process is necessary to achieve a visually smooth mirror surface. To put the two previous seemingly conflicting experiments together, we need to have information on both length scales in the same sample. Fortunately, this is exactly what we have observed in our experiment. In Fig. 1 we show the critical scattering profiles along the H, K, and L directions at T$`{}_{SP}{}^{}+0.1`$K and T$`{}_{SP}{}^{}+0.3`$K for undoped CuGeO<sub>3</sub>. At T$`{}_{SP}{}^{}+0.1`$K, there clearly exist two distinct scattering profiles along all three directions, with a sharp central peak superimposed upon a broader peak. This corresponds to archetypal two-length scale behavior. However, a closer examination of the data reveals that even though there are clearly two features along all three directions, the central peak along the H direction is much sharper in comparison to the broad one than is observed along the other two directions. In other words, the ratio of the correlation lengths for these two length scales are significantly different along one of the three crystal axis directions. In Fig. 2 , we show the inverse correlation lengths of the broad component as functions of temperature along the H, K, and L directions. Several features can be recognized immediately. First, the correlation length diverges rapidly as the temperature approaches T<sub>SP</sub> from above, which demonstrates that the SP transition in our CuGeO<sub>3</sub> crystal is a well defined second-order phase transition. Second, the correlation length also diverges anisotropically along the three crystal axes. In the temperature range $`T_{SP}<T<T_{SP}+0.4K`$, the anisotropy ratio remains $`\xi _c/\xi _a5`$ and $`\xi _c/\xi _b3`$, which is consistent with the high temperature data taken both by x-ray and neutron scattering. Thus we do not observe any evidence for the presumed 2D crossover , in which there should exist a dramatic change in the anisotropy ratio about 1K above T<sub>SP</sub>. Specifically, in Ref it is argued that below this crossover temperature, the correlation length along the $`b`$-axis direction equals the correlation length along $`c`$-axis direction. Our results clearly demonstrate that the correlation length anisotropy ratio remains unchanged from high temperatures to very near T<sub>SP</sub>, that is, there is no evidence for any crossover. Fig. 3 shows the inverse correlation length of the sharp component as a function of temperature. One of the salient features is that, although the correlation length of the sharp component diverges in a manner similar to that of the broader component, the anisotropy ratio of the correlation lengths along the three axes directions is modified to $`\xi _c/\xi _a1.5`$ and $`\xi _c/\xi _b4.4`$. This is reminiscent of the high resolution results reported by Harris et al.. Instead of the relationship $`\xi _c>\xi _b>\xi _a`$, the large length scale fluctuations exhibit the hierarchy $`\xi _c>\xi _a>\xi _b`$. The change of the order of the correlation lengths is informative, since there are not many physical mechanisms which could induce such a directional preference. The confirmation of the change of the order of the correlation lengths in our experiment proves that this is a general phenomenon instead of an irreproducible singular case. Furthermore, if we directly compare the magnitude of the two length scales along $`a`$, $`b`$ and $`c`$ crystal axes, ratios of $`28:6:8`$ would result, with the maximum along the $`a`$ axis and similar values along the $`b`$\- and $`c`$-axes. The other feature worth mentioning is the relative importance of the second length scale in both studies. In Harris et. al. ’s case, only the long length scale fluctuations were clearly observable over the temperature range studied. On the other hand, in our experiments, the fluctuations associated with both length scales are clearly observable, which proves that the relative amplitude of the second length scale fluctuations is sample dependent. Over the last several years, we and others have carried out detailed studies of the effects of dopants on the CuGeO<sub>3</sub> magnetic and structural phase transitions with a focus on the overall phase diagram. Such studies can be regarded as a systematic exploration of the effects of point defects on the CuGeO<sub>3</sub> structural phase transition. Thus, as a byproduct of our Cu<sub>1-x</sub>(Zn,Mg)<sub>x</sub>GeO<sub>3</sub> phase diagram studies, we also are able to test the hypothesis that the long length scale fluctuations are caused by point defects. Fig. 4 shows the inverse correlation lengths along the three crystal axes as functions of temperature for 1% Zn-doped CuGeO<sub>3</sub>. The dramatic effects of the Cu ion dilution on the phase transition are apparent: the transition temperature has been suppressed by more than 1K upon only 1% Zn doping, and the critical exponent associated with the correlation length appears to be different from that of the undoped sample. We are uncertain currently whether this apparently different critical behavior is intrinsic or merely due to a trivial concentration gradient effect. Further experiments are needed to clarify this issue. However, the ratios of the inverse correlation lengths are $`\xi _c:\xi _b:\xi _a=5.9:2.0:1`$, which are essentially identical to those of the undoped samples. This consistency of the anisotropy ratios between the doped and undoped sample naturally excludes models for the second length scale based on point defects. ## IV Discussion Before we present our interpretion of our experimental observations, we first briefly summarize the results of previous experimental and theoretical studies on the two length scale phenomenon. Most high resolution critical scattering studies of both structural and magnetic phase transitions reveal two length scales. Further, an elegant neutron scattering study by Shirane and coworkers reveals that the long length scale fluctuations are located in the “skin” of the sample. A subsequent study on the same single crystals by transmission electron microsopy\[TEM\] finds that the density of dislocations has a steep increase within a few microns of the sample surface, which coincides with the onset of the long length scale. Based on the spatial coexistence of the second length scale and dislocations, the authors of Ref conclude that the second length scale originates from dislocations, albeit in an indirect way. As more and more experimental evidence turns up, a gradual consensus is emerging that the origin of the second length scale fluctuations is the random strain fields caused by defects in the sample skins. However, an intrinsic effect explanation can not be excluded. Moreover, even if the idea that the second length scale originates from defects is taken for granted, there exists additional complexity because the defects can either be point defects or line defects such as dislocations. A recent study suggests that point defects are responsible for the occurrence of the second length scale. The dislocation theory, on the other hand, has been less favored. One of the key objections used against it is the lack of directional preference in all the previous studies, that is, dislocations are line defects and they should inevitably favor particular directions. From the results of our study, we believe that CuGeO<sub>3</sub> serves as a model system to study the origin of the second length scale and provides strong evidence that dislocation defects are responsible for the occurence of the second length scale fluctuations. Using dislocation theory, in the following, we explain our experimental observations by a phenomenological model. One of the marked differences between our results and those reported by Harris is the relative importance of the long length scale fluctuations. This can easily be explained, since the density and spatial distribution of dislocations naturally depend on sample preparation and surface processing such as chemical etching, so they would unavoidably vary from sample to sample. The most determinant piece of information to support a dislocation model is the occurrence of a directional preference. The experimental results on the Zn-doped sample provide additional support by demonstrating the irrelevancy of point defects. To understand qualitatively the experimentally observed direction preference, we refer to the theoretical work by Altarelli et al., in which the effect of dislocations has been treated on a qualitative level. As discussed by Altarelli et al., in real crystals, surface treatment always induces slipping parallel to the surface. These defects are edge dislocations parallel to the sample surface but randomly oriented in the plane. They induce anisotropic stress fields in the surrounding crystal since they are line defects by nature. The stress field produced by dislocations can be well modeled by dipole fields with the maximum in the plane perpendicular to the dipole, which is the Burgers vector direction in our case. The whole problem can then be mapped into that of a group of randomly oriented dipoles lying in a plane. The stress field can lower the free energy of the structually ordered phase, thus increasing the phase transition temperature in the stressed region. This is used to account for the emergence of the second length scale. Using this theory, the different ratio between the two length scales can be qualitatively explained. We recall that CuGeO<sub>3</sub> crystals naturally cleave in the $`a`$ plane.The difference in magnitude of the two length scales is most prominent along the $`a`$-axis because the fluctuation amplitude is presumed to be proportional to the average stress field. The random orientation of the dipoles in the surface would results in an isotropic stress field distribution in the plane. However, the maximum average stress field would be produced along the surface normal direction due to the dipole nature. We believe that the stress field is responsible for the creation of pretransitional ordered domain structures. These domains order at a higher temperature than the bulk and have an anisotropic structure owing to the anisotropy of the stress field. This can naturally explain the unusual sharp feature of the critical scattering along $`a`$-axis and also why the ratio of the two length scales remains relatively unmodified in the other two directions. We find a tiny difference in the ratio along the $`b`$ and $`c`$ directions, in agreement with Harris et al.. We speculate that this subtle anisotropy originates in a slight anisotropic distribution of the dislocations in the plane. Indeed, a closer inspection of a naturally cleaved CuGeO<sub>3</sub> sample surface reveals that the apparently smooth surface is actually composed of some stripes running along the $`c`$-axis direction. These are most likely formed during crystal growth. When the single crystals are grown using the floating zone method, the seed rod is oriented with the easy growth direction coinciding with the travelling zone direction. Stripes are then naturally formed along the direction of the crystal growth, which is $`c`$-axis. From the theory of dislocations, structural line defects are preferentially created along the same direction. These defects are normally edge dislocations with the Burgers’ vector perpendicular to the dislocation line and lying in the slip plane, the $`b`$-axis direction in our case. Hence , this could create a tiny preference for the dipoles to lie in the $`b`$ direction: a resulting minimum ratio along the $`b`$-axis is expected. We would further comment that even though dislocation theory offers a satisfactory heuristic explanation of our critical scattering results, many open questions still exist, for example, why is there a similar ratio between these two length scales in many different physical systems and why there does there exist a clear phase transition for the long length scale fluctuations despite the fact we are assuming a spread of transition temperatures. More theoretical and experimental work is needed to address these issues. In conclusion, we have studied the critical fluctuations in pure and Zn doped CuGeO<sub>3</sub>. Two length scales have been observed with different anisotropy ratios for the correlation lengths along the three crystal axes. The maximum of the magnitude of the two length scales is found to be along the $`a`$-axis direction, which is the surface normal of the crystal. We argue that dislocation theory would serve as the best explanation of the origin of the second length scale fluctuations. ## V Acknowledgment We thank G. Shirane for insightful comments. This work was supported by the NSF-LTP Program under Grant No. DMR97-04532.
warning/0004/gr-qc0004031.html	ar5iv	text	# Exact radiative spacetimes: some recent developments ## 1 Introduction Gernot Franz Sebastian Neugebauer suggested that I should give a fairly broad review on radiative spacetimes. Johann Wolfgang von Goethe believes that this is not easy: “It is extremely difficult to report on the opinions of others … If the reporter goes into detail, he creates impatience and boredom; if he wants to summarize, he risks giving his own point of view; if he avoids judgements, the reader does not know where to begin, and if he organizes his materials according to principles, the presentation becomes one-sided and arouses opposition, and the history itself creates new histories.” (J.W.Goethe, “Materialien zur Geschichte der Farbenlehre”.) My “Materialien to the Exact Gravitational-Waves Lehre” will be: plane waves and their collisions; cylindrical waves and null infinity in (2+1)-dimensional spacetimes; Robinson-Trautman solutions and type N twisting spacetimes; boost-rotation symmetric spacetimes and spinning C-metric; “cosmological” waves (Gowdy models) and the approach to cosmological singularity. Some parts of the “Materialien” are taken from ; see also for the more detailed review and references until 1995. ## 2 Plane waves and their collisions By the definition (see e.g. ) a vacuum spacetime is a “plane-fronted gravitational wave” if it contains a shearfree geodesic null congruence (with tangents $`k^\alpha `$), and if it admits “plane wave surfaces” (spacelike 2-surfaces orthogonal to $`k^\alpha `$). Because of the existence of plane wave surfaces, the expansion and twist must vanish as well. The best known subclass of these waves are “plane-fronted gravitational waves with parallel rays” (pp-waves) which are defined by the condition that the null vector $`k^\alpha `$ is covariantly constant, $`k_{\alpha ;\beta }=0`$. In suitable null coordinates with a null coordinate $`u`$ such that $`k_\alpha =u_{,\alpha }`$ and $`k^\alpha =\left(/v\right)^\alpha `$, the metric has the form $$ds^2=2d\zeta d\overline{\zeta }2dudv2H(u,\zeta ,\overline{\zeta })du^2,$$ (1) where $`H`$ is a real function dependent on $`u`$, and on the complex coordinate $`\zeta `$ which spans the wave 2-surfaces $`u=\mathrm{constant}`$, $`v=\mathrm{constant}`$. The vacuum field equations imply $`2H=f(u,\zeta )+\overline{f}(u,\overline{\zeta })`$, where $`f`$ is an arbitrary function of $`u`$, analytic in $`\zeta `$. In general the pp-waves have only the single isometry generated by the Killing vector $`k^\alpha =\left(/v\right)^\alpha `$. However, a much larger group of symmetries may exist for various particular choices of the function $`H(u,\zeta ,\overline{\zeta })`$. Jordan, Ehlers and Kundt (see also ) gave a complete classification of the pp-waves in terms of their symmetries and corresponding special forms of $`H`$. For example, in the best known case of plane waves $`H(u,\zeta ,\overline{\zeta })=A(u)\zeta ^2+\overline{A}(u)\overline{\zeta }^2`$, with $`A(u)`$ being an arbitrary function of $`u`$. This spacetime admits five Killing vectors. Recently, Aichelburg and Balasin generalized the classification given in by admitting distribution-valued profile functions and allowing for non-vacuum spacetimes. They have shown that with $`H`$ in the form of delta-like pulses, $`H(u,\zeta ,\overline{\zeta })=f(\zeta ,\overline{\zeta })\delta (u)`$, new symmetry classes arise even in the vacuum case. The main motivation to consider impulsive pp-waves stems from the metrics describing a black hole or a “particle” boosted to the speed of light. The simplest metric of this type, given by Aichelburg and Sexl , is a Schwarzschild black hole with mass $`m`$ boosted in such a way that $`\mu =m/\sqrt{1w^2}`$ is held constant as $`w1`$. It reads $$ds^2=2d\zeta d\overline{\zeta }2dudv4\mu \mathrm{log}(\zeta \overline{\zeta })\delta (u)du^2,$$ (2) with $`H`$ clearly in the form of a delta pulse. This is not a vacuum metric: the energy-momentum tensor $`T_{\alpha \beta }=\mu \delta (u)\delta (\zeta )k_\alpha k_\beta `$ indicates that there is a “point-like particle” moving with the speed of light along $`u=0`$. The interest in impulsive waves generated by boosting a “particle” at rest to the velocity of light by means of an appropriate limiting procedure persists up to the present. The ultrarelativistic limits of Kerr and Kerr-Newman black holes were obtained , and recently, boosted static multipole (Weyl) particles were studied . Impulsive gravitational waves were also generated by boosting the Schwarzschild-de Sitter and Schwarzschild-anti de Sitter metrics to the ultrarelativistic limit ; see the contribution by J. Griffiths and J. Podolský in these Proceedings. These types of spacetimes, especially the simple Aichelburg-Sexl metrics, have been employed in current problems of the generation of gravitational radiation from axisymmetric black hole collisions and black hole encounters. The recent monograph by d’Eath gives a comprehensive survey, including the author’s new results. There is good reason to believe that spacetime metrics produced in high speed collisions will be simpler than those corresponding to (more realistic) situations in which black holes start to collide with low relative velocities. The spacetimes corresponding to the collisions at exactly the speed of light is an interesting limit which can be treated most easily. Aichelburg-Sexl metrics are used to describe limiting “incoming states” of two black holes, moving one against the other with the speed of light. Great interest has been stimulated by ’t Hooft’s work on the quantum scattering of two pointlike particles at centre-of-mass energies higher or equal to the Planck energy. This quantum process has been shown to have close connection with classical black hole collisions at the speed of light (see and references therein). Recently, the Colombeau algebra of generalized functions, which enables one to deal with singular products of distributions, has been brought to general relativity and used in the description of impulsive pp-waves in various coordinate systems , and also for a rigorous solution of the geodesic and geodesic deviation equations for impulsive waves . The investigation of the equations of geodesics in non-homogeneous pp-waves (with $`f\zeta ^3`$) has shown that the motion of test particles is chaotic (see and the contribution by J. Podolský in these Proceedings). Plane-fronted waves have been used as simple metrics in various other contexts, for example, in quantum field theory on a given background (see e.g. ). As emphasized very recently by Gibbons , since for pp-waves and type N Kundt’s class all possible invariants formed from the Weyl tensor and its covariant derivatives vanish , these metrics suffer no quantum corrections to all loop orders. Thus they may offer insights into the behaviour of a full quantum theory. ### Colliding plane waves The first detailed study of colliding plane waves was undertaken independently by Khan and Penrose and by Szekeres (see for references). Szekeres formulated the problem as a characteristic initial value problem for a system of hyperbolic equations in two variables (null coordinates) $`u,v`$ with data specified on the pair of null hypersurfaces, say $`u=0,v=0`$ intersecting in a spacelike 2-surface. Although Szekeres’ formulation of a general solution for the problem of colliding parallel-polarized waves is difficult to use for constructing explicit solutions, it has been employed in a general analysis of the structure of the singularities produced by the collision . It has also inspired the work developed at the beginning of the 1990s by Hauser and Ernst . Their new method of analyzing the initial value problem can be used also when the polarization of the approaching waves is not aligned. They formulated the initial value problem in terms of the equivalent matrix Riemann-Hilbert problem. Their techniques are related to those used by Neugebauer and Meinel to construct and analyze the rotating disk solution as a boundary value problem (see their contributions to these Proceedings). Most recently, Hauser and Ernst prepared an extensive treatise in which they give a general description and detailed mathematical proofs of their study of the solutions of the hyperbolic Ernst equation. The papers on colliding plane waves published until 1991 are reviewed in (see also ). New developments have been mostly involved with “non-classical” issues like the inclusion of dilatonic fields or the discussion of the particle production. One of the few exceptions has been the analysis of colliding waves in the expanding backgrounds as e.g. in Friedmann-Robertson-Walker universes filled by stiff fluid . In contrast to the waves propagating and colliding on the “flat backgrounds”, no singularities arise in the expanding backgrounds. ## 3 Cylindrical waves Despite the fact that cylindrically symmetric waves cannot describe exactly the radiation from bounded sources, they even recently played an important role in clarifying a number of complicated issues, such as testing the quasilocal mass-energy , testing codes in numerical relativity , investigation of the cosmic censorship , and quantum gravity (see for more details and references). In recent work with Ashtekar and Schmidt , we considered gravitational waves with a space-translation Killing field (“generalized Einstein-Rosen waves”). In the (2+1)-dimensional framework the Einstein-Rosen subclass forms a simple instructive example of explicitly given spacetimes which admit a smooth global null (and timelike) infinity even for strong initial data. 4-dimensional vacuum gravity which admits a spacelike hypersurface Killing vector $`/z`$ is equivalent to 3-dimensional gravity coupled to a scalar field. In 3 dimensions, there is no gravitational radiation. Hence, the local degrees of freedom are all contained in the scalar field. One therefore expects that Cauchy data for the scalar field will suffice to determine the solution. For data which fall off appropriately, we thus expect the 3-dimensional Lorentzian geometry to be asymptotically flat in the sense of Penrose, i.e. that there should exist a 2-dimensional boundary representing null infinity. In general cases, this is analyzed in . Restricting here ourselves to the Einstein-Rosen waves by assuming that there is a further spacelike, hypersurface orthogonal Killing vector $`/\phi `$ which commutes with $`/z`$, we find the 3-metric given by $$d\sigma ^2=g_{ab}dx^adx^b=e^{2\gamma }(dt^2+d\rho ^2)+\rho ^2d\phi ^2,$$ (3) where $`\gamma =\gamma (t,\rho )`$. The field equations for the scalar field $`\psi `$ coupled to this metric become $$\ddot{\psi }+\psi ^{\prime \prime }+\rho ^1\psi ^{}=0,\gamma ^{}=\rho (\dot{\psi }^2+\psi ^2),\dot{\gamma }=2\rho \dot{\psi }\psi ^{}.$$ (4) Thus, we can first solve the axisymmetric wave equation for $`\psi `$ on Minkowski space and then solve for $`\gamma `$ – the only unknown metric coefficient – by quadratures. The “method of descent” from the Kirchhoff formula in 4 dimensions gives the representation of the solution of the wave equation in 3 dimensions in terms of Cauchy data $`\mathrm{\Psi }_0=\psi (t=0,x,y),\mathrm{\Psi }_1=\psi _{,t}(t=0,x,y)`$ (see ). We assume that the Cauchy data are axially symmetric and of compact support. Investigating the behaviour of the solution at future null infinity $`𝒥^+`$, one finds $$\psi (u,\rho )=\frac{f_0(u)}{\sqrt{\rho }}+\frac{1}{\sqrt{\rho }}\underset{k=1}{\overset{\mathrm{}}{}}\frac{f_k(u)}{\rho ^k},$$ (5) where $`u=t\rho `$ and the coefficients $`f`$’s are determined by the Cauchy data. The field equations imply $$\gamma =\gamma _02_{\mathrm{}}^u\left[\dot{f}_0(u)\right]^2𝑑u\underset{k=1}{\overset{\mathrm{}}{}}\frac{h_k(u)}{(k+1)\rho ^{k+1}}.$$ (6) Thus, $`\gamma `$ also admits an expansion in $`\rho ^1`$. It is straightforward to show that the spacetime admits a smooth future null infinity by setting $`\stackrel{~}{\rho }=\rho ^1,\stackrel{~}{u}=u,\stackrel{~}{\phi }=\phi `$ and rescaling $`g_{ab}`$ by a conformal factor $`\mathrm{\Omega }=\stackrel{~}{\rho }`$. Hence, the (2+1)-dimensional curved spacetime has a smooth (2-dimensional) null infinity. Penrose’s picture works for arbitrarily strong initial data $`\mathrm{\Psi }_0`$, $`\mathrm{\Psi }_1`$. We can thus conclude that cylindrical waves in (2+1)-dimensions give an explicit model of the Bondi-Penrose radiation theory which admits smooth null and timelike infinity for arbitrarily strong initial data. There is no other such model available. The general results on the existence of $`𝒥`$ in 4 dimensions assume weak data. ## 4 On the Robinson-Trautman and type N twisting solutions These spacetimes have attracted increased attention in the last decade – most notably in the work by Chruściel, and Chruściel and Singleton . In these studies the Robinson-Trautman spacetimes have been shown to exist globally for all positive “times”, and to converge asymptotically to a Schwarzschild metric. Interestingly, the extension of the spacetimes across the “Schwarzschild-like” event horizon can only be made with a finite degree of smoothness. These studies are based on the derivation and analysis of an asymptotic expansion describing the long-time behaviour of the solutions of the nonlinear parabolic Robinson-Trautman equation. In our recent work , we studied Robinson-Trautman spacetimes with a positive cosmological constant $`\mathrm{\Lambda }`$. The results proving the global existence and convergence of the solutions of the Robinson-Trautman equation can be taken over from the previous studies since $`\mathrm{\Lambda }`$ does not explicitly enter this equation. We have shown that, starting with arbitrary, smooth initial data at $`u=u_0`$, these cosmological Robinson-Trautman solutions converge exponentially fast to a Schwarzschild-de Sitter solution at large retarded times ($`u\mathrm{}`$). The interior of a Schwarzschild-de Sitter black hole can be joined to an “external” cosmological Robinson-Trautman spacetime across the horizon $`^+`$ with a higher degree of smoothness than in the corresponding case with $`\mathrm{\Lambda }=0`$. In particular, in the extreme case with $`9\mathrm{\Lambda }m^2=1`$, in which the black hole and cosmological horizons coincide, the Robinson-Trautman spacetimes can be extended smoothly through $`^+`$ to the extreme Schwarzschild-de Sitter spacetime with the same values of $`\mathrm{\Lambda }`$ and $`m`$. However, such an extension is not analytic (and not unique). We have also demonstrated that the cosmological Robinson-Trautman solutions are explicit models exhibiting the cosmic no-hair conjecture. As far as we are aware, these models represent the only exact analytic demonstration of the cosmic no-hair conjecture under the presence of gravitational waves. They also appear to be the only exact examples of a black hole formation in nonspherical spacetimes which are not asymptotically flat. ### Type N twisting spacetimes Since diverging, non-twisting Robinson-Trautman spacetimes of type N have singularities, there has been hope that if one admits a nonvanishing twist a more realistic radiative spacetime may exist. Stephani , however, indicated, by constructing a general solution of the linearized equations, that singularities at infinity probably occur. More recently, Finley et al found an approximate twisting type N solution up to the third order of iteration on the basis of which they suggested that it seems that the twisting, type N fields can describe a radiation field outside bounded sources. However, employing the Newman-Penrose formalism and MAPLE we succeeded in discovering a nonvanishing quartic invariant in the 2nd derivatives of the Riemann tensor , which shows that solutions of both Stephani and Finley et al contain singularities at large $`r`$. Very recently, Mac Alevey argued that an approximate solution at any finite order can be calculated without occurrence of singularities. It is very likely, however, that a corresponding exact solution must contain singularities since Mason proved that the only vacuum algebraically special spacetime that is asymptotically simple is the Minkowski space. Even if a radiative solution with a complete smooth null infinity may be out of reach, it is of interest to construct radiative solutions which admit at least a global null infinity in the sense that its smooth cross sections exist although this null infinity is not necessarily complete. The only explicit examples of such solutions are spacetimes with boost-rotation symmetry. ## 5 The boost-rotation symmetric radiative spacetimes I reviewed these spacetimes representing “uniformly accelerated objects” in various places (see e.g. and references therein); here I shall just mention some new results. The unique role of the boost-rotation symmetric spacetimes is exhibited by a theorem which roughly states that in axially symmetric, locally asymptotically flat electrovacuum spacetimes (in the sense that a null infinity satisfying Penrose’s requirements exists, but it need not necessarily exist globally), the only additional symmetry that does not exclude radiation is the boost symmetry. To prove such a result we start from the metric $`ds^2`$ $`=`$ $`\left(r^1Ve^{2\beta }r^2e^{2\gamma }U^2\mathrm{cosh}2\delta r^2e^{2\gamma }W^2\mathrm{cosh}2\delta 2r^2UW\mathrm{sinh}2\delta \right)du^2`$ $`2e^{2\beta }dudr2r^2\left(e^{2\gamma }U\mathrm{cosh}2\delta +W\mathrm{sinh}2\delta \right)dud\theta `$ $`2r^2\left(e^{2\gamma }W\mathrm{cosh}2\delta +U\mathrm{sinh}2\delta \right)\mathrm{sin}\theta dud\varphi `$ $`+r^2\left[\mathrm{cosh}2\delta \left(e^{2\gamma }d\theta ^2+e^{2\gamma }\mathrm{sin}^2\theta d\varphi ^2\right)+2\mathrm{sinh}2\delta \mathrm{sin}\theta d\theta d\varphi \right],`$ where all funtions describing the metric and electromagnetic field tensor $`F_{\mu \nu }`$ are independent of $`\varphi `$. Assuming asymptotic expansions of these functions at large $`r`$ with $`u`$, $`\theta `$, $`\varphi `$ fixed to guarantee asymptotic flatness, and using the outgoing radiation condition and the field equations, one finds the expansions to have specific forms. For example, $$\gamma =\frac{c}{r}+(C{\scriptscriptstyle \frac{1}{6}}c^3{\scriptscriptstyle \frac{3}{2}}cd^2)\frac{1}{r^3}+\mathrm{},V=r2M+\mathrm{},$$ $$F_{02}=X+(ϵ_{,\theta }e_{,u})\frac{1}{r}+\mathrm{},F_{03}=Y\frac{f_{,u}}{r}+\mathrm{},$$ (8) where the ‘coefficients’ $`c`$, $`d`$, … are functions of $`u`$ and $`\theta `$. The expansions are needed to further orders – see for their rather lengthy forms. Let us only recall that the decrease of the Bondi mass, $`m(u)=\frac{1}{2}_0^\pi M(u,\theta )\mathrm{sin}\theta d\theta ,`$ is given by $$m_{,u}={\scriptscriptstyle \frac{1}{2}}\underset{0}{\overset{\pi }{}}(c,_u^2+d,_u^2+X^2+Y^2)\mathrm{sin}\theta d\theta 0,$$ (9) where $`c_{,u}`$, $`d_{,u}`$, $`X`$, $`Y`$ are the gravitational and electromagnetic news functions. Now one writes down the Killing equations and solves them asymptotically in $`r^1`$. One arrives at the following theorem : Suppose that an axially symmetric electrovacuum spacetime admits a “piece” of $`𝒥^+`$ in the sense that the Bondi-Sachs coordinates can be introduced in which the metric takes the form (5), with the asymptotic form of the metric and electromagnetic field given by (8). If this spacetime admits an additional Killing vector forming with the axial Killing vector a two-dimensional Lie algebra, then the additional Killing vector has asymptotically the form $$\eta ^\alpha =[ku\mathrm{cos}\theta +\alpha (\theta ),kr\mathrm{cos}\theta +𝒪(r^0),k\mathrm{sin}\theta +𝒪(r^1),𝒪(r^1)],$$ (10) where $`k`$ is a constant. For $`k=0`$ it generates asymptotically translations (function $`\alpha `$ has then a specific form). For $`k0`$ it is the boost Killing field. The case of translations is analyzed in detail in . Theorem 1, precisely formulated and proved there, states that if an asymptotically translational Killing vector is spacelike, then null infinity is singular at some $`\theta 0,\pi `$; if it is null, null infinity is singular at $`\theta =0`$ or $`\pi `$. The first case corresponds to cylindrical waves, the second case to a plane wave propagating along the symmetry axis. We refer to for the case when there is also a cosmic string present along the symmetry axis. The case of timelike Killing vector is described by Theorem 2 (proved also in ): If an axisymmetric electrovacuum spacetime with a non-vanishing Bondi mass admits an asymptotically translational Killing vector and a complete cross section of $`𝒥^+`$, then the translational Killing vector is timelike and spacetime is thus stationary. The case of the boost Killing vector $`(k0)`$ is thoroughly analyzed in . The general functional forms of the news functions (both gravitational and electromagnetic), and of the mass aspect and total Bondi mass of boost-rotation symmetric spacetimes are there given. Very recently these results were obtained by using the Newman-Penrose formalism and under more general assumptions (for example, $`𝒥`$ could in principle be polyhomogeneous). The general structure of the boost-rotation symmetric spacetimes with hypersurface orthogonal Killing vectors was analyzed in detail in . Their radiative properties, including explicit construction of radiation patterns and of Bondi mass for the specific boost-rotation symmetric solutions were investigated in several works – we refer to the reviews and for details. There also the role of the boost-rotation symmetric spacetimes in such diverse fields like numerical relativity and quantum production of black-hole pairs is noticed and references are given. Here I would like to mention yet a recent progress in understanding specific boost-rotation symmetric spacetimes with Killing vectors which are not hypersurface orthogonal. This is the spinning $`C`$-metric (see e.g. ). It was discovered by Plebański and Demaiński as a generalization of the standard $`C`$-metric which is known to represent uniformly accelerated non-rotating black holes. In we first transformed the metric into Weyl coordinates, and then found a transformation which brings it into the canonical form of the radiative spacetimes with the boost-rotation symmetry: $`ds^2`$ $`=`$ $`e^\lambda d\rho ^2+\rho ^2e^\mu d\varphi ^2`$ $`+`$ $`(z^2t^2)^1\left[(e^\lambda z^2e^\mu t^2)dz^22zt(e^\lambda e^\mu )dzdt+(e^\lambda t^2e^\mu z^2)dt^2\right]`$ $``$ $`2𝒜e^\mu (zdttdz)d\varphi 𝒜^2e^\mu (z^2t^2)d\varphi ^2,`$ (11) where functions $`e^\mu ,e^\lambda `$ and $`𝒜`$ are given in terms of $`(t,\rho ,z)`$ in a somewhat complicated but explicit manner. This metric can represent two uniformly accelerated, spinning black holes, either connected by a conical singularity, or with conical singularities extending from each of them to infinity. The behaviour of the curvature invariants clearly indicates the presence of a non-vanishing radiation field (see Figure 5 in ). The spinning $`C`$-metric is the only explicitly known example with two Killing vectors which are not hypersurface orthogonal, in which one can give arbitrarily strong initial data on a hyperboloid “above the roof” ($`t>\|z\|`$) which evolve into the radiative spacetime with smooth $`𝒥^+`$. ## 6 Inhomogeneous cosmologies and gravitational waves Among the known vacuum inhomogeneous models, the Gowdy solutions (see e.g. ) have played the most distinct role. They belong to the class of solutions with two commuting spacelike Killing vectors. Within a cosmological context, they form a subclass of a wider class of $`G_2`$ cosmologies – as are now commonly denoted models which admit an Abelian group $`G_2`$ of isometries with orbits being spacelike 2-surfaces. A 2-surface with a 2-parameter isometry group must be a space of constant curvature, and since neither a 2-sphere nor a 2-hyperboloid possess 2-parameter subgroups, it must be intrinsically flat. If the 2-surface is an Euclidean plane or a cylinder, then one speaks about planar or cylindrical universes. Gowdy universes are compact – the group orbits are 2-tori $`T^2`$. The metrics with two spacelike Killing vectors are often called the generalized Einstein-Rosen metrics as, for example, by Carmeli, Charach and Malin in their comprehensive survey of inhomogeneous cosmological models of this type. In dimensionless coordinates ($`t,z,x^1,x^2`$), the line element can be written as ($`A,B=1,2`$) $$ds^2/L^2=e^F(dt^2+dz^2)+\gamma _{AB}dx^Adx^B,$$ (12) where $`L`$ is a constant length, $`F`$ and $`\gamma _{AB}`$ depend on $`t`$ and $`z`$ only, and thus the spacelike Killing vectors are $`{}_{}{}^{(1)}\xi _{}^{\alpha }=(0,0,1,0),{}_{}{}^{(2)}\xi _{}^{\alpha }=(0,0,0,1).`$ Let us mention some recent developments in which the Gowdy models have played a role. Gowdy-type models have been used to study the propagation and collision of gravitational waves with toroidal wavefronts in the FRW closed universes with a stiff fluid . In the standard Gowdy spacetimes it is assumed that the “twists” associated with the isometry group on $`T^2`$ vanish. In the generalized Gowdy models without this assumption are considered, and their global existence in time is proved. As both interesting and non-trivial models, the Gowdy spacetimes have recently attracted the attention of mathematical and numerical relativists. Chruściel, Isenberg and Moncrief proved that Gowdy spacetimes developed from a dense subset in the initial data set cannot be extended past their singularities, i.e. in “most” Gowdy models the strong cosmic censorship is satisfied. On cosmic censorship and spacetime singularities, especially in the context of compact cosmologies, we refer to . This review shows how intuition gained from such solutions as the Gowdy models or the Taub-NUT spaces, when combined with new mathematical ideas and techniques, can produce rigorous results with a generality out of reach until recently. To such results belongs also the very recent work of Kichenassamy and Rendall on the sufficiently general class of solutions (containing the maximum number of arbitrary functions) representing unpolarized Gowdy spacetimes. The new mathematical technique, the so called Fuchsian algorithm, enables one to construct singular (and nonsingular) solutions of partial differential equations with a large number of arbitrary functions, and thus provide a description of singularities. Applying the Fuchsian algorithm to Einstein’s equations for Gowdy spacetimes with topology $`T^3`$, Kichenassamy and Rendall have proved that general solutions behave at the (past) singularity in a Kasner-like manner, i.e. they are asymptotically velocity dominated with a diverging curvature invariant. One needs an additional magnetic field not aligned with the two Killing vectors of the Gowdy unpolarized spacetimes in order to get a general oscillatory (Mixmaster) approach to a singularity, as shown by the numerical calculations . Some metrics can be considered as exact “gravitational solitons” propagating on a cosmological background. Verdaguer prepared a very complete review of solitonic cosmological solutions admitting two spacelike Killing vector fields. Recently, differential conservation laws for large perturbations of gravitational field with respect to a given curved background have been fomulated . They should bring more light also on various solitonic models in cosmology. I thank the organizers for inviting me to the interesting meeting and Tomáš Ledvinka for the help with the manuscript. Support from the grant No. GAČR 202/99/0261 of the Czech Republic is acknowledged.
warning/0004/math0004132.html	ar5iv	text	# 1 Introduction ## 1 Introduction The purpose of this paper is to study the space of 1-point correlation functions, or trace functions, which arise from certain vertex operator algebras. Let us suppose that $`V`$ is a vertex operator algebra \[FLM\] (also see \[B\], \[MN\]) with standard $`L(0)`$-grading $$V=\underset{n}{}V_n.$$ The most basic trace function is the formal graded character $$\mathrm{ch}_qV=\mathrm{tr}_Vq^{L(0)c/24}=q^{c/24}\underset{n}{}(dimV_n)q^n$$ where $`c`$ is the central charge of $`V.`$ For many well-known VOAs, the graded character has certain modular-invariance properties. For example, if $`V(𝔤,l)`$ is the generalized Verma module of level $`l`$ associated to a (complex) Lie algebra $`𝔤`$ of dimension $`d`$ (cf. \[FZ\], \[L\]) and trivial $`𝔤`$-module, then $`V(𝔤,l)`$ is a vertex (operator) algebra satisfying $$\mathrm{ch}_qV(𝔤,l)=\eta (q)^d$$ where $`\eta (q)`$ is the Dedekind eta-function. Similarly, if $`L`$ is a positive-definite, even lattice of rank $`d`$ and $`\theta _L(q)`$ is the corresponding theta-function, and if $`V_L`$ is the associated VOA (cf. \[B\], \[FLM\]), then $$\mathrm{ch}_qV_L=\theta _L(q)/\eta (q)^d.$$ So $`\mathrm{ch}_qV_L`$ is a modular function (i.e., of weight zero) on some congruence subgroup of the modular group $`SL(2,),`$ while $`\mathrm{ch}_qV(𝔤,l)`$ is a modular form of weight $`d/2`$ on the full modular group<sup>4</sup><sup>4</sup>4We will have no need to concern ourselves with the fact that $`d`$ may be odd, so that the weight may be a half integer.. On the other hand, the VOA $`M(c,0)`$ associated to the Virasoro algebra of central charge $`c`$ (cf. \[FZ\], \[L\]) satisfies $$\mathrm{ch}_qM(c,0)=(1q)q^{c/24+1}/\eta (q)$$ and hence is not modular. Here, of course, we are interpreting $`q`$ in the usual way to be equal to $`e^{2\pi i\tau }`$ with $`\tau `$ in the complex upper-half plane $`H.`$ In the fundamental paper of Zhu \[Zh\] general correlation functions, which generalize the graded character, were studied. Recall (loc cit) that if $`v`$ lies in $`V_k`$ with vertex operator $`Y(v,z)=_nv(n)z^{n1},`$ then the so-called zero mode $`o(v)=v(k1)`$ is a linear operator on V which leaves invariant each homogeneous space $`V_n,`$ so that one can form the expression $$Z(v,q)=\mathrm{tr}_Vo(v)q^{L(0)c/24}=q^{c/24}\underset{n}{}(\mathrm{tr}_{V_n}o(v))q^n.$$ One calls the linear extension of $`Z(v,q)`$ to all of V the (1-point) correlation function determined by $`V.`$ It is the purpose of this paper to understand the nature of this function in the case of the lattice VOAs $`V_L`$ and the Heisenberg VOA $`M(1),`$ which is the space $`V(𝔤,l)`$ in the case that $`𝔤`$ is an abelian Lie algebra equipped with a non-degenerate symmetric bilinear form $`(,)`$ and $`l=1.`$ $`M(1)`$ is often referred to as the VOA of $`d`$ free bosons if $`dim𝔤=d.`$ In particular, we ask when $`Z(v,q)`$ is modular in a suitable sense and, if not, how does it deviate from being modular? In order to explain our results and methods, we need to introduce so-called quasi-modular forms \[KZ\], which play an interesting role in the present paper. For a positive integer $`N`$ and for a Dirichlet character $`ϵ`$ modulo $`N`$, let $`M(N,ϵ)`$ denote the ring of modular forms on the congruence subgroup $`\mathrm{\Gamma }_0(N)`$ which transform according to the character $`ϵ`$ and which are holomorphic in $`H.`$ Let (1.1) $$Q(N,ϵ)=M(N,ϵ)[E_2]$$ be the space obtained by adjoining to $`M(N,ϵ)`$ the Eisenstein series $`E_2(q).`$ For example, if $`N=1,`$ so that also $`ϵ=1,`$ then we have (1.2) $$Q=Q(1,1)=C[E_2,E_4,E_6]$$ the full ring of quasi-modular forms on $`SL(2,).`$ We can now state our first result: ###### Theorem 1. Let $`V`$ be the VOA of $`d`$ free bosons, with $`Q`$ as in (1.2). (a) If $`v`$ is in $`V,`$ then $`Z(v,q)`$ converges to a holomorphic function $`f(v,q)/\eta (q)^d`$ in the upper half plane for some $`f(v,q)Q.`$ (b) Every $`f(q)`$ in $`Q`$ may be realized as $`f(v,q)`$ for some $`v`$ in $`V.`$ We can restate Theorem 1 in the following form: there is a linear surjection (1.3) $$\begin{array}{cccc}t:& V& & Q\hfill \\ & v& & (\mathrm{ch}_qV)^1Z(v,q).\hfill \end{array}$$ In fact we will construct a section of the map t, more precisely we explicitly describe a subspace $`W`$ of $`V`$ such that the restriction of $`t`$ to $`W`$ is an isomorphism onto $`Q.`$ ###### Theorem 2. Let $`L`$ be a positive-definite even lattice of rank $`d`$ and level $`N,`$ so that the theta function $`\theta _L(q)`$ lies in $`M(N,ϵ)`$ for suitable $`ϵ.`$ Then for every element $`v`$ in $`V_L,`$ we have (1.4) $$Z(v,q)=f(v,q)/\eta (q)^d$$ for some $`f(v,q)`$ in $`M(N,ϵ).`$ In particular, each $`Z(v,q)`$ is a sum of modular forms (of varying weights). The role of quasi-modular forms in the proof of Theorem 2 is more hidden. To explain, we recall here (and discuss in greater detail in Section 4.1) that in \[Zh\], Zhu shows how to identify vectors in $`V`$ such that the corresponding trace function $`Z(v,q)`$ essentially behaves like a modular form of weight $`k`$ for some $`k.`$ More precisely, if we think of $`Z(v,q)`$ as a function of $`\tau ,`$ then under the action of the modular group $$\gamma =\left(\begin{array}{cc}a& b\\ c& d\end{array}\right):Z(v,\tau )(c\tau +d)^kZ(v,\gamma \tau ),\gamma SL(2,),$$ the $`\gamma `$-transforms of $`Z(v,\tau )`$ spans a finite-dimensional $`SL(2,)`$-module of holomorphic functions on $`H.`$ Moreover each such $`\gamma `$ transform has a suitable q-expansion. Extending the language of \[KM\], we call such functions generalized modular forms of weight $`k.`$ A modular form has this property, however it is shown in \[KM\] that there are generalized modular forms which are not modular. Thus one cannot deduce from Zhu’s results alone that the trace functions $`Z(v,q)`$ are modular, or even sums of modular forms. However, we will show below that the trace functions occurring in Theorem 2 also have the property that $`\eta (q)^dZ(v,q)`$ lies in $`Q(N,ϵ)`$ so that in this regard they behave similarly to the trace functions in the free boson case. It is easy to see that a generalized modular form that is also a quasi-modular form is necessarily an ordinary modular form, and Theorem 2 then follows. Our third main theorem has a slightly different flavor. In the paper \[DM1\] the space of trace functions $`Z(v,q)`$ was determined in the case of the Moonshine module $`V^{\mathrm{}}`$ \[FLM\], a prominent role being played by the primary fields i.e., the (homogeneous) vectors of $`V^{\mathrm{}}`$ which are highest weight vectors for the Virasoro algebra. Based on this work and the calculations in \[HL\], we conjecture that if $`V`$ is a holomorphic VOA then each cusp form on $`SL(2,)`$ (possibly with character) can be realized by a trace function $`Z(v,q)`$ in which $`v`$ is a primary field. The conjecture seems to be non-trivial for any holomorphic VOA. We will prove it for the lattice VOA $`V_{E_8}`$ based on the $`E_8`$ root lattice. Our approach uses the theory of spherical harmonics, and through this mechanism one sees that the conjecture may be viewed as a conformal field theoretic analog of the following problem in number theory: given a positive-definite self-dual even lattice $`L,`$ describe the space of modular forms $`\theta _L(P,\tau )`$ obtained by modifying $`\theta _L(\tau )`$ by $`P,`$ where $`P`$ ranges over the homogeneous spherical harmonic functions with respect to $`L.`$ (Replacing “holomorphic” by “rational” in our conjecture corresponds to eliminating the self-duality of $`L.`$) Waldspurger showed \[W\] that all cusp forms of level one can be obtained as $`\theta _{E_8}(P,\tau )`$ for suitable $`P,`$ and this leads to our result about $`V_{E_8}`$ because of the following. ###### Theorem 3. Let $`P`$ be a homogeneous spherical harmonic of degree $`k`$ with respect to the lattice $`L`$ of rank $`d.`$ Then there is a primary field $`v_P`$ in the lattice VOA $`V_L`$ with the property that $$Z(v_P,q)=\theta _L(P,q)/\eta (q)^d.$$ There is a circle of ideas relating our results: in \[EZ\], Eichler and Zagier gave an approach to the result of Waldspurger using ideas from the theory of Jacobi forms. Jacobi forms and Jacobi-like forms \[Za\] are closely related to quasi-modular forms, and they also play a role in the results of \[DM2\] which identify the quasi-modular forms underlying the proof of theorem 2 and relate them to theta functions modified by a spherical harmonic. And the occurrence of Jacobi forms in string theory has been observed frequently in the past few years. ## 2 Vertex operator algebras and graded traces In this section we briefly review from \[Zh\] the “bracket” vertex operator algebra $`(V,Y[],\mathrm{𝟏},\omega c/24)`$ constructed from $`(V,Y,\mathrm{𝟏},\omega ).`$ We also list several formulas involving the trace functions from \[Zh\] and present an easy corollary which is used frequently in the later sections. ### 2.1 Vertex operator algebras of genus one Let $`V=(V,Y,\mathrm{𝟏},\omega )`$ be a vertex operator algebra. Then $`V`$ may be regarded as a vertex operator algebra on the sphere. In order to study modular invariance in the theory of vertex operator algebras, a new vertex operator algebra on the torus was introduced in \[Zh\]. The new vertex operator algebra is $`(V,Y[],\mathrm{𝟏},\omega c/24)`$ where $`c`$ is the central charge of $`V.`$ The new vertex operator associated to a homogeneous element $`a`$ is given by $$Y[a,z]=\underset{n}{}a[n]z^{n1}=Y(a,e^z1)e^{z\mathrm{wt}(a)}$$ while a Virasoro element is $`\stackrel{~}{\omega }=\omega c/24`$. Thus $$a[m]=\mathrm{Res}_z\left(Y(a,z)(\mathrm{ln}(1+z))^m(1+z)^{\mathrm{wt}(a)1}\right)$$ and $$a[m]=\underset{i=m}{\overset{\mathrm{}}{}}c(\mathrm{wt}(a),i,m)a(i)$$ for some scalars $`c(\mathrm{wt}(a),i,m)`$ such that $`c(\mathrm{wt}(a),m,m)=1.`$ In particular, $$a[0]=\underset{i0}{}\left(\genfrac{}{}{0pt}{}{\mathrm{wt}(a)1}{i}\right)a(i).$$ We also write $$L[z]=Y[\omega ,z]=\underset{n}{}L[n]z^{n2}.$$ Then the $`L[n]`$ again generate a copy of the Virasoro algebra with the same central charge $`c.`$ Now $`V`$ is graded by the $`L[0]`$-eigenvalues, that is $$V=\underset{n}{}V_{[n]}$$ where $`V_{[n]}=\{vV\|L[0]v=nv\}.`$ We also write $`\mathrm{wt}[a]=n`$ if $`aV_{[n]}.`$ It should be pointed out that for any $`n`$ we have $$\underset{mn}{}V_n=\underset{mn}{}V_{[n]}.$$ We also recall the notion of $`V`$-module briefly. A $`V`$-module $`M=(M,Y^M)`$ is a $``$-graded vector space $$M=_\lambda M_\lambda $$ such that each $`M_\lambda `$ is finite-dimensional and $`M_{\lambda +n}`$ is zero if n is a small enough integer. Furthermore $`M_\lambda `$ is the eigenspace of $`L(0)`$ with eigenvalue $`\lambda `$ where $`L(0)`$ is the component operator of $`Y^M(\omega ,z)=_nL(n)z^{n2}.`$ As in the case of vertex operator algebras, the component operator $`o(a)=a^M(\mathrm{wt}(a)1)`$ preserves each homogeneous space $`M_\lambda `$ if $`aV`$ is homogeneous and $`Y^M(a,z)=_na^M(n)z^{n1}.`$ Again using the $``$-linearity, we extend the operator $`o(a)`$ to all $`aV.`$ ### 2.2 Graded traces Next we state some results from \[Zh\]. We use the Eisenstein series $`E_{2k}(\tau )`$ normalized as in \[DLM\] equation (4.28). Thus: $$E_{2k}(\tau )=\frac{B_{2k}}{2k!}+\frac{2}{(2k1)!}\underset{n=1}{\overset{\mathrm{}}{}}\sigma _{2k1}(n)q^n$$ where $`\sigma _k(n)`$ is the sum of the $`k`$-powers of the divisors of $`n`$ and $`B_{2k}`$ a Bernoulli number. Let $`M=(M,Y^M)`$ be a $`V`$-module. For any $`aV`$ we define a formal power series in $`q:`$ $$Z_M(a,q)=\mathrm{tr}_Mo(a)q^{L(0)c/24}=q^{c/24}\underset{\lambda }{}(\mathrm{tr}_{M_\lambda }o(a))q^\lambda .$$ Then $`Z_V(a,q)`$ is exactly $`Z(a,q)`$ defined before. ###### Proposition 2.2.1 (\[Zh\], Proposition 4.3.5, 4.3.6). Let $`M`$ be a $`V`$-module. Then the following identities hold as formal power series for any $`a,bV.`$ (2.2.1) $`Z_M(a[0]b,q)=0,`$ (2.2.2) $`\mathrm{tr}_Mo(a)o(b)q^{L(0)c/24}=Z_M(a[1]b,q){\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}E_{2k}(q)Z_M(a[2k1]b,q),`$ (2.2.3) $`Z_M(a[2]b,q)={\displaystyle \underset{k=2}{\overset{\mathrm{}}{}}}(2k1)E_{2k}(q)Z_M(a[2k2]b,q).`$ The following corollary plays an important role in our arguments. ###### Corollary 2.2.2. Let notation be as before. Then for each positive integer $`r`$ we have (2.2.4) $`Z_M(a[r]b,q)`$ $`=\delta _{1,r}\mathrm{tr}_Mo(a)o(b)q^{L(0)c/24}`$ $`+(1)^{r+1}{\displaystyle \underset{k>r/2}{\overset{\mathrm{}}{}}}h(k,r)E_{2k}(q)Z_M(a[2kr]b,q)`$ where $`h(k,r)=\left(\genfrac{}{}{0pt}{}{2k1}{r1}\right).`$ ###### Proof. We prove the result by induction on $`r`$. The cases $`r=1,2`$ are nothing but (2.2.2) and (2.2.3) respectively. Suppose the statement is true for $`r2`$. Then replacing $`a`$ by $`L[1]a`$ and noting $`(L[1]a)[n]=na[n1]`$, we see that $$rZ_M(a[r1]b,q)=(1)^{r+1}\underset{k>r/2}{}(r2k)h(k,r)E_{2k}(q)Z_M(a[2kr1]b,q).$$ Then by (2.2.1), we conclude that $$Z_M(a[r1]b,q)=(1)^{r+2}\underset{k>(r+1)/2}{\overset{\mathrm{}}{}}h(k,r+1)E_{2k}(q)Z_M(a[2kr1]b,q).$$ ## 3 Quasi-modularity of trace functions for free boson VOAs In this section, after reviewing from \[FLM\] the free boson vertex operator algebra $`M(1),`$ we establish Theorem 1. ### 3.1 Free boson VOAs Let $`𝔥`$ be a $`d`$-dimensional vector space with a non-degenerate symmetric bilinear form $`(,)`$ and $`\widehat{𝔥}`$ be the corresponding affinization viewing $`𝔥`$ as an abelian Lie algebra: $`\widehat{𝔥}=𝔥[t,t^1]K`$ with commutator relations $`[ht^m,h^{}t^n]=(h,h^{})\delta _{m+n,0}K,(h,h^{}𝔥,m,n),`$ $`[K,𝔥[t,t^1]]=0.`$ Consider the induced module $$M(1)=U(\widehat{𝔥})_{𝔥[t]K}$$ where $`𝔥[t]`$ acts trivially on $``$ and $`K`$ acts as $`1`$. We denote by $`h(n)`$ the action of $`ht^n`$ on $`M(1)`$. The space $`M(1)`$ is linearly isomorphic to the symmetric algebra $`S(𝔥t^1[t^1])`$. Thus setting $`\mathrm{𝟏}=11`$, any element in $`M(1)`$ is a linear combination of elements of type $$v=a_1(n_1)\mathrm{}a_k(n_k)\mathrm{𝟏},(a_1,\mathrm{},a_k𝔥,n_1,\mathrm{},n_k_+).$$ For such $`v`$, we define (3.1.1) $$Y(v,z)=^{_{}}_{^{}}^{(n_11)}a_1(z)\mathrm{}^{(n_k1)}a_k(z)^{_{}}_{^{}},^{(n)}=\frac{1}{n!}\left(\frac{d}{dz}\right)^n$$ where (3.1.2) $$a(z)=\underset{n}{}a(n)z^{n1}$$ and $`^{_{}}_{^{}}`$ $`^{_{}}_{^{}}`$ indicates the normal ordering procedure. Now let $`\{h_i\}_{i=1}^d`$ be an orthonormal basis of $`𝔥`$ and set $`\omega =\frac{1}{2}_{i=1}^dh_i(1)^2\mathrm{𝟏}`$. Then $`(M(1),Y,\mathrm{𝟏},\omega )`$ is a vertex operator algebra with a vacuum $`\mathrm{𝟏}`$ and Virasoro element $`\omega `$ (see \[FLM\]). In particular, $$M(1)=\underset{n0}{}M(1)_n$$ where $`M(1)_n=a_1(n_1)\mathrm{}a_k(n_k)\mathrm{𝟏}\|a_1,\mathrm{},a_k𝔥,n_1,\mathrm{},n_k_+,n_i=n.`$ We will identify $`M(1)_1`$ with $`𝔥`$ in an obvious way. Note that $`o(a)=0`$ for $`a𝔥,`$ so that (2.2.4) simplifies in the free boson case. Moreover, we see from Subsection 2.1 that $`a[0]=a(0)=0.`$ Thus we have the following commutator relation by noting that $`a[1]b=a(1)b=(a,b)`$ for $`a,b𝔥:`$ (3.1.3) $$[a[m],b[n]]=m\delta _{m+n,0}(a,b).$$ It is easy to see from this that $`M(1)_{[n]}`$ is spanned by vectors $`a_1[n_1]\mathrm{}a_k[n_k]\mathrm{𝟏}`$ for $`a_1,\mathrm{},a_k𝔥,n_1,\mathrm{},n_k_+`$ such that $`n_i=n.`$ ### 3.2 Proof of Theorem 1 We first prove part (a) of Theorem 1. It is enough to prove it for an element of the spanning set defined in Section 3.1. Let $`v=a_1[r_1]\mathrm{}a_k[r_k]\mathrm{𝟏}`$ and define the length of $`v`$ to be $`k,`$ denoting this by $`l(v)=k.`$ We prove the statement by induction on $`k.`$ If $`k=0`$ then $`v=\mathrm{𝟏},`$ $`Z(v,q)=1/\eta (q)^d`$ and there is nothing to prove. Now we assume that the statement holds for $`v`$ with $`l(v)k,`$ and consider $`v=a[r]a_1[r_1]\mathrm{}a_k[r_k]\mathrm{𝟏}=a[r]b`$ where $`b=a_1[r_1]\mathrm{}a_k[r_k]\mathrm{𝟏}.`$ Since $`o(a)=0`$, we see from equation (2.2.4) that (3.2.1) $$Z(a[r]b,q)=\underset{m>r/2}{\overset{\mathrm{}}{}}h(m,r)E_{2m}(q)Z(a[2mr]b,q).$$ Since $`2mr>0,`$ $`a[2mr]b`$ is a linear combination of homogeneous vectors whose lengths are no greater than $`k.`$ By induction, each $`Z(a[2mr]b,q)`$ converges to a holomorphic function which is a quotient of a quasi-modular form by $`\eta (q)^d.`$ Note that $`E_{2m}(q)`$ is in $`Q.`$ It is immediate that $`Z(v,q)`$ converges to a holomorphic function of the same kind. In order to prove part (b) of Theorem 1 we need several lemmas. From now on, we fix $`a𝔥`$ such that $`(a,a)=1.`$ ###### Lemma 3.2.1. If $`n>0,`$ $`r0`$ are integers, we have $$Z(a[n]^{2r}\mathrm{𝟏},q)=(1)^{(n+1)r}n^r(2r1)!!h(n,n)^rE_{2n}(q)^r/\eta (q)^d$$ where $`(2r1)!!=135\mathrm{}(2r3)(2r1).`$ ###### Proof. Recall Corollary 2.2.2 and the fact that $`a[s]a[t]=s\delta _{s,t}`$ for positive integers $`s,t`$ (cf. (3.1.3)). We have $`Z(a[n]^{2r}\mathrm{𝟏},q)`$ $`=(1)^{n+1}h(n,n)E_{2n}(q)Z(a[2nn]a[n]^{2r1}\mathrm{𝟏},q)`$ $`=(1)^{n+1}n(2r1)h(n,n)E_{2n}(q)Z(a[n]^{2r2}\mathrm{𝟏},q).`$ Now the result follows by induction, noting that $`Z(\mathrm{𝟏},q)=1/\eta (q)^d.`$ ###### Lemma 3.2.2. Let $`r,s0`$ be integers. Then (3.2.2) $$Z(a[1]^{2r}a[2]^{2s}\mathrm{𝟏},q)=(6)^s(2r1)!!(2s1)!!E_2(q)^rE_4(q)^s/\eta (q)^d$$ and (3.2.3) $$Z(a[2]^{2s}a[3]^{2t}\mathrm{𝟏},q)=(6)^s(2s1)!!(2t1)!!(30)^tE_4(q)^sE_6(q)^t/\eta (q)^d.$$ ###### Proof. Since the proofs of (3.2.2) and (3.2.3) are almost identical we only prove (3.2.2). The special case $`r=0`$ follows from Lemma 3.2.1. So we can proceed by induction on $`r`$ and assume that $`r0.`$ Then by Corollary 2.2.2 and Lemma 3.2.1 we have $`Z(a[1]^{2r}a[2]^{2s}\mathrm{𝟏},q)`$ $`=(2r1)h(1,1)E_2(q)Z(a[1]^{2r2}a[2]^{2s}\mathrm{𝟏},q)`$ $`=(6)^s(2r1)!!(2s1)!!E_2(q)^rE_4(q)^s/\eta (q)^d,`$ as required. ∎ ###### Lemma 3.2.3. Let $`r,s,t0`$ be integers. Then (3.2.4) $$Z(a[1]^{2r}a[2]^{2s}a[3]^{2t}\mathrm{𝟏},q)=c_{r,s,t}E_2(q)^rE_4(q)^sE_6(q)^t/\eta (q)^d+lowerterms$$ where $`c_{r,s,t}`$ is a nonzero constant and “lower terms” means a linear combination of terms of the form $`E_2(q)^r^{}E_4(q)^s^{}E_6(q)^t^{}/\eta (q)^d`$ with $`0r^{}<r`$ if $`r1,`$ and $`0`$ if $`r=0.`$ ###### Proof. If $`r=0`$ this is (3.2.3), so we may assume that $`r1.`$ First we deal with the case $`s=0.`$ Following Corollary 2.2.2 and Lemma 3.2.1 we have $`Z(a[1]^{2r}a[2]^{2s}a[3]^{2t}\mathrm{𝟏},q)`$ $`=(2r1)E_2(q)Z(a[1]^{2r2}a[2]^{2s}a[3]^{2t}\mathrm{𝟏},q)`$ $`+6tE_4(q)Z(a[1]^{2r1}a[2]^{2s}a[3]^{2t1}\mathrm{𝟏},q).`$ Continuing in this fashion, the result follows. ∎ Now we can prove part (b) of Theorem 1. Since the space $`Q`$ of quasi-modular forms on $`SL(2,)`$ is an algebra generated by $`E_2(q),E_4(q)`$ and $`E_6(q),`$ it is enough to prove (b) for $`f(q)=E_2(q)^rE_4(q)^sE_6(q)^t`$ with non-negative integers $`r,s,t.`$ But this follows from Lemmas 3.2.2 and 3.2.3. Indeed, this shows that if $`W`$ is the subspace of $`M(1)`$ spanned by $`a[1]^{2r}a[2]^{2s}a[3]^{2t}\mathrm{𝟏}`$ for $`r,s,t0`$ then the map $`t`$ of (1.3) induces an isomorphism from $`W`$ to $`Q.`$ ## 4 Modularity of trace functions for lattice vertex operator algebras We study trace functions for lattice vertex operator algebras $`V_L`$ (\[B\], \[FLM\]). We prove that $`Z(v,q)`$ is quasi-modular for all $`vV_L.`$ This result together with a result in \[DM2\] yields a proof of Theorem 3. It should be mentioned that in fact, we compute the trace function $`Z_M(v,q)`$ corresponding to an irreducible module $`M`$. It is clear from our proof that each $`Z_M(v,q)`$ is also modular. ### 4.1 Vertex operator algebras associated to lattices We now work in the setting of Chapter 8 of \[FLM\]. In particular, $`L`$ is a positive definite lattice, and $`𝔥=L_{}.`$ Let $`[L]`$ be the group algebra with a basis $`\{e^\alpha \|\alpha L\}.`$ Then the vertex operator algebra $`V_L`$ associated to $`L`$ is $$V_L=M(1)[L]$$ as vector spaces. The vertex operator $`Y(v,z)`$ for $`vM(1)`$ is defined as in the case of $`M(1)`$ (see formulas (3.1.1) and (3.1.2)). The operator $`a(n)`$ for $`a𝔥`$ and $`n0`$ acts on $`V_L`$ via its action on $`M(1).`$ The operator $`a(0)`$ acts on $`V_L`$ by acting on $`[L]`$ in the following way: (4.1.1) $$a(0)e^\alpha =(a,\alpha )e^\alpha $$ for $`\alpha L.`$ We identify $`M(1)`$ with $`M(1)e^0.`$ Then the vacuum of $`V_L`$ is the vacuum $`\mathrm{𝟏}`$ of $`M(1)`$ and the Virasoro element is the same as before. For the purposes of this paper we do not make explicit the expression for vertex operators $`Y(v,z)`$ for general $`v`$ except for those $`vM(1).`$ We remark that $`V_L`$ is a rational vertex operator algebra in the sense of \[DLM\]. Let $`L^{}=\{\alpha L_{}\|(\alpha ,L)\}`$ be the dual lattice of $`L`$ with coset decomposition $`L^{}=_{iL^{}/L}(L+\lambda _i).`$ Then each space $`V_{L+\lambda _i}=M(1)[L+\lambda _i]`$ is an irreducible $`V_L`$-module \[FLM\] where $`[L+\lambda _i]`$ is the subspace of the group algebra $`[L^{}]`$ corresponding the coset $`L+\lambda _i.`$ Moreover $`\{V_{L+\lambda _i}\|iL^{}/L\}`$ constitute the complete set of non-isomorphic irreducible $`V_L`$-modules \[D\]. Since $`V_L`$ also satisfies the $`C_2`$-condition (cf. \[Zh\], \[DLM\]) we have the following modular property for $`V_L`$ by \[Zh\]: ###### Proposition 4.1.1. Let $`v(V_L)_{[n]}`$ and $`\gamma =\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)SL(2,).`$ Then $`Z_i(v,q)`$ converges to a holomorphic function in the upper half plane and there exist scalars $`c_{ij}^\gamma `$ independent of $`v`$ and $`\tau `$ such that $$Z_i(v,\frac{a\tau +b}{c\tau +d})=(c\tau +d)^n\underset{jL^{}/L}{}c_{i,j}^\gamma Z_j(v,\tau )$$ where $`Z_i(v,\tau )=Z_i(v,q)=Z_{V_{L+\lambda _i}}(v,q).`$ ### 4.2 Graded traces in $`V_L`$ In this subsection we determine $`Z_i(v,q)`$ for $`iL^{}/L`$ and $`vV_L,`$ and prove Theorem 2. ###### Lemma 4.2.1. Let $`vM(1)e^\alpha `$ for nonzero $`\alpha L.`$ Then $`Z_i(v,q)=0.`$ ###### Proof. Recall from \[FLM\] the vertex operator $`Y(v,z).`$ It is easy to see that $`v(n)M(1)e^\beta M(1)e^{\alpha +\beta }`$ for $`\beta L+\lambda _i`$ for all $`n.`$ So it is immediate that $`Z_i(v,q)=0.`$ It remains to determine $`Z_i(v,q)`$ for $`vM(1).`$ Note that as $`M(1)`$-module, $`V_{L+\lambda _i}=_{\alpha L+\lambda _i}M(1)e^\alpha `$ and each $`M(1)e^\alpha `$ is an irreducible $`M(1)`$-module. For any $`vM(1)`$ and $`\alpha L^{}`$ we set $$Z_\alpha (v,q)=Z_{M(1)e^\alpha }(v,q).$$ Then $$Z_i(v,q)=\underset{\alpha L+\lambda _i}{}Z_\alpha (v,q).$$ ###### Lemma 4.2.2. Let $`a𝔥`$ such that $`(a,a)=1`$ and $`\alpha L^{}.`$ Then for any non-negative integer $`r`$, $`Z_\alpha (a[1]^r\mathrm{𝟏},q)`$ converges to holomorphic function in the upper half plane. Moreover there exist scalars $`c_{r,r2i}`$ with $`0ir/2`$ and $`c_{r,r}=1`$ independent of $`\alpha `$ and $`a`$ such that $$Z_\alpha (a[1]^r\mathrm{𝟏},q)=\left(\underset{0ir/2}{}c_{r,r2i}(a,\alpha )^{r2i}E_2(q)^i\right)q^{(\alpha ,\alpha )/2}/\eta (q)^d.$$ ###### Proof. We prove this by induction on $`r.`$ If $`r=0`$ the assertion is clear. If $`r=1`$ the assertion follows from the fact that $`o(a[1]\mathrm{𝟏})=(a,\alpha )`$ on $`M(1)e^\alpha .`$ Now we assume that $`r2.`$ Using Corollary 2.2.2 gives $$Z_\alpha (a[1]^r\mathrm{𝟏},q)=(a,\alpha )Z_\alpha (a[1]^{r1}\mathrm{𝟏},q)+(r1)E_2(q)Z_\alpha (a[1]^{r2}\mathrm{𝟏},q).$$ By induction, both $`Z_\alpha (a[1]^{r1}\mathrm{𝟏},q)`$ and $`Z_\alpha (a[1]^{r2}\mathrm{𝟏},q)`$ converge to holomorphic functions in $`H,`$ whence so does $`Z_\alpha (a[1]^r\mathrm{𝟏},q).`$ Indeed, $`Z_\alpha (a[1]^r\mathrm{𝟏},q)`$ $`=(a,\alpha )\left({\displaystyle \underset{0i(r1)/2}{}}c_{r1,r12i}(a,\alpha )^{r12i}E_2(q)^i\right)q^{(\alpha ,\alpha )/2}/\eta (q)^d`$ $`+(r1)E_2(q)\left({\displaystyle \underset{0i(r2)/2}{}}c_{r2,r22i}(a,\alpha )^{r22i}E_2(q)^i\right)q^{(\alpha ,\alpha )/2}/\eta (q)^d`$ $`=\left({\displaystyle \underset{0ir/2}{}}c_{r,r2i}(a,\alpha )^{r2i}E_2(q)^i\right)q^{(\alpha ,\alpha )/2}/\eta (q)^d.`$ It is easy to express $`c_{r,r2i}`$ in terms of the $`c_{r1,r2j}`$ and $`c_{r2,r2k}`$, and so obtain a recursive formula. We leave this to the reader. Since $`c_{r1,r2j}`$ and $`c_{r2,r2k}`$ are independent of $`a,\alpha ,`$ $`c_{r,r2i}`$ is also independent of $`a,\alpha .`$ For convenience, we set $$f_{a,\alpha ,r}(q)=\underset{0ir/2}{}c_{r,r2i}(a,\alpha )^{r2i}E_2(q)^i.$$ Recall that $`\{h_1,\mathrm{},h_d\}`$ is an orthonormal basis of $`𝔥.`$ Let $`r_1,\mathrm{},r_d`$ be non-negative integers. Then by Corollary 2.2.2 and Lemma 4.2.2 we obtain ###### Lemma 4.2.3. Let $`v=h_1[1]^{r_1}\mathrm{}h_d[1]^{r_d}\mathrm{𝟏}.`$ Then $`Z_\alpha (v,q)`$ converges to a holomorphic function in the upper half plane and $$Z_\alpha (v,q)=f_{h_1,\alpha ,r_1}(q)\mathrm{}f_{h_d,\alpha ,r_d}(q)q^{(\alpha ,\alpha )/2}/\eta (q)^d.$$ Recall that $`V_{L+\lambda _i}=_{\alpha L+\lambda _i}M(1)e^\alpha .`$ The following corollary is immediate: ###### Corollary 4.2.4. Let $`v`$ be as before. The function $`Z_i(v,q)`$ is equal to $$(\underset{\alpha L+\lambda _i}{}f_{h_1,\alpha ,r_1}(q)\mathrm{}f_{h_d,\alpha ,r_d}(q)q^{(\alpha ,\alpha )/2})/\eta (q)^d.$$ ###### Theorem 4.2.5. Let $`vM(1).`$ Then $`Z_i(v,q)`$ can be expressed as a finite sum $$Z_i(v,q)=\underset{j}{}f_j(q)\mathrm{\Theta }_{L+\lambda _i,k_j}(a_j,q)/\eta (q)^d$$ where $`f_j(q)Q,`$ $`k_j0,`$ $`a_j𝔥`$ and $$\mathrm{\Theta }_{L+\lambda _i,k_j}(a_j,q)=\underset{\alpha L+\lambda _i}{}(a_j,\alpha )^{k_j}q^{(\alpha ,\alpha )/2}.$$ ###### Proof. First we take $`v=a[n]^{r_n}\mathrm{}a[1]^{r_1}\mathrm{𝟏}`$ for $`a𝔥`$ with $`(a,a)=1.`$ Then from Corollary 2.2.2 we have $$Z_i(v,q)=\underset{s0}{}c_sg_s(q)Z_i(a[1]^s\mathrm{𝟏},q)$$ for some scalars $`c_s`$ and quasi-modular forms $`g_s(q)Q.`$ Now the result follows from Corollary 4.2.4. For arbitrary $`v`$ we can assume that $`v`$ is a monomial in $`h_j(n)`$ for $`j=1,\mathrm{},d`$ and $`n>0.`$ Use the result for $`a[n]^{r_n}\mathrm{}a[1]^{r_1}\mathrm{𝟏}`$ and Corollary 2.2.2 again to see that $`Z_i(v,q)`$ is a linear combination of functions of the form $$f(q)\mathrm{\Theta }_{L+\lambda _i,t_1,\mathrm{},t_d}(h_1,\mathrm{},h_d,q)/\eta (q)^d$$ where $`f(q)`$ belong to $`Q,`$ $`t_j`$ are non-negative integers and $$\mathrm{\Theta }_{L+\lambda _i,t_1,\mathrm{},t_d}(h_1,\mathrm{},h_d,q)=\underset{\alpha L+\lambda _i}{}(h_1,\alpha )^{t_1}\mathrm{}(h_d,\alpha )^{t_d}q^{(\alpha ,\alpha )/2}.$$ It is easy to see that $`\mathrm{\Theta }_{L+\lambda _i,t_1,\mathrm{}t_d}(h_1,\mathrm{},h_d,q)`$ is a linear combination of $`\mathrm{\Theta }_{L+\lambda _i,t_1+\mathrm{}+t_d}(b,q)`$ for $`b𝔥.`$ This completes the proof of the theorem. ∎ We now concentrate on the case $`i=0.`$ Note that $`Z_0(v,q)=Z(v,q).`$ We prove ###### Theorem 4.2.6. Let $`Q(N,ϵ)`$ be as in (1.1) in the introduction. For each $`v`$ in $`V_L,`$ the function $`Z(v,q)\eta (q)^d`$ lies in $`Q(N,ϵ).`$ ###### Proof. After theorem 4.2.5 it suffices to show that if $`a`$ in $`𝔥`$ satisfies $`(a,a)=1`$ then the function $`\mathrm{\Theta }_{L,k}(a,q)`$ lies in $`Q(N,ϵ).`$ This is essentially established in the paper \[DM2\]. In the notation of \[DM2\], the present function $`\mathrm{\Theta }_{L,k}(a,q)`$ is denoted $`\theta (Q,a,k,\tau )`$ where $`Q`$ is the integral quadratic form which corresponds to $`L.`$ If $`k`$ is odd then the function is identically zero. If $`k=2l`$ is even then Theorem 2 of (loc cit) says that the function (4.2.1) $$\psi (Q,a,2l,\tau )=\underset{t=0}{\overset{l}{}}\gamma (t,2l)E_2(\tau )^t\theta (Q,a,2l2t,\tau )$$ is a holomorphic modular form in $`M(N,ϵ).`$ In (4.2.1), $`\gamma (t,2l)`$ is a certain non-zero constant equal to $`1`$ when $`t=0.`$ Moreover $`\theta (Q,a,0,\tau )`$ is just the theta function of $`L.`$ Using this information, it follows from (4.2.1) and induction on $`l`$ that $`\theta (Q,a,2l,\tau )`$ indeed lies in $`M(N,ϵ),`$ as required. ∎ We are now in a position to complete the proof of Theorem 2. In Theorem 4.2.6 we have shown that each trace function $`Z(v,q)`$ satisfies (1.4) with $`f(v,q)`$ some element of $`Q(N,ϵ).`$ We have to show that in fact $`f(v,q)`$ lies in $`M(N,ϵ).`$ First note that we may assume without loss that $`v`$ lies in $`(V_L)_{[k]}`$ for some integer $`k`$ (cf. section 2.1). This puts us in a position to apply Proposition 4.1.1, which tell us that $`Z(v,q)`$ is a generalized modular form in the sense of the introduction. Thus the following result completes the proof of Theorem 2: ###### Proposition 4.2.7. An element of $`Q(N,ϵ)`$ is a modular form of weight $`k`$ if, and only if, it is a generalized modular form of weight $`k.`$ ###### Proof. It is enough to prove sufficiency. Let $`f(\tau )Q(N,ϵ),`$ of weight $`k,`$ be expressed in the form (4.2.2) $$f(\tau )=\underset{i=0}{\overset{m}{}}f_i(\tau )E_2(\tau )^i$$ with each $`f_i(\tau )`$ a form in $`M(N,ϵ)`$ of weight $`k2i.`$ As $`f(\tau )`$ is also a generalized modular form of weight $`k,`$ each of its $`\gamma `$-transforms ($`\gamma SL(2,)`$) has a $`q`$-expansion. In particular, the $`S`$-transform of $`f(\tau )`$ yields an equality of the shape (4.2.3) $$\tau ^kf(S\tau )=\underset{n}{}b(n)q^{n/N}.$$ The $`\gamma `$-transform of $`E_2(\tau )`$ is well-known. In particular, we have (4.2.4) $$E_2(S\tau )=\tau ^2E_2(\tau )\tau /2\pi i.$$ Comparing (4.2.2) - (4.2.4) we obtain an equality (4.2.5) $$\underset{n}{}b(n)q^{n/N}=\underset{i=0}{\overset{m}{}}\tau ^{k+2i}f_i(S\tau )(E_2(\tau )1/\mathrm{log}q)^i.$$ Since each $`\tau ^{k+2i}f_i(S\tau )`$ has a $`q`$-expansion, an equality like (4.2.5) can only hold if all $`f_i(\tau )`$ are zero for $`i>0.`$ Thus $`f(\tau )=f_0(\tau )`$ is indeed modular. ∎ This completes the proof of Proposition 4.2.7, and hence that of Theorem 2, with the following caveat: in Theorem 4.2.6 and Proposition 4.2.7 we were implicitly dealing with modular forms of integral weight. But it is easy to see that these results and those of \[DM2\] extend to the half-integral case i.e., the case in which the lattice $`L`$ has odd rank. We leave the details to the interested reader. ## 5 Spherical harmonics and trace functions In this section we first discuss the relation between spherical harmonics and highest weight vectors for the Virasoro algebra. Each spherical harmonic $`P`$ of degree $`k`$ gives rise to a highest weight vector $`v_P`$ of weight $`k`$ in an canonical way. We then prove Theorem 3. ### 5.1 Spherical harmonics and primary vectors We continue our discussion of the vertex operator algebra $`V_L.`$ In particular, $`M(1)`$ is a vertex operator subalgebra of $`V_L.`$ An element in $`V_L`$ is called a primary state (or singular vector or highest weight vector) if it satisfies $`L(n)v=0`$ for all $`n_+`$. Thanks to the Virasoro commutator relations, this is equivalent to $`L(1)v=L(2)v=0`$. ###### Lemma 5.1.1. Let $`vM(1)`$ be a polynomial in the variables $`h_i(1),(1id)`$. Then $`v`$ is quasiprimary, i.e., $`L(1)v=0`$. ###### Proof. Recall that $`\{h_1,\mathrm{},h_d\}`$ is an orthonormal basis of $`𝔥.`$ Note that $$L(1)=\frac{1}{2}\underset{i=1}{\overset{d}{}}\underset{k}{}^{_{}}_{^{}}h_i(1k)h_i(k)^{_{}}_{^{}}=\underset{i=1}{\overset{d}{}}\underset{k1}{}h_i(1k)h_i(k)$$ and both $`h_i(k)`$ for $`k2`$ and $`h_i(0)`$ annihilate $`v`$. Since all summands of $`L(1)`$ contain one of these as a factor, we immediately see that $`L(1)v=0`$. ∎ ###### Lemma 5.1.2. Let $`vM(1)`$ be a polynomial in the variables $`h_i(1),(1id)`$. Then $`v`$ is primary if, and only if, $`v`$ is a spherical harmonic, i.e., $`\mathrm{\Delta }v=0`$ where $$\mathrm{\Delta }=\underset{i=1}{\overset{d}{}}\frac{^2}{h_i(1)^2}.$$ ###### Proof. It suffices to consider the condition $`L(2)v=0`$. Since $$L(2)=\frac{1}{2}\underset{i=1}{\overset{d}{}}h_i(1)^2+\underset{i=1}{\overset{d}{}}\underset{k2}{}h_i(2k)h_i(k)$$ and $`h_i(k)`$ annihilates $`v`$ if $`k2,`$ we see that $`L(2)v=0`$ is equivalent to $`_{i=1}^dh_i(1)^2v=0`$. Now since $`h_i(1)h_i(1)^n\mathrm{𝟏}=nh_i(1)^{n1}\mathrm{𝟏}`$, $`h_i(1)`$ acts as $`/h_i(1)`$ on $`M(1)`$. ∎ ### 5.2 Proof of Theorem 3 Let $`P`$ be a spherical harmonic of degree $`k`$: $$P=P(x_1,\mathrm{},x_d),\mathrm{\Delta }P=0$$ where $`\mathrm{\Delta }=_{i=1}^d^2/x_i^2`$. Let $`v_P`$ be the corresponding primary state in $`M(1)`$, i.e., $$v_P=P(h_1(1),\mathrm{},h_d(1))\mathrm{𝟏}.$$ We note that $`o(v_P)=v_P(k1).`$ Now let $`P=_Ic_Ix^I`$ where $`c_I`$ and $`x^I=x_1^{a_1}\mathrm{}x_d^{a_d}`$ and where $`I=(a_1,\mathrm{},a_d)`$ is some ordered tuple of non-negative integers such that $`a_i=k`$. Then from Section 3.1, $$Y(v_P,z)=\underset{I}{}c_IY(v_I,z)=\underset{I}{}c_I^{_{}}_{^{}}h_I(z)^{_{}}_{^{}}$$ where $`h_I(z)=_{i=1}^dh_i(z)^{a_i}`$ and $`h_i(z)=_nh_i(n)z^{n1}`$. Therefore (5.2.1) $$o(v_P)=\underset{I}{}c_I^{_{}}_{^{}}h_1(n_{11})\mathrm{}h_1(n_{1a_1})h_2(n_{21})\mathrm{}h_2(n_{2a_2})\mathrm{}h_d(n_{d1})\mathrm{}h_1(n_{da_d})^{_{}}_{^{}}$$ where the indices in the second sum range over all integers such that $`_{i=1}^d_{j=1}^{a_i}n_{ij}=0`$. In particular, this includes the case with all $`n_{ij}=0`$: $$\sigma =\underset{J}{}c_Jh_1(0)^{a_1}\mathrm{}h_d(0)^{a_d}=P(h_1(0),\mathrm{},h_d(0)).$$ As $`h_i(0)`$ operates on $`ue^\alpha `$ as in (4.1.1), we see that $$\sigma :ue^\alpha \left(\underset{J}{}c_J(h_1,\alpha )^{a_1}\mathrm{}(h_d,\alpha )^{a_d}\right)ue^\alpha .$$ But this is exactly $`P(\alpha )ue^\alpha `$. It follows that $$\mathrm{tr}_{V_L}\sigma q^{L_0\frac{d}{24}}=\underset{\alpha L}{}P(\alpha )q^{\frac{(\alpha ,\alpha )}{2}}/\eta (\tau )^d=\frac{\theta _L(P,\tau )}{\eta (\tau )^d}.$$ To prove the theorem, it remains to show that all other terms in (5.2.1) have trace $`0`$ in their action on $`V_L.`$ To see this, we adopt a different approach and note that to prove the theorem it is enough to establish it for a spanning set of spherical harmonics of degree $`k`$. These are given, for example \[H\], by $$P=(t_1x_1+\mathrm{}+t_dx_d)^k$$ for $`(t_1,\mathrm{},t_d)^d`$ such that $`_{i=1}^dt_i^2=0`$. Set $`h_0=_{i=1}^dt_ih_i𝔥`$, and note that $`(h_0,h_0)=_{i=1}^dt_i^2=0`$. Then $$Y(v_P,z)=^{_{}}_{^{}}(\underset{i=1}{\overset{d}{}}t_ih_i(z))^k^{_{}}_{^{}}=^{_{}}_{^{}}h_0(z)^k^{_{}}_{^{}}.$$ Clearly $$o(v_P)=h_0(0)^k+\underset{n0}{}\left(\genfrac{}{}{0pt}{}{k}{2}\right)h_0(0)^{k2}h_0(n)h_0(n)+\mathrm{}.$$ Then it is enough to show that each operator of the form (5.2.2) $$U=h_0(n_1)h_0(n_2)\mathrm{}h_0(n_r),(r2,notalln_iarezero,n_i=0)$$ has trace $`0`$ on $`M(1)`$. Let us take a basis $`h_0,k_0,k_1,\mathrm{},k_{d2}`$ of $`𝔥`$ with $$(h_0,k_0)=1,(h_0,k_i)=0,(1id2).$$ Then $`h_0(m)`$ commutes with all $`h_0(m),k_i(m),(1id,m)`$. Now consider a vector in $`M(1)`$ of the form $`pv`$ where $`p`$ is a polynomial in $`k_0(m),m_+`$ and $`v`$ does not involve any mode $`k_0(m),m`$. Since some $`n_i`$ occurring in (5.2.2) is positive, we see that the operator $`U`$ maps $`pv`$ to $`0`$ or $`p^{}v^{}`$ where $`p^{}`$ has lower degree than $`p`$. Therefore the trace of $`U`$ on $`M(1)`$ is $`0`$ as the operator $`U`$ never has non-zero eigenvalues. This completes the proof of Theorem 3. As explained in the introduction, the result of Waldspurger \[W\] together with Theorem 3 implies the ###### Corollary 5.2.1. Let $`f(q)`$ be any cusp-form of level $`1.`$ There there is a primary field $`v`$ in the lattice VOA $`V_{E_8}`$ such that $`Z(v,q)=f(q)/\eta (q)^8.`$
warning/0004/math0004158.html	ar5iv	text	# On the relation between the A-polynomial and the Jones polynomial ## 1. Introduction In 1984, V. Jones introduced a polynomial invariant of knots \[J\] through skein relations. Another version of this invariant, the Kauffman bracket \[K\], was introduced shortly after. Colored versions of these invariants were defined, via quantum groups \[RT\], and via Jones-Wenzl idempotents \[L\], \[BHMV\]. In 1993, Cooper, Culler, Gillet, Long, and Shalen defined a two variable polynomial invariant of knots, the A-polynomial, using the character variety of $`SL(2,)`$-representations of the fundamental group of the knot complement. This invariant was generalized in \[FGL\] to a finitely generated ideal of polynomials in the quantum plane. The construction is done in the context of skein modules, and is based on the fact that the Kauffman bracket skein modules represent deformations of function rings on character varieties \[B\], \[PS\], and on the relationship between the skein algebra of the cylinder over a torus and the noncommutative torus \[FG\]. As shown in \[FGL\], each element in the noncommutative A-ideal defines a matrix that annihilates the vector whose entries are the colored Jones polynomials of the knot (or, more precisely, the colored Kauffman brackets of the knot; they differ from the colored Jones polynomials by the change of variable $`tit`$). The orthogonality between the rows of the matrix and the vector whose entries are the colored Jones polynomials of the knot has been called the “orthogonality relation (between the Jones polynomial and the A-polynomial)”. In the present paper it is shown that the noncommutative A-ideal together with a finite number (depending on the A-ideal of the knot) of colored Kauffman brackets of the knot determine all other colored Kauffman brackets of the knot. Also, it is shown that, under certain technical conditions on the A-ideal, the noncommutative A-ideal determines all colored Kauffman brackets of the knot. As an example, any knot having the same A-ideal as the unknot, respectively trefoil knot, has the same colored Kauffman brackets as the unknot, respectively trefoil knot. ## 2. The action of $`K_t(𝕋^2\times I)`$ on $`K_t(𝔻^2\times I)`$ The Kauffman bracket skein module of the three manifold $`M`$ is defined in the following way. Let $`[t,t^1]`$ be the $`[t,t^1]`$-module freely spanned by the isotopy classes of framed links in $`M`$ including the empty link, and let $`𝒮`$ be the submodule spanned by the relations $`\text{ }\text{}\text{ }t\text{ }\text{}\text{ }t^1\text{ }\text{}\text{ }`$ and $`+t^2+t^2`$. The Kauffman bracket skein module of $`M`$ is $`K_t(M)=[t]/𝒮`$. In the case where $`M`$ is the cylinder over a surface, $`K_t(M)`$ has a natural algebra structure, with product defined by placing one link on top of another. If $`M`$ is a manifold with boundary, the operation of gluing a cylinder to the boundary induces a $`K_t(M\times I)`$-module structure on $`K_t(M)`$. As an example it is known that the Kauffman bracket skein algebra of the cylinder over an annulus (i.e., that of the solid torus), is $`[t,t^1,\alpha ]`$, where $`\alpha `$ is the curve that runs once around the annulus and has framing parallel to the annulus. Another, more complicated example is that of the Kauffman bracket skein algebra of $`K_t(𝕋^2\times I)`$. Its multiplication rule and action on the skein module of the solid torus are described by means of two families of Chebyshev polynomials, $`\{T_n\}_n`$ defined by $`T_0=2,T_1=x,`$ $`T_{n+1}=xT_nT_{n1}`$ for $`n`$ and $`\{S_n\}_n`$ defined by $`S_0=2,S_1=x,`$ $`S_{n+1}=xS_nS_{n1}`$ for $`n`$. Let $`p`$ and $`q`$ be two integers with $`p=np^{}`$, $`q=nq^{}`$, $`p^{}`$, $`q^{}`$ coprime. We define $`(p,q)_T=T_n((p^{},q^{}))`$, where $`(p^{},q^{})`$ is the corresponding curve on the torus, with framing parallel to the torus, and its powers are defined by parallel copies. The elements $`(p,q)_T`$, $`p0,q`$ span $`K_t(M)`$ as a $`[t,t^1]`$-module. In \[FG\] we proved the following product-to-sum formula $$(p,q)_T(r,s)_T=t^{\|_{rs}^{pq}\|}(p+r,q+s)_T+t^{\|_{rs}^{pq}\|}(pr,qs)_T.$$ As a consequence of this formula, the Kauffman bracket skein algebra of the cylinder over a torus is isomorphic to the subalgebra of the noncommutative torus generated by noncommutative cosines. Let us recall that the algebra of trigonometric polynomials in the noncommutative torus is $`_t[l,l^1,m,m^1]`$, with multiplication $``$, satisfying $`lm=t^2ml`$. The elements $`e_{p,q}=t^{pq}l^pm^q`$ are the noncommutative exponentials; they satisfy $`e_{p,q}e_{r,s}=t^{\|_{rs}^{pq}\|}e_{p+r,q+s}.`$ The noncommutative cosines are $`\frac{1}{2}(e_{p,q}+e_{p,q})`$. The map $`(p,q)_Te_{p,q}+e_{p,q}`$ gives the isomorphism between $`K_t(𝕋^2\times I)`$ and the algebra of noncommutative cosines. Let $`K`$ be a knot in $`S^3`$, and $`M`$ the complement of a regular neighborhood of $`K`$. Recall the left action of $`K_t(𝕋^2\times I)`$ on $`K_t(M)`$. The peripheral ideal of $`K`$ is the left ideal of $`K_t(𝕋^2\times I)`$ which annihilates the empty link. The noncommutative A-ideal of $`K`$, denoted by $`𝒜_t(K)`$ is the left ideal obtained by extending $`I_t(K)`$ to $`_t[l,l^1,m,m^1]`$ then contracting it to $`_t[l,m]`$. As explained in \[FGL\], this is a noncommutative generalization of the A-polynomial. The A-polynomial is obtained by setting $`t=1`$, replacing $`l`$ and $`m`$ by $`l`$ and $`m`$ and taking the generator of the radical of the one-dimensional part of the A-ideal (divided by $`(l1)`$). There is a left and a right action of $`K_t(𝕋^2\times I)`$ on $`K_t(𝔻^2\times S^1)`$, one for the positive, the other one for the negative orientation of the boundary torus. To understand them, let us denote by $`x_{p,q}`$ the image in $`K_t(𝔻^2\times S^1)`$ of $`(p,q)_T`$ on the boundary torus (with the positive orientation). It is not hard to see that $`x_{0,q}=(t^2)^q+(t^2)^q`$ and the product-to-sum formula yields $`x_{p+1,q}=t^q(1,0)x_{p,q}t^{2q}x_{p1,q}.`$ The second order recurrence relation for $`t^{pq}x_{p,q}`$ has fixed coefficients, and hence a formula for the general term can be found. It is $`x_{p,q}=t^{pq}((t^2)^qS_p(\alpha )(t^2)^qS_{p2}(\alpha )).`$ Lifting the skeins $`T_n(\alpha )`$ to the boundary torus and using the product-to-sum formula we get the following ###### Lemma 1. The left action is described by $`(p,q)_TT_n(\alpha )`$ $`=`$ $`t^{(2n+p)q}[(t^2)^qS_{n+p}(\alpha )(t^2)^qS_{n+p2}(\alpha )]`$ $`+t^{(2np)q}[(t^2)^qS_{pn}(\alpha )(t^2)^qS_{pn2}(\alpha )]`$ while the right action is given by $`T_n(\alpha )(p,q)_T`$ $`=`$ $`(p,q)_TT_n(\alpha )`$ $`=`$ $`t^{(2n+p)q}[(t^2)^qS_{p+n}(\alpha )(t^2)^qS_{p+n2}(\alpha )]`$ $`+t^{(2np)q}[(t^2)^qS_{pn}(\alpha )(t^2)^qS_{pn2}(\alpha )].`$ ## 3. The results Gluing a solid torus to the complement $`M`$ of a regular neighborhood of a knot $`K`$, in such a way that the longitude is glued to the longitude and the meridian to the meridian, induces a pairing $`K_t(𝔻\times S^1)\times K_t(M)[t,t^1].`$ The basis $`\{S_n(\alpha )\}_n`$ induces a family of functionals $`<S_n(\alpha ),>`$, $`n=0,1,2,\mathrm{}`$. If we denote by $`\mathrm{}`$ the empty link, then $`<S_n(\alpha ),\mathrm{}>=\kappa _n(K),`$ where $`\kappa _n(K)`$ is the $`n`$th colored Kauffman bracket of $`K`$ with zero framing \[L\], \[T\] (the $`n`$th colored Kauffman bracket is a “twisted” version of the $`n`$th colored Jones polynomial as defined in \[RT\]). Indeed, the recurrence relation for $`S_n`$ shows that the link in $`S^3`$ obtained from the pairing is $`K`$ colored by the Jones-Wenzl idempotent. The pairing is compatible with the actions of $`K_t(𝕋^2\times I)`$ on both modules, i.e. $`<u(p,q)_T,v>=<u,(p,q)_Tv>`$ for any skeins $`u`$ and $`v`$. In particular, if $`a`$ is in the peripheral ideal $`I_t(K)`$ of $`K`$, then $`<ua,\mathrm{}>=0`$. So, if $`u=T_n(\alpha )`$, and $`a=_ic_i(p_i,q_i)_T`$, then by Lemma 1, $`<aT_n(\alpha ),\mathrm{}>={\displaystyle \underset{i}{}}c_i(t^{p_iq_i}(t^{(2n+p_i)q_i}[(t^2)^{q_i}<S_{p_i+n}(\alpha ),\mathrm{}>`$ $`(t^2)^{q_i}<S_{p_i+n2}(\alpha ),\mathrm{}>]`$ $`+t^{(2np_i)q_i}[(t^2)^{q_i}<S_{p_in}(\alpha ),\mathrm{}>(t^2)^{q_i}<S_{p_in2}(\alpha ),\mathrm{}>])`$ $`=`$ $`{\displaystyle \underset{i}{}}c_i(t^{(2n+p_i)q_i}[(t^2)^{q_i}\kappa _{p_i+n}(K)(t^2)^{q_i}\kappa _{p_i+n2}(K)]`$ $`+t^{(2np_i)q_i}[(t^2)^{q_i}\kappa _{p_in}(K)(t^2)^{q_i}\kappa _{p_in2}(K)]).`$ This relation has been called the orthogonality relation in \[FGL\] since it expresses the orthogonality between the vector with entries equal to the colored Kauffman brackets of the knot and the rows of the matrix of the linear transformation induced by $`a`$ between the module $`K_t(𝔻\times I)`$ with basis $`\{T_n(\alpha )\}_n`$ and the same module with basis $`\{S_n(\alpha )\}_n`$. Since $`a`$ arises from an element in the noncommutative A-ideal (through an extension and a contraction), orthogonality expresses a relationship between the the elements of the A-ideal and the vector whose entries are the colored Kauffman brackets. ###### Theorem 1. For every knot $`K`$ there is a number $`\nu (K)`$ such that if $`K^{}`$ is a knot with $`𝒜_t(K)=𝒜_t(K^{})`$ and $`\kappa _j(K)=\kappa _j(K^{})`$ for $`j=1,2,\mathrm{},\nu (K)`$, then $`\kappa _j(K)=\kappa _j(K^{})`$ for all $`j`$. Moreover, $`\nu (K)`$ depends only on the A-ideal of $`K`$. ###### Proof. Choose $`a=_jc_j(p_j,q_j)_T`$ some element in $`I_t(K)`$, let $`p`$ be the maximum of $`p_j`$ and assume $`p_j=p`$ if $`j=1,2,\mathrm{},m`$, $`p_jp`$ if $`j>m`$. Then, coefficient of $`\kappa _{n+p}`$ in the orthogonality relation written for $`a`$ is $`{\displaystyle \underset{j=1}{\overset{m}{}}}c_j(1)^{q_j}t^{(2n+2+p)q_j}`$ Since the $`q_j`$ appearing in this expression are distinct (the $`p_i`$’s being the same), this expression is identically equal to zero only for finitely many $`n`$. Hence the orthogonality relation provides a recurrence relation that determines uniquely $`\kappa _n`$ for large $`n`$. ∎ As the result below shows, in certain situations the A-ideal determines the colored Kauffman brackets of the knot. ###### Theorem 2. Assume that $`K`$ is a knot with the property that $`𝒜_t(K)`$ contains a polynomial $`_{p,q}\gamma _{p,q}l^pm^q`$ of degree $`2`$ in $`l`$ such that there exists no $`n0`$ for which the expression $`_q\gamma _{2,q}(1)^qt^{(2n+2)q}`$ is identically equal to zero. Then for any knot $`K^{}`$ with the property that $`𝒜_t(K)=𝒜_t(K^{})`$, it follows that $`\kappa _n(K)=\kappa _n(K^{})`$ for all $`n=1,2,3,\mathrm{}`$. ###### Proof. The polynomial gives rise to an element $`a=_ic_i(1,q_i)_T+u`$ in $`I_t(K)`$, with $`c_i=t^{q_i}\gamma _{2,q_i}`$ and $`u`$ a polynomial in $`(0,1)`$. By Lemma 1, $`T_n(\alpha )u`$ is of the form $`\lambda S_n(\alpha )+\mu S_{n2}(\alpha )`$, $`\lambda ,\mu [t,t^1]`$. On the other hand, the same lemma shows that $`T_n(\alpha ){\displaystyle \underset{i}{}}c_i(1,q_i)_T={\displaystyle \underset{i}{}}c_i[(t)^{(2n+3)q_i}S_{n+1}(\alpha )`$ $`(t)^{(2n1)q_i}S_{n1}(\alpha )+(t)^{(2n+3)}S_{1n}(\alpha )(t)^{(2n1)q_i}S_{n1}(\alpha )].`$ Since $`S_k=S_{k2}`$ for all $`k`$, this is further equal to $`{\displaystyle \underset{i}{}}c_i[(t)^{(2n+3)q_i}S_{n+1}(\alpha )[(t)^{(2n1)q_i}+(t)^{(2n1)q_i}]S_{n1}(\alpha )`$ $`+(t)^{(2n+3)q_i}S_{n3}(\alpha )].`$ Hence the orthogonality relation applied to $`a`$ yields a $`[t,t^1]`$-linear equation in $`\kappa _{n+1}(K),\kappa _n(K),\kappa _{n1}(K),\kappa _{n2}(K)`$, and $`\kappa _{n3}(K)`$. For $`n1`$, the coefficient of $`\kappa _{n+1}(K)`$ is $`_ic_i(t)^{(2n+3)q_i}`$, and the condition from the statement translates to the fact that for no $`n`$ this is identically equal to zero. Therefore, the orthogonality relation provides a recursion that uniquely determines $`\kappa _n(K)`$ from $`\kappa _1(K)`$. For $`n=0`$, since $`\kappa _0(K)=1`$, $`\kappa _1(K)=0`$, $`\kappa _2(K)=1`$, $`\kappa _3(K)=\kappa _1(K)`$, the orthogonality relation gives a linear equation in $`\kappa _1(K)`$. The coefficient of $`\kappa _1(K)`$ is $`2_ic_i(t)^{3q_i}`$. Again this is not equal to zero. So the equation can be solved uniquely for $`\kappa _1(K)`$. It follows that the A-ideal determines the colored Kauffman brackets of the knot, and we are done.∎ Observe that the degree in $`l`$ of any polynomial in $`𝒜_t(K)`$ is at least $`2`$. ## 4. Examples ### 4.1. The unknot The A-ideal of the unknot is generated by $`(l+t^2)(l+t^2)`$ and $`lm^2(l+t^2)+t^2(l+t^2)`$ \[FGL\]; hence it satisfies the conditions in Theorem 2. The orthogonality relation for $`(l+t^2)(l+t^2)`$, that is, for $`(1,0)_T+t^2+t^2I_t(K)`$, gives $`\kappa _0(K)=1,\kappa _1(K)=t^2t^2,`$ $`\kappa _{n+1}(K)=(t^2t^2)\kappa _n(K)\kappa _{n1}(K),n1.`$ From this we obtain the well known formula $`\kappa _n(K)=(1)^n(t^{2n+2}t^{2n2})/(t^2t^2).`$ The orthogonality relation for the other element leads to a different recurrence relation with the same solution. ### 4.2. The trefoil The A-ideal of the left-handed trefoil is generated by $`[m^4(l+t^{10})t^4(l+t^2)](lt^6m^6)`$, $`(l+t^{24})(l+t^{10})(l+t^2)(lt^6m^6)`$ and $`(m^2t^{22})(l+t^{10})(l+t^2)(lt^6m^6)`$ \[G\]. A quick look at the element $`[m^4(l+t^{10})t^4(l+t^2)](lm^6t^6)`$ shows that the conditions in the statement of Theorem 2 are fulfilled. This element corresponds to $`(1,5)_Tt^8(1,1)_T+t^3(0,5)_Tt(0,1)_T`$ in the peripheral ideal. The orthogonality relation produces the following recursion $`(t^{10n15}+t^{2n11})\kappa _{n+1}(K)+(t^{10n+7}t^{10n13}+t^{2n+3}`$ $`+t^{2n1})\kappa _n(K)+(t^{10n+5}t^{10n+5}t^{2n7}+t^{2n7})\kappa _{n1}(K)+(t^{10n13}`$ $`+t^{10n+7}t^{2n1}t^{2n+3})\kappa _{n2}(K)+(t^{10n15}t^{2n11})\kappa _{n3}(K)=0.`$ In particular, for $`n=0`$, $`(t^{11}t^{15})\kappa _1(K)t^7t^{13}+t^3+t^1=0,`$ and hence $`\kappa _1(K)=t^{18}t^{10}t^6t^2`$, the well known formula for the Kauffman bracket of the trefoil knot with framing zero.
warning/0004/quant-ph0004101.html	ar5iv	text	# Effects of Local Fields on Spontaneous Emission in Dielectric Media ## Abstract The local-field renormalization of the spontaneous emission rate in a dielectric is explicitly obtained from a fully microscopic quantum-electrodynamical, many-body derivation of Langevin–Bloch operator equations for two-level atoms embedded in an absorptive and dispersive, linear dielectric host. We find that the dielectric local-field enhancement of the spontaneous emission rate is smaller than indicated by previous studies. In the formative period of nonlinear optics, Bloembergen taught us that the nonlinear optical effects of a dilute collection of atoms that are embedded in a dielectric host are enhanced by local-field effects . Now, in the era of quantum optics, researchers are again looking at the interaction of dielectric materials, the radiation field, and resonant atoms. Central to these investigations are efforts to quantize the electromagnetic field in dielectrics. One widely used technique is to quantize the macroscopic Maxwell equations in which the classical constitutive relations have been assumed . In the microscopic approach, a generalized Hopfield transformation, based on Fano diagonalization, is used to obtain the coupled polariton modes of the coupled field–oscillator system . These quantization methods are well-established for dielectrics with negligible absorption and techniques to deal with the special requirements of quantizing the electromagnetic field in absorbing dielectrics are beginning to emerge . Since Purcell first predicted the alteration of the emission rate of an excited atom due to an optical cavity , it has become well known that the observed spontaneous emission rate of an atom depends on its environment. When the quantized coupled field–dielectric theories are applied to the spontaneous emission of a two-level atom embedded in an absorptionless dielectric, the relation $$\mathrm{\Gamma }_{SE}^{diel}=n\mathrm{}^2\mathrm{\Gamma }_0$$ (1) is obtained . Here, $`n`$ is the linear index of refraction and $`\mathrm{}`$ is the dielectric local-field enhancement factor, $`\mathrm{\Gamma }_0`$ is the vacuum spontaneous emission rate, and $`\mathrm{\Gamma }_{SE}^{diel}`$ is the enhanced spontaneous emission rate in the dielectric. Both the Lorentz virtual-cavity model $`\mathrm{}=(n^2+2)/3`$ and the Onsanger real-cavity model $`\mathrm{}=3n^2/(2n^2+1)`$ have been utilized in various studies of local-field effects on spontaneous emission. One of the key features of these approaches of applying the quantization of fields in dielectrics to spontaneous emission is that the oscillators that comprise the dielectric host are assumed to be unaffected by the presence of the embedded atom. The dielectric medium, as well as the field, is treated as a local condition at the site of a resonant atom such that the atom interacts with an all-pervasive, nonlocal, quantized effective field, the vacuum polariton modes, rather than the local vacuum radiation field modes and the oscillators. Because the near-dipole–dipole interaction is the fundamental interaction underlying the Lorentz local field, a many-body approach that explicitly deals with the vacuum radiation field modes and the interactions of the atom with the nearby polarizable particles of the host is clearly needed to accurately evaluate the effects of local fields on spontaneous emission in dielectric media. In this letter, we develop Langevin–Bloch operator equations of motion for a dense collection of two-level atoms embedded in a dielectric host medium. The Heisenberg picture is used since this provides the most direct correspondence between the operator equations and the Bloch equations. We begin the development from a fully microscopic many-body viewpoint in which the material is treated as a disordered mixture of two different species of two-level systems and derive the Heisenburg equations of motion for the material and field mode operators. Adiabatically eliminating the variables associated with the quantized field modes results in coupled equations of motion for the material variables. We take the harmonic oscillator limit for one species by assuming that its resonance frequency is sufficiently detuned from the primary species that the atoms remain in the ground state. Adiabatically eliminating the harmonic oscillators results in a Langevin-Bloch formulation for two-level systems embedded in a dielectric host that exhibits local-field renormalization of the fluctuations, the near-dipole–dipole (NDD) interaction of the two-level atoms, the radiation field, the dephasing rate, and the population decay rate. We obtain $$\mathrm{\Gamma }_{SE}^{diel}=Re(\mathrm{})\mathrm{\Gamma }_0$$ (2) for the renormalized spontaneous emission rate. In our many-body derivation, the dielectric local-field enhancement factor $`\mathrm{}=(n^2+2)/3`$ arises from the interaction of the embedded atom with the nearby polarizable particles of the host via the electromagnetic field. The linear index of refraction is complex and frequency dependent and properly accounts for the frequency dispersion and absorption of the dielectric. We consider a disordered mixture of two species, $`a`$ and $`b`$, of two-level systems. The two-level systems are coupled only via the electromagnetic field. We allow for the possibility of an externally applied probe or driving field that is taken to be in a coherent state. The constituents of the total Hamiltonian are the Hamiltonians for the free atoms of species $`a`$ and $`b`$, the free quantized radiation field, and the interaction of the two-level systems with the quantized electromagnetic field. We have, in the electric-dipole and rotating-wave approximations, $$H=\underset{j}{}\frac{\mathrm{}\omega _a}{2}\sigma _3^j+\underset{n}{}\frac{\mathrm{}\omega _b}{2}\zeta _3^n+\mathrm{}\underset{l,\sigma }{}\omega _la_l^{}a_l$$ $$i\mathrm{}\underset{j}{}\underset{l,\sigma }{}\left(g_l^ja_l\sigma _+^je^{i\stackrel{}{k}_l\stackrel{}{r}_j}g_{l}^{j}{}_{}{}^{}a_l^{}\sigma _{}^je^{i\stackrel{}{k}_l\stackrel{}{r}_j}\right)$$ $$i\mathrm{}\underset{n}{}\underset{l,\sigma }{}\left(h_l^na_l\zeta _+^ne^{i\stackrel{}{k}_l\stackrel{}{r}_n}h_{l}^{n}{}_{}{}^{}a_l^{}\zeta _{}^ne^{i\stackrel{}{k}_l\stackrel{}{r}_n}\right)$$ $$\frac{i\mathrm{}}{2}\underset{j}{}\left(\mathrm{\Omega }_a\sigma _+^je^{i(\omega _pt\stackrel{}{k}_p\stackrel{}{r}_j)}\mathrm{\Omega }_a^{}\sigma _{}^je^{i(\omega _pt\stackrel{}{k}_p\stackrel{}{r}_j)}\right)$$ $$\frac{i\mathrm{}}{2}\underset{n}{}\left(\mathrm{\Omega }_b\zeta _+^ne^{i(\omega _pt\stackrel{}{k}_p\stackrel{}{r}_n)}\mathrm{\Omega }_b^{}\zeta _{}^ne^{i(\omega _pt\stackrel{}{k}_p\stackrel{}{r}_n)}\right),$$ where $`a_l^{}`$ and $`a_l`$ are the creation and destruction operators for the field modes and $`\omega _l`$ is the frequency of the field in the mode $`l`$. The vacuum dispersion relation is used throughout, e.g. $`\stackrel{}{k}_l=\widehat{k}_l\omega _l/c`$, where $`\widehat{k}_l`$ is a unit vector in the direction of $`\stackrel{}{k}_l`$. For atoms of species $`a`$, $`\sigma _3^j`$ is the population inversion operator and $`\sigma _\pm ^j`$ are the raising and lowering operators for the $`j^{th}`$ atom, $`g_l^j=(2\pi \omega _l/\mathrm{}V)^{1/2}\mu _a\widehat{𝐩}_j\widehat{𝐞}_{\stackrel{}{k}_l,\sigma }`$ is the coupling between the atom at position $`\stackrel{}{r}_j`$ and the quantized radiation field, $`\widehat{𝐩}_j`$ is a unit vector in the direction of the dipole moment at $`\stackrel{}{r}_j`$, $`\omega _a`$ is the transition frequency, $`\mu _a`$ is the dipole moment, $`\mathrm{\Omega }_a=\mu _a/\mathrm{}`$ is the Rabi rate, and $``$ is the field envelope associated with the coherent field with carrier frequency $`\omega _p`$. For species $`b`$, $`\zeta _3^n`$, $`\zeta _\pm ^n`$, $`h_l^n`$, $`\stackrel{}{r}_n`$, $`\widehat{𝐩}_n`$, $`\omega _b`$, $`\mu _b`$, and $`\mathrm{\Omega }_b=\mu _b/\mathrm{}`$ respectively perform the same functions. Finally, $`V`$ is the quantization volume, $`\widehat{𝐞}_{\stackrel{}{k}_l,\sigma }`$ is the polarization vector, and $`\sigma `$ denotes the state of polarization. The Heisenberg equations of motion for the material and field mode operators $$\frac{da_l}{dt}=i\omega _la_l+\underset{j}{}g_{l}^{j}{}_{}{}^{}\sigma _{}^je^{i\stackrel{}{k}_l\stackrel{}{r}_j}+\underset{n}{}h_{l}^{n}{}_{}{}^{}\zeta _{}^ne^{i\stackrel{}{k}_l\stackrel{}{r}_n}$$ $$\frac{d\sigma _{}^j}{dt}=i\omega _a\sigma _{}^j+\underset{l,\sigma }{}g_l^j\sigma _3^ja_le^{i\stackrel{}{k}_l\stackrel{}{r}_j}+\frac{1}{2}\sigma _3^j\mathrm{\Omega }_ae^{i(\omega _pt\stackrel{}{k}_p\stackrel{}{r}_j)}$$ $$\frac{d\sigma _3^j}{dt}=2\underset{l,\sigma }{}(g_l^j\sigma _+^ja_le^{i\stackrel{}{k}_l\stackrel{}{r}_j}+g_{l}^{j}{}_{}{}^{}a_l^{}\sigma _{}^je^{i\stackrel{}{k}_l\stackrel{}{r}_j})$$ $$\sigma _+^j\mathrm{\Omega }_ae^{i(\omega _pt\stackrel{}{k}_p\stackrel{}{r}_j)}\mathrm{\Omega }_a^{}\sigma _{}^je^{i(\omega _pt\stackrel{}{k}_p\stackrel{}{r}_j)}$$ $$\frac{d\zeta _{}^n}{dt}=i\omega _b\zeta _{}^n+\underset{l,\sigma }{}h_l^n\zeta _3^na_le^{i\stackrel{}{k}_l\stackrel{}{r}_n}+\frac{1}{2}\zeta _3^n\mathrm{\Omega }_be^{i(\omega _pt\stackrel{}{k}_p\stackrel{}{r}_n)}$$ $$\frac{d\zeta _3^n}{dt}=2\underset{l,\sigma }{}(h_l^n\zeta _+^na_le^{i\stackrel{}{k}_l\stackrel{}{r}_n}+h_{l}^{n}{}_{}{}^{}a_l^{}\zeta _{}^ne^{i\stackrel{}{k}_l\stackrel{}{r}_n})$$ $$\zeta _+^n\mathrm{\Omega }_be^{i(\omega _pt\stackrel{}{k}_p\stackrel{}{r}_n)}\mathrm{\Omega }_b^{}\zeta _{}^ne^{i(\omega _pt\stackrel{}{k}_p\stackrel{}{r}_n)}$$ are derived from $`i\mathrm{}(dO/dt)=[O,H]`$. From this point, we adopt normal ordering in which $`a_l^{}`$ appears to the left of the atomic operators and $`a_l`$ appears to the right. The Heisenberg equations of motion for the material variables are obtained by elimination of the variables associated with the quantized field modes . Then, $$\frac{d\sigma _{}^j}{dt}=i\omega _a\sigma _{}^j(t)+\frac{\mu _a}{2\mathrm{}}\sigma _3^j(t)(t)e^{i(\omega _pt\stackrel{}{k}_p\stackrel{}{r}_j)}$$ $$+\underset{l,\sigma }{}g_l^j\sigma _3^j(t)a_l(0)e^{i(\omega _lt\stackrel{}{k}_l\stackrel{}{r}_j)}$$ $$+\underset{l,\sigma }{}_0^t𝑑t^{}e^{i\omega _l(tt^{})}\underset{i}{}g_{l}^{i}{}_{}{}^{}g_l^j\sigma _3^j(t)\sigma _{}^i(t^{})e^{i\stackrel{}{k}_l(\stackrel{}{r}_i\stackrel{}{r}_j)}$$ $$+\underset{l,\sigma }{}_0^t𝑑t^{}e^{i\omega _l(tt^{})}\underset{m}{}h_{l}^{m}{}_{}{}^{}g_l^j\sigma _3^j(t)\zeta _{}^m(t^{})e^{i\stackrel{}{k}_l(\stackrel{}{r}_m\stackrel{}{r}_j)}.$$ Using standard QED methods in the Markovian approximation , this immediately reduces to $$\frac{d\sigma _{}^j}{dt}=i\mathrm{\Delta }_a\sigma _{}^j+\frac{\mu _a}{\mathrm{}}\sigma _3^j\left(\frac{}{2}+f^+\right)iϵ_a\sigma _3^j\overline{\sigma }_{}\frac{1}{2}\gamma _a\sigma _{}^j+$$ $$\underset{l,\sigma }{}_0^t𝑑t^{}e^{i(\omega _l\omega _p)(tt^{})}\underset{m}{}h_{l}^{m}{}_{}{}^{}g_l^j\sigma _3^j(t)\zeta _{}^m(t^{})e^{i\stackrel{}{k}_l(\stackrel{}{r}_j\stackrel{}{r}_m)}$$ (3) in a frame of reference rotating at $`\omega _p`$, such that $`\mathrm{\Delta }_a=\omega _p\omega _a`$. In the context of the analysis of Ref. , with normal ordering, $`f^+`$ is a Langevin force operator arising from fluctuations of the vacuum field, the dephasing rate $`\gamma _a/2`$ is half of the population decay rate $`\gamma _a=4\omega _p^3\|\mu _a\|^2/(3c^3\mathrm{})`$ that is asssociated with the self-field under the condition $`i=j`$, and $`ϵ_a=4\pi N_a\|\mu _a\|^2/(3\mathrm{})`$ is the strength of the near-dipole–dipole (NDD) interaction due to the reaction field of all other atoms, $`ij`$, of species $`a`$, where $`N_a`$ is the relevant number density and $`\overline{\sigma }_{}`$ is a local spatial average of the operator . The last term in Eq. (3) is the contribution of the reaction field from all the atoms of species $`b`$ to the $`j^{th}`$ atom of species $`a`$. A similar analysis yields $$\frac{d\sigma _3^j}{dt}=2[iϵ_a\sigma _+^j\overline{\sigma }_{}\frac{\mu _a}{2\mathrm{}}\sigma _+^j\frac{\mu _a}{\mathrm{}}\sigma _+^jf^+$$ $$\underset{l,\sigma }{}_0^t𝑑t^{}e^{i(\omega _l\omega _p)(tt^{})}\underset{m}{}h_{l}^{m}{}_{}{}^{}g_l^j\sigma _+^j(t)\zeta _{}^m(t^{})e^{i\stackrel{}{k}_l(\stackrel{}{r}_j\stackrel{}{r}_m)}$$ $$+H.c.]\gamma _a(\sigma _3^j+1)$$ (4) for the equation of motion of the inversion operator. In the harmonic oscillator limit for species $`b`$, we have $$\frac{d\zeta _{}^n}{dt}=i\mathrm{\Delta }_b\zeta _{}^n\frac{\mu _b}{\mathrm{}}\left(\frac{}{2}+f^+\right)+iϵ_b\overline{\zeta }_{}\frac{1}{2}\gamma _b\zeta _{}^n+$$ $$\underset{l,\sigma }{}_0^t𝑑t^{}e^{i(\omega _l\omega _p)(tt^{})}\underset{i}{}h_l^ng_{l}^{i}{}_{}{}^{}\sigma _{}^i(t^{})e^{i\stackrel{}{k}_l(\stackrel{}{r}_n\stackrel{}{r}_i)},$$ (5) where $`\mathrm{\Delta }_b=\omega _p\omega _b`$, $`\gamma _b=4\omega _p^3\|\mu _b\|^2/(3c^3\mathrm{})`$, and $`ϵ_b=4\pi N_b\|\mu _b\|^2/(3\mathrm{})`$. Equations (3)–(5) are coupled operator equations for a material composed of two-level systems and harmonic oscillators. None of the relevant parameters for the two-level systems, i.e. frequency, field strength, fluctuations, dephasing rate, population decay rate and NDD, are renormalized by the prescence of the host medium as long as we retain separate equations of motion for the harmonic oscillators. The next step is to adiabatically eliminate the equations of motion for the harmonic oscillators by substituting the formal integral of Eq. (5) into Eqs. (3) and (4). Thus $$\frac{d\sigma _{}^j}{dt}=i\mathrm{\Delta }_a\sigma _{}^j+\frac{\mu _a}{\mathrm{}}\sigma _3^j\left(\frac{}{2}+f^+\right)iϵ_a\sigma _3^j\overline{\sigma }_{}\frac{1}{2}\gamma _a\sigma _{}^j+\underset{l,\sigma }{}_0^t𝑑t^{}e^{i(\omega _l\omega _p)(tt^{})}\underset{m}{}h_{l}^{m}{}_{}{}^{}g_l^je^{i\stackrel{}{k}_l(\stackrel{}{r}_m\stackrel{}{r}_j)}$$ $$_0^t^{}𝑑t^{\prime \prime }e^{\alpha (t^{}t^{\prime \prime })}\left[\frac{\mu _b}{\mathrm{}}\sigma _3^j(t)\left(\frac{(t^{\prime \prime })}{2}+f^+(t^{\prime \prime })\right)+\underset{l^{},\sigma ^{}}{}_0^{t^{\prime \prime }}𝑑t^{\prime \prime \prime }e^{i(\omega _l^{}\omega _p)(t^{\prime \prime }t^{\prime \prime \prime })}\underset{i}{}h_l^{}^mg_{l^{}}^{i}{}_{}{}^{}\sigma _3^j(t)\sigma _{}^i(t^{\prime \prime \prime })e^{i\stackrel{}{k}_l^{}(\stackrel{}{r}_m\stackrel{}{r}_i)}\right],$$ (6) where $`\zeta _+^m(0)=0`$ and $`\alpha =i(\mathrm{\Delta }_b+ϵ_b+i\gamma _b/2)`$. The last term is the part of the reaction field that is due to the prescence of the harmonic oscillators and their subsequent adiabatic elimination. The large square bracket contains terms that are largely equivalent to all of the original field components, the coherent field, vacuum fluctuations, the self-field, and the reaction field, and will lead to the renormalization of each. The self-field contribution, $`i=j`$, in which the atom couples to itself via the linear particles, can be evaluated using the transverse delta function . The reaction field contribution, the near dipole–dipole interaction of all the $`i`$ atoms with atom $`m`$, is of the same form as the NDD interaction of the $`i`$ atoms with atom $`j`$ that was studied in Ref. and used in obtaining Eq. (3). We refer to this type of interaction as a many-atom Milonni-Knight problem . The sum in the large square brackets becomes $$i\frac{4\pi }{3\mathrm{}}N_a\mu _a^{}\mu _b\sigma _3^j(t)\overline{\sigma }_{}(t^{\prime \prime })+\frac{1}{2}\frac{4\omega _p^3}{3c^3\mathrm{}}\mu _a^{}\mu _b\sigma _3^j(t)\sigma _{}^j(t^{\prime \prime }).$$ (7) In the last term of Eq. (6), note that $$_0^t^{}𝑑t^{\prime \prime }e^{\alpha (t^{}t^{\prime \prime })}f(t^{\prime \prime })\frac{1}{\alpha }f(t^{}),$$ (8) since the exponential is strongly peaked near $`t^{}=t^{\prime \prime }`$. Because $`\stackrel{}{r}_m\stackrel{}{r}_j`$, the remaining part of the last term in Eq. (6) is again the many-atom Milonni-Knight problem. Then the portion exterior to the large square brackets can be written as $$\frac{4\pi }{3\mathrm{}}\frac{N_b\|\mu _b\|^2}{\mathrm{\Delta }_b+ϵ_b+i\gamma _b/2}\frac{\mu _a}{\mu _b^{}}\left[\mathrm{}\right]=(\mathrm{}1)\frac{\mu _a}{\mu _b^{}}\left[\mathrm{}\right],$$ (9) where $`N_b`$ is the number density of oscillators, species $`b`$, $$\mathrm{}=\frac{n_b^2+2}{3}$$ (10) is the complex Lorentz dielectric local-field enhancement factor, and $`n_b`$ is the complex index of refraction of the host material. Using (7) and (9) in Eq. (6), we obtain $$\frac{d\sigma _{}^j}{dt}=i\mathrm{\Delta }_a\sigma _{}^j\frac{\mu _a}{\mathrm{}}\sigma _3^j\left(\frac{\mathrm{}}{2}+\mathrm{}f^+\right)+i\mathrm{}ϵ_a\sigma _3^j\overline{\sigma }_{}$$ $$\frac{1}{2}\mathrm{}\gamma _a\sigma _{}^j.$$ (11) A similar calculation for the equation of motion of the inversion operator, Eq. (4), yields $$\frac{d\sigma _3^j}{dt}=2[i\mathrm{}ϵ_a\sigma _+^j\overline{\sigma }_{}\frac{\mu _a}{\mathrm{}}\sigma _+^j(\frac{\mathrm{}}{2}+\mathrm{}f^+)$$ $$\frac{1}{4}\mathrm{}\gamma _a(\sigma _3^j+1)+H.c.]$$ (12) Equations (11) and (12) can be considered as operator Langevin–Bloch equations of motion for a dense collection of two-level atoms embedded in a dielectric medium. The effect of the adiabatically eliminated damped linear oscillators is contained in the complex Lorentz dielectric enhancement factor $`\mathrm{}`$ that renormalizes the coherent field, the Langevin force operator, the NDD interaction, the dephasing rate and the population decay rate. We have neglected Cauchy principle parts throughout and there will be local-field enhancement effects from these, as well, e.g. renormalization of the Lamb shift by $`Re(\mathrm{})`$. The results of our fully microscopic many-body QED treatment agree with semiclassical results for local-field enhancement of the coherent field and the NDD interaction , if expectation values are taken in the limit of classical factorization. The purely quantum mechanical aspects are the enhancement of the Langevin force operator and the corresponding damping rates. At resonance, the population decay rate that appears in the Langevin–Bloch equation of motion for the inversion operator is the renormalized spontaneous emission rate $$\mathrm{\Gamma }_{SE}^{diel}=Re(\mathrm{}(\omega _a))\gamma _a(\omega _a)=Re(\mathrm{})\mathrm{\Gamma }_0.$$ (13) The spontaneous emission rate is renormalized by $`Re(\mathrm{})`$. In addition, there is a level shift, or frequency renormalization, in the amount of $`Im(\mathrm{})\gamma _a/2`$ due to the product of the imaginary part of $`\mathrm{}`$ with the dephasing rate. There have been a number of experimental measurements of the spontaneous emission rate of embedded atoms, or atom-like particles, in a dielectric . However, these experiments typically involve non-trivial boundary conditions, such as organic ligand cages or nanospheres, that can profoundly affect the observed spontaneous emission rate. To date, we know of no experimental measurements of the index dependence of the spontaneous emission rate in a bulk dielectric. The complete theory presented here makes it possible to bring the entire arsenal of laser spectroscopy to bear on the measurement of local-field effects due to a dielectric background. For example, because even small frequency shifts can be resolved spectroscopically, it might be possible to verify our results by measuring the level shift $`Im(\mathrm{})\gamma _a/2`$ or the renormalization of the Lamb shift by $`Re(\mathrm{})`$ as a function of the density of a buffer gas. Operator equations of motion, including the local-field renormalization of the spontaneous emission rate, for a single atom, or a tenuous collection of atoms, embedded in a dielectric are contained in these, more general, results in the limit $`ϵ_a0`$. The NDD interaction of densely embedded atoms in a dielectric is included primarily because it provides a very useful theoretical backdrop for evaluating the many-body effects. The same type of many-body summations that occur in the local-field renormalization of the spontaneous emission rate have been studied previously in the evaluation of local-field effects for dense atoms in vacuo . Further, it was our semiclassical derivation of the ‘anomalous’ renormalization of NDD interaction in a dielectric host that indicated a need to examine the problem in its entirety. Because the NDD interaction and the spontaneous emission rate have the same dependence on the dipole moment and arise in the same way from the elimination of the field operator, one would expect them to have the same renormalization in a dielectric. We have shown that this is the case. In conclusion, we obtained the renormalization of the spontaneous emission rate of an atom embedded in a dielectric host. This result was obtained from a fully microscopic, many-body derivation of Langevin–Bloch operator equations for two-level atoms embedded in an absorptive and dispersive, linear dielectric host. The dielectric local-field enhancement of the coherent field, the Langevin force operator, the NDD interaction, and the damping rates all stem from the same reaction field that arises from the nearby harmonic oscillators, necessitating the full many-body derivation. We found that the dielectric enhancement of the spontaneous emission rate is much smaller than indicated by previous studies. This is an enabling result, paving the way for application of high-index materials to enhance quantum optical effects.
warning/0004/math0004024.html	ar5iv	text	# A small transitive family of continuous functions on the Cantor set ## 1. Introduction In \[Dow93\] Dow gave a proof of the Rudin-Shelah theorem about the existence of $`2^𝔠`$ points in $`\beta `$ that are Rudin-Keisler incomparable. The proof actually shows that whenever a family $``$ of $`𝔠`$ continuous self-maps of $`\beta `$ (or $`^{}`$) are given there is a set $`S`$ of $`2^𝔠`$ many $``$-independent points in $`\beta `$ (or $`^{}`$). This suggests that we measure the complexity of a space $`X`$ by the cardinal number $`𝔱𝔣(X)`$ defined as the minimum cardinality of a set $``$ of continuous self maps such that for all $`x,yX`$ there is $`f`$ such that $`f(x)=y`$ or $`f(y)=x`$. Let us call such an $``$ transitive. Thus Dow’s proof shows $`𝔱𝔣(\beta )`$, $`𝔱𝔣(^{})𝔠^+`$. We investigate $`𝔱𝔣(C)`$, where $`C`$ denotes the Cantor set. Van Mill observed that $`𝔱𝔣(C)\mathrm{}_1`$; a slight extension of his argument shows that $`\mathrm{MA}(\mathrm{countable})`$ implies $`𝔱𝔣(C)=𝔠`$. Our main result states that in the Sacks model the continuous functions on the Cantor set that are coded in the ground model form a transitive set. Thus we get the consistency of $`𝔱𝔣(C)=\mathrm{}_1<\mathrm{}_2=𝔠`$. The gap between $`𝔱𝔣(C)`$ and $`𝔠`$ cannot be arbitrarily wide, because Hajnal’s free set lemma implies that for any space $`X`$ one has $`\|X\|𝔱𝔣(X)^+`$. In \[Mil84\] Miller showed that it is consistent with $`ZFC`$ that for every set of reals of size continuum there is a continuous map from that set onto the closed unit interval. In fact he showed that the iterated perfect set model of Baumgartner and Laver (see \[BauLav79\]) is such a model, and noted that the continuous map can even be coded in the ground model. Here we will show that in the iterated perfect set model, for every two reals $`x`$ and $`y`$ there exists a continuous function with code in the ground model that maps $`x`$ onto $`y`$ or $`y`$ onto $`x`$. ###### Definition 1. By a transitive set of functions $``$ we mean a set of continuous functions such that for every two reals $`x`$ and $`y`$ there exists an element $`f`$ such that $`f(x)=y`$ or $`f(y)=x`$ holds. Let us also define the cardinal number $`𝔱𝔣`$ by $$𝔱𝔣=\mathrm{min}\{\|\|:\text{is a transitive set of functions}\},\text{i.e.}𝔱𝔣=𝔱𝔣(C).$$ The paper is organized as follows: in section 2 we prove some simple facts on $`𝔱𝔣`$, the minimal size of transitive sets of functions. We also state and prove the main theorem of this paper in section 2, using theorems proved later on in section 3. As a corollary to the main theorem we have the consistency of $`𝔱𝔣<𝔠`$ with $`ZFC`$. Finally in section 4 we will make a remark on the effect on $`𝔱𝔣`$ when we add $`\kappa `$ many Sacks reals side-by-side to a model of $`ZFC+CH`$. ## 2. Notation and Preliminaries For the rest of this paper let $`V`$ be a model of $`ZFC`$. We will use the same notations and definitions as Baumgartner and Laver in \[BauLav79\], so for any ordinal $`\alpha `$ we let $`_\alpha `$ denote the poset that iteratively adds $`\alpha `$ Sacks reals to the model $`V`$, using countable support. Let $`_1=`$, where $``$ denotes the ’normal’ Sacks poset for the addition of one Sacks real. Let $`G_\alpha `$ be $`_\alpha `$-generic over $`V`$, we define $`V_\alpha `$ by $`V_\alpha =V[G_\alpha ]`$ for every ordinal $`\alpha `$. Note that if $`\beta <\alpha `$ we have that $`G_\alpha \beta `$ is a $`_\beta `$-generic subset over $`V`$. If we denote the $`(\alpha +1)`$-th added Sacks real by $`s_\alpha `$ then we can also write $`V_\alpha =V[s_\beta :\beta <\alpha ]`$. Assuming $`VCH`$, the proof of the following facts can be found in \[BauLav79\]: 1. Forcing with $`_\alpha `$ does not collapse cardinals. 2. $`V_{\omega _2}`$ is a model of $`ZFC+2^\mathrm{}_0=\mathrm{}_2`$. 3. Let $`\dot{}_\beta `$ denote the result of defining $`_\beta `$ in $`V_\alpha `$. Then for any $`\alpha ,\beta 1`$, $`_\alpha \mathrm{`}\mathrm{`}_{\alpha ,\alpha +\beta }\text{is isomorphic to}\dot{}_\beta \mathrm{"}`$. We will now prove some facts on the cardinal $`𝔱𝔣`$. The first is Van Mill’s observation alluded to above. ###### Theorem 2. $`𝔱𝔣\mathrm{}_1`$. Proof Suppose $``$ is a countable set of functions. Let $`A_f`$ denote the set $`\{x:\mathrm{int}(f^1(x))\mathrm{}\}`$ for every $`f`$. Every $`A_f`$ is at most countable because $`2^\omega `$ is separable. So choose an $`x`$ in $`2^\omega _fA_f`$, then we know that for every $`f`$ the set $`f^1(x)`$ is nowhere dense in $`2^\omega `$. For such an $`x`$ the set $`\{f^1(x):f\}`$ is countable. Because the set $`\{f(x):f\}`$ is also countable the Baire category theorem tells us that the set $`2^\omega _f(\{f(x)\}f^1(x))`$ is nonempty, thus showing that $``$ is not transitive. ###### Theorem 3. $`𝔱𝔣𝔠𝔱𝔣^+`$. ###### Remark 4. The proof of theorem 2 shows that $`𝔱𝔣`$ is at least the minimum number of nowhere dense sets needed to cover $`C`$. Then theorem 3 and $`\mathrm{MA}(\mathrm{countable})`$ imply $`𝔱𝔣=𝔠`$. The second inequality is a consequence of the following lemma. The proof of this lemma can be found in \[Wil79\]. For this we need some more notation. Let $`S`$ be an arbitrary set. By a set mapping on $`S`$ we mean a function $`f`$ mapping $`S`$ into the power set of $`S`$. The set map is said to be of order $`\lambda `$ if $`\lambda `$ is the least cardinal such that $`\|f(x)\|<\lambda `$ for each $`x`$ in $`S`$. A subset $`S^{}`$ of $`S`$ is said to be free for $`f`$ if for every $`xS^{}`$ we have $`f(x)S^{}\{x\}`$. ###### Lemma 5 (Free set lemma). Let $`S`$ be a set with $`\|S\|=\kappa `$ and $`f`$ a set map on $`S`$ of order $`\lambda `$ where $`\lambda <\kappa `$. Then there is a free set of size $`\kappa `$ for $`f`$. Proof\[Theorem 3\] The proof of the first inequality is easy. We simply have to observe that the set of all constant functions on the reals is a transitive set of functions. Now for the second inequality. Striving for a contradiction suppose that $`𝔠𝔱𝔣^{++}`$. Let $``$ be a transitive set of functions such that $`\|\|=𝔱𝔣`$. We define a set map $`F`$ on the reals by $`F(x)=\{f(x):f\}`$ for every $`x2^\omega `$. Because $`\|F(x)\|𝔱𝔣`$, this set map $`F`$ is of order $`𝔱𝔣^+`$, which is less than $`𝔠`$. According to the free set lemma there exists a set $`X2^\omega `$ such that $`\|X\|=𝔠`$ and for every $`xX`$ we have $`F(x)X\{x\}`$. This is a contradiction, because every two reals in $`X`$ provide a counter example of $``$ being a transitive set. Closed subsets of the Cantor set can be coded by sub trees of $`{}_{}{}^{<\omega }2`$, as follows: if $`A`$ is closed then let $`T_A=\{xn:xA,n\omega \}`$; one can recover $`A`$ from $`T_A`$ by observing that $`A=\{x{}_{}{}^{\omega }2:n\omega ,xnT_A\}`$. When we say that a closed set $`A`$ is coded in the ground model we mean that $`T_A`$ belongs to the ground model. We shall always construct a continuous function $`f`$ between closed sets $`A`$ and $`B`$ by specifying an order-preserving map $`\varphi `$ from $`T_A^{}`$ to $`T_B`$, where $`T_A^{}`$ denotes the set of splitting nodes of $`T_A`$. Once $`\varphi `$ is found one defines $`f`$ by $$f(x)=\mathrm{`}\mathrm{`}\text{the path through}T_B\text{determined by the restriction of}\varphi \text{to}\{xn:n\omega \}\mathrm{"}.$$ We say that $`f`$ is coded in the ground model if $`\varphi `$ belongs to $`V`$. In what follows we shall denote the map $`\varphi `$ by $`f`$ as well. Let us define the set $`𝒢`$ (in any $`V_\alpha `$) by $$𝒢=\{f:f\text{is a continuous function with code in}V\}$$ Now we can explicitly state the main theorem of this paper. Section 3 is completely devoted to the proof of this theorem by parts, so we will prove the theorem here and refer to the needed theorems proved in that section. ###### Theorem 6 (Main Theorem). The set $`𝒢`$ is transitive in $`V_\alpha `$ for every ordinal $`\alpha `$. Proof We will show by transfinite induction that $`𝒢`$ is a transitive set in $`V_\alpha `$ for all $`\alpha `$. For $`\alpha =0`$ this is obvious. Suppose the theorem is true for all $`\beta <\alpha `$. Let $`x`$ and $`y`$ be reals in $`V_\alpha `$. If $`\alpha `$ is a successor ordinal, $`\alpha =\beta +1`$, then we use theorem 11 in the case that at least one of $`x`$ and $`y`$ is not in $`V_\beta `$ to show that there exist a continuous function $`f`$ defined in $`V`$ (so $`f𝒢`$) such that in $`V_\alpha `$ we have $`f(x)=y`$ or $`f(y)=x`$. Since we are forcing with countable support and because reals are countable objects, there are no new reals added by $`_\alpha `$ for $`\mathrm{cf}(\alpha )>\mathrm{}_0`$. So if $`\alpha `$ is a limit ordinal we only have to consider the case where $`\mathrm{cf}(\alpha )=\mathrm{}_0`$ and at least one of $`x,y`$ is not in $`_{\beta <\alpha }V_\beta `$. Then we use theorem 17 to show the existence of an continuous function $`f`$ defined in $`V`$ such that in $`V_\alpha `$ $`f(x)=y`$ or $`f(y)=x`$ holds. As is well-known, if $`VCH`$ then $`V_{\omega _2}𝔠=\mathrm{}_2`$. This enables us to show that $`𝔱𝔣<𝔠`$ is consistent. ###### Corollary 7. If $`VCH`$ then $`V_{\omega _2}𝔱𝔣<𝔠`$. In this paper we shall repeatedly use the fact that any homeomorphism $`h`$ between two closed nowhere dense subsets of the Cantor set can be extended to a homeomorphism of the Cantor set onto itself (see \[KnaRei53\]). Furthermore it is straight forward to extend a continuous function between to closed nowhere dense (disjoint) subsets of the Cantor set to a continuous self map of the Cantor set. Because we can make sure that the subsets of the Cantor set that define the added reals $`x`$ and $`y`$ are nowhere dense and closed, when we show that there exists a homeomorphism (or a continuous function) $`f`$ mapping of one of these sets onto the other, in such a way that in the extension $`x`$ is mapped onto $`y`$ or vice versa, we actually have shown that there exists a self map of the Cantor that is a homeomorphism (continuous function) mapping, in the extension, $`x`$ onto $`y`$ or $`y`$ onto $`x`$. ## 3. The continuous functions with code in the ground model $`V`$ form a transitive set in $`V_\alpha `$ In this section we prove that for every $`\alpha `$ and any new real $`x`$ in the Baumgartner and Laver model $`V_\alpha `$ (i.e. $`xV_\alpha _{\beta <\alpha }V_\beta `$) and $`y`$ any real in $`V_\alpha `$ there exists a function $`f`$ defined in the ground model $`V`$ such that in $`V_\alpha `$ the equation $`f(x)=y`$ holds. We make the following definition. For any $`\sigma {}_{}{}^{<\omega }2`$ we let $`l(\sigma )\omega `$ denote the length of $`\sigma `$. So for every $`\sigma {}_{}{}^{<\omega }2`$ we have $`\sigma {}_{}{}^{l(\sigma )}2`$. To show how we construct our continuous maps we reprove the familiar fact that Sacks reals are minimal, see \[Jec78\]. ###### Lemma 8. Suppose $`x`$ is a real in $`V[G]V`$, where $`G`$ is a $``$-generic filter over $`V`$, and that $`p`$ is such that $`p\mathrm{`}\mathrm{`}\dot{x}V\mathrm{"}`$. Then there exists a $`qp`$ and a homeomorphism $`f`$ defined in $`V`$ such that $`q\mathrm{`}\mathrm{`}f(\dot{s})=\dot{x}\mathrm{"}`$. Here $`\dot{s}`$ denotes the name of the added Sacks real. Proof We will construct a fusion sequence $`\{p_i,n_i:i\omega \}`$ such that each $`p_{i+1}`$ will know all the first $`i`$ splitting nodes of every branch of the perfect tree $`p_i`$ and $`(p_{i+1},n_{i+1})>(p_i,n_i)`$ for every $`i`$. Because $`p`$ forces that $`\dot{x}`$ is a new real, there exists an element $`u_{\mathrm{}}{}_{}{}^{<\omega }2`$ with maximal length $`m_{\mathrm{}}`$, such that $`p\mathrm{`}\mathrm{`}\dot{x}m_{\mathrm{}}=u_{\mathrm{}}\mathrm{"}`$ and $`p`$ does not decide $`\dot{x}(m_{\mathrm{}})`$. There exist $`p_0,p_1p_0`$ such that $`p_k\mathrm{`}\mathrm{`}\dot{x}(m_{\mathrm{}})=k\mathrm{"}`$ for $`k\{0,1\}`$. Without loss of generality the stems of $`p_0`$ and $`p_1`$ are incompatible. Let $`n_0=\mathrm{min}\{n\omega :p_0np_1n\}`$ and let $`p_0`$ denote the element $`p_0p_1`$. Now assume we have $`p_i=\{p_\sigma :\sigma {}_{}{}^{i+1}2\}`$. Consider $`\tau {}_{}{}^{i+1}2`$, we have an element $`u_\tau {}_{}{}^{<\omega }2`$ of maximal length $`m_\tau `$ such that $`p_\tau \mathrm{`}\mathrm{`}\dot{x}m_\tau =u_\tau \mathrm{"}`$. There exist $`p_{\tau {}_{}{}^{}0}`$, $`p_{\tau {}_{}{}^{}1}p_\tau `$ such that $`p_{\tau {}_{}{}^{}k}\mathrm{`}\mathrm{`}\dot{x}(m_\tau )=k\mathrm{"}`$ for $`k\{0,1\}`$. Again without loss of generality the stems of $`p_{\tau {}_{}{}^{}0}`$ and $`p_{\tau {}_{}{}^{}1}`$ are incompatible. Let $`n_\tau `$ denote the integer $`\mathrm{min}\{n\omega :p_{\tau {}_{}{}^{}0}np_{\tau {}_{}{}^{}0}n\}`$ and $`n_{i+1}=\mathrm{max}\{n_\sigma :\sigma {}_{}{}^{i+1}2\}`$. We let $`p_{i+1}`$ denote the element $`\{p_\sigma :\sigma {}_{}{}^{i+2}2\}`$. Now the induction step is completed, because $`p_{i+1}`$ knows all the first $`i+1`$ splitting nodes of every branch in $`p_i`$ and $`(p_{i+1},n_{i+1})>(p_i,n_i)`$ for every $`i\omega `$. We define the function $`f`$ by $$f^1([u_\sigma ])[\mathrm{stem}(p_\sigma )]\text{for}\sigma {}_{}{}^{<\omega }2.$$ As $`\mathrm{stem}(p_\sigma )`$ is a finite approximation of the added Sacks real $`\dot{s}`$, we have by the construction of our $`p_\sigma `$ for $`\sigma {}_{}{}^{<\omega }2`$ and the function $`f`$ that $`p_\sigma \mathrm{`}\mathrm{`}f(\dot{s})[u_\sigma ]\mathrm{"}`$ for every $`\sigma {}_{}{}^{<\omega }2`$. And so the fusion $`q`$ of the sequence $`\{p_i,n_i:i\omega \}`$ forces that in the extension $`V[G]`$ the equality $`f(s)=x`$ holds. This $`f`$, being a continuous bijection between two Cantor sets, is (of course) a homeomorphism. ###### Remark 9. In the lemma we have also defined a map $`\varphi `$ from the finite sub trees of the fusion $`q`$ onto the finite sub trees of $`T=_{\sigma {}_{}{}^{<\omega }2}u_\sigma `$ which induces our homeomorphism. We have $`\varphi (q)=T`$ and $$\varphi ([q\sigma ])=\{u_\tau :\sigma \tau \text{and}\tau {}_{}{}^{<\omega }2\}.$$ We note that $`[T]`$ is the set of all the possible interpretations of $`\dot{x}`$ in $`V[G]`$ and that $`T`$ depends on $`\varphi `$ and $`q`$ only. In theorem 11 we will use this interpretation of the previous lemma. As a warming up exercise we prove the following. ###### Theorem 10. The set $`𝒢`$ is transitive in $`V_1`$. Proof Suppose $`x`$ and $`y`$ are two reals of $`V_1(=V[s_0])`$. We consider two cases. Case 1: $`x`$ is a real in $`V`$. The constant function $`c_x=\{y,x:y\text{a real in}V_1\}`$ is a continuous function defined in $`V`$, thus a member of $`𝒢`$, and in $`V_1`$ it maps $`y`$ onto $`x`$. Case 2: both $`x`$ and $`y`$ are reals not in $`V`$. Let $`p`$ be a witness of this, so $`p\mathrm{`}\mathrm{`}\dot{x},\dot{y}V\mathrm{"}`$. According to lemma 8 there exists a $`qp`$ and a homeomorphism $`f`$ defined in $`V`$ such that $`q\mathrm{`}\mathrm{`}f(\dot{s}_0)=\dot{x}\mathrm{"}`$, where $`\dot{s}_0`$ denotes the added Sacks real. If we apply the lemma again we get an $`rq`$ and a homeomorphism $`g`$ defined in $`V`$ such that $`q\mathrm{`}\mathrm{`}g(\dot{s}_0)=\dot{y}\mathrm{"}`$. But now we have that $`r\mathrm{`}\mathrm{`}(gf^1)(\dot{x})=\dot{y}\mathrm{"}`$ and we see that $`gf^1`$ is the element of $`𝒢`$ we are looking for. ###### Theorem 11. For $`\alpha `$ an ordinal and $`x`$ and $`y`$ reals in $`V_{\alpha +1}`$ such that $`xV_\alpha `$ there exists an $`f𝒢`$ such that in $`V_{\alpha +1}`$ $`f(x)=y`$ holds. Moreover if also $`yV_\alpha `$ then $`f`$ can be chosen to be a homeomorphism. Proof This is an immediate consequence of the lemmas 14 and 15. We make the following definitions. For $`p`$ and $`s{}_{}{}^{<\omega }2`$ we let $`p_s`$ denote the sub-tree $`\{tp:st\text{or}ts\}`$ of $`p`$. Of course $`p_s`$ is a perfect tree if and only if $`s{}_{}{}^{<\omega }2p`$. To generalize this to $`_\alpha `$, suppose $`p`$ is an element of $`_\alpha `$, $`F`$ is a finite subset of $`\mathrm{dom}(p)`$ and $`n\omega `$, we say that a function $`\tau :F{}_{}{}^{n}2`$ is consistent with $`p`$ if the following holds for every $`\beta F`$: $$(p\tau )\beta _\beta \mathrm{`}\mathrm{`}\tau (\beta )p(\beta )\mathrm{"}.$$ So we have for every $`\beta F`$ that $`(p\tau )\beta _\beta \mathrm{`}\mathrm{`}(p(\beta ))_{\tau (\beta )}\text{is a perfect tree}\mathrm{"}`$. Furthermore let us suppose that $`F`$ and $`H`$ are two sets such that $`FH`$, and $`n`$ and $`m`$ are two integers such that $`m<n`$, if $`\tau `$ is a function mapping $`F`$ into $`{}_{}{}^{m}2`$ then we say that a function $`\sigma :H{}_{}{}^{n}2`$ extends the function $`\tau `$ if for every $`iF`$ we have $`\sigma (i)m=\tau (i)`$. For later use we will prove the following: ###### Lemma 12. Let $`p_\alpha `$, $`F[\mathrm{dom}(p)]^{<\omega }`$ and $`n\omega `$. Suppose $`\tau :F{}_{}{}^{n}2`$ is consistent with $`p`$ then for every $`rp\tau `$ there exists a $`qp`$ such that $`q\tau =r`$ and $`q\beta _\beta \mathrm{`}\mathrm{`}(p(\beta ))_s=(q(\beta ))_s\text{for every}s{}_{}{}^{n}2`$ such that $`s\tau (\beta )\mathrm{"}`$ for every $`\beta F`$. Proof Define the element $`q_\alpha `$ as follows for $`\beta <\alpha `$: $$q\beta _\beta \mathrm{`}\mathrm{`}q(\beta )=\{\begin{array}{cc}r(\beta )\hfill & \beta F\hfill \\ r(\beta )\{(p(\beta ))_s:s{}_{}{}^{n}2p(\beta )\text{such that}s\tau (\beta )\hfill & \beta F\mathrm{"}.\hfill \end{array}$$ In this way we strengthen the tree $`p(\beta )`$ above $`\tau (\beta )`$ keeping the rest of the perfect tree intact (according to $`F`$ anyway). We need the following lemma to make sure that the maps we will construct in the lemmas 14 and 15 are well-defined and continuous. ###### Lemma 13. Let $`p_{\alpha +1}`$. Suppose $`F,H[\mathrm{dom}(p)]^{<\omega }`$ are such that $`FH`$ and $`m,n\omega `$ are such that $`m<n`$. If $`\tau :F{}_{}{}^{m}2`$ is consistent with $`p`$, $`N`$ is an integer and $`T`$ is a finite tree such that $$(p\tau )\alpha \mathrm{`}\mathrm{`}p(\alpha ){}_{}{}^{N}2=T\mathrm{"},$$ then there exist a $`(q,j)>_H(p\tau ,n)`$ and an $`M>N`$ such that for every $`\sigma :H{}_{}{}^{n}2`$ extending $`\tau `$, if $`\sigma `$ is consistent with $`q`$, then there exists $`T_\sigma `$ such that $`q\sigma \mathrm{`}\mathrm{`}q(\alpha ){}_{}{}^{M}2=T_\sigma \mathrm{"}`$. Also $`\|(T_\sigma )_t{}_{}{}^{M}2\|2`$ for every $`tT`$ and $`[T_\sigma ][T_\varsigma ]=\mathrm{}`$ whenever $`\sigma `$ and $`\varsigma `$ are distinct and consistent with $`q`$. Proof Let $`\mathrm{\Sigma }_\tau `$ denote the set of all $`\sigma :H{}_{}{}^{n}2`$ extending $`\tau `$. Because $`p(\alpha )`$ is a perfect tree there exists a $`_\alpha `$-name $`\dot{M}`$ such that for every $`tT`$ we have $$(p\tau )\alpha \mathrm{`}\mathrm{`}\|(p(\alpha ))_t{}_{}{}^{\dot{M}}2\|2\|\mathrm{\Sigma }_\tau \|\mathrm{"}.$$ According to lemma 2.3 of \[BauLav79\] there exists a $`(q^{},j^{})>_H((p\tau )\alpha ,n)`$ such that if $`\sigma \mathrm{\Sigma }_\tau `$ is consistent with $`q^{}`$ we have an $`M_\sigma `$ such that $`q^{}\sigma \mathrm{`}\mathrm{`}\dot{M}=M_\sigma \mathrm{"}`$. Put $`M=\mathrm{max}\{M_\sigma :\sigma \mathrm{\Sigma }_\tau \text{consistent with}q^{}\}`$. We have $`q^{}\mathrm{`}\mathrm{`}\|(p(\alpha ))_t{}_{}{}^{M}2\|2\|\mathrm{\Sigma }_\tau \|\mathrm{"}`$ for every $`tT`$. Enumerate $`\{\sigma \mathrm{\Sigma }_\tau :\sigma \text{consistent with}q^{}\}`$ as $`\{\sigma _k:k<K\}`$. Let $`rq^{}\sigma _0`$ be such that $`r\mathrm{`}\mathrm{`}p(\alpha ){}_{}{}^{M}2=S_{\sigma _0}\mathrm{"}`$, where $`S_{\sigma _0}`$ is such that $`\|(S_{\sigma _0})_t{}_{}{}^{M}2\|2\|\mathrm{\Sigma }_\tau \|`$ for every $`tT`$. Use lemma 12 to find a $`q_0q^{}`$ such that $`q_0\sigma _0=r`$. We continue this procedure with all the $`\sigma _k\mathrm{\Sigma }_\tau `$. So if $`\sigma _k`$ is consistent with $`q_{k1}`$ we find an $`rq_{k1}\sigma _k`$ such that $`r\mathrm{`}\mathrm{`}p(\alpha ){}_{}{}^{M}2=S_{\sigma _k}\mathrm{"}`$, and also that $`\|(S_{\sigma _k})_t{}_{}{}^{M}2\|2\|\mathrm{\Sigma }_\tau \|`$ for every $`tT`$. And we use lemma 12 to define $`q_kq_{k1}`$ such that $`q_k\sigma _k=r`$. If $`\sigma _k`$ is not consistent with $`q_{k1}`$ we choose $`q_k=q_{k1}`$. We now have for every $`\sigma \mathrm{\Sigma }_\tau `$ consistent with $`q_{K1}`$ a finite tree $`S_\sigma {}_{}{}^{M}2`$ extending the tree $`T`$ such that every branch in $`T`$ has (at least) $`2\|\mathrm{\Sigma }_\tau \|`$ different extensions in $`S_\sigma {}_{}{}^{M}2`$ and $`q_{K1}\sigma \mathrm{`}\mathrm{`}p(\alpha ){}_{}{}^{M}2=S_\sigma \mathrm{"}`$. As $`q_{K1}`$ forces that, for each $`yT`$ the size of the set $`p(\alpha ))_t^M2`$ is at least $`2\|\mathrm{\Sigma }_\tau \|`$ we can find for $`\sigma \mathrm{\Sigma }_\tau `$ consistent with $`q_{K1}`$ a sub tree $`T_\sigma `$ of $`S_\sigma `$ such that $`\|(T_\sigma )_t{}_{}{}^{M}2\|2`$ and whenever $`\sigma `$ and $`\varsigma `$ are distinct and consistent with $`q_{K1}`$ we have $`[T_\sigma ][T_\varsigma ]=\mathrm{}`$. Define $`q_{\alpha +1}`$ such that $`q\alpha =q_{K1}`$ and choose $`q(\alpha )`$ such that for every consistent $`\sigma \mathrm{\Sigma }_\tau `$ we have $`q\sigma \mathrm{`}\mathrm{`}q(\alpha )=p(\alpha )[T_\sigma ]\mathrm{"}`$. If we let $`j`$ be equal to $`\mathrm{max}\{j^{},M\}`$ the proof is complete. ###### Lemma 14. Given an ordinal $`\alpha `$, a $`p_{\alpha +1}`$ and $`_{\alpha +1}`$-names $`\dot{x}`$ and $`\dot{y}`$ such that $`p\mathrm{`}\mathrm{`}\dot{x}V_\alpha \text{and}\dot{y}V_\alpha \mathrm{"}`$ then there exists a continuous function $`f`$ defined in $`V`$ and a $`qp`$ such that $`q\mathrm{`}\mathrm{`}f(\dot{x})=\dot{y}\mathrm{"}`$. Proof By remark 9 we know that there is an $`rp\alpha `$ and there exist $`_{\alpha +1}`$ names $`\dot{\varphi }`$ for a map on the finite sub trees of $`p(\alpha )`$ and $`\dot{T}`$ for a perfect tree such that $`r\mathrm{`}\mathrm{`}\dot{\varphi }(p(\alpha ))=\dot{T}\mathrm{"}`$. Without loss of generality we assume that $`p\alpha =r`$. Let us construct a fusion sequence $`\{p_i,n_i,F_i:i\omega \}`$. Let $`p_0=p_1=p`$, $`n_0=n_1=0`$, $`F_0=\mathrm{}`$ and choose $`F_1[\mathrm{dom}(p)]^{<\omega }`$ in such a way that we are building a fusion sequence. Suppose we have constructed the sequence up to $`i`$, let us construct the next element of the fusion sequence. We let $`\{\tau _k:k<K\}`$ denote all $`\tau :F_{i1}{}_{}{}^{n_{i1}}2`$ consistent with $`p_i`$. If we choose in lemma 13 $`\tau =\tau _0`$, $`F=F_{i1}`$ and $`m=n_{i1}`$ we get a $`(q_0,m_0)>_{F_i}(p_i\tau _0,n_i)`$ such that for every $`\sigma :F_i{}_{}{}^{n_i}2`$ extending $`\tau _0`$, consistent with $`q_0`$, we have a finite sub tree $`T_\sigma {}_{}{}^{M(\tau _0)}2`$ ($`M(\tau _0)\omega `$ follows from lemma 13) of $`p_i(\alpha )=p(\alpha )`$ such that 1. $`T_\sigma `$ is an extension of $`T_{\tau _0}`$, 2. for every branch $`t`$ in $`T_{\tau _0}`$ there exist at least two different branches of length $`M(\tau _0)`$ in $`T_\sigma `$ extending $`t`$, 3. if $`\sigma `$ and $`\varsigma `$ are two distinct members of $`\mathrm{\Sigma }_{\tau _0}`$ consistent with $`q_0`$ we have $`[T_\sigma ][T_\varsigma ]=\mathrm{}`$. We choose $`r_0_{\alpha +1}`$ with lemma 12 such that $`r_0q_0`$ and $`r_0\tau _0=q_0`$. We iteratively consider all the $`\tau :F_{i1}{}_{}{}^{n_{i1}}2`$. In the general case if $`\tau _k`$ is consistent with $`r_{k1}`$ then lemma 13 gives us a $`q_k`$ and an $`m_k\omega `$ such that $`(q_k,m_k)>_{F_i}(r_{k1}\tau _k,n_i)`$. We choose $`r_k`$ in the same way as above, using lemma 12 such that $`r_kq_k`$ and $`r_k\tau _k=q_k`$. If $`\tau _k`$ is inconsistent with $`r_{k1}`$ then we choose $`r_k=r_{k1}`$ and $`m_k=m_{k1}`$. After considering all the $`\tau _k`$’s we define $`p_{i+1}=r_{K1}`$ and $`n_{i+1}=\mathrm{max}\{m_k:k<K\}`$. This ends the construction of the next element of the fusion sequence. For every $`i<\omega `$ if $`\sigma :F_i{}_{}{}^{n_i}2`$ is consistent with $`p_{i+1}`$ and extends $`\tau :F_{i1}{}_{}{}^{n_{i1}}2`$ then $$p_{i+1}\sigma \mathrm{`}\mathrm{`}p(\alpha ){}_{}{}^{M(\tau )}2=T_\sigma \mathrm{"}.$$ Considering our function $`\dot{\varphi }`$, let us denote the finite tree $`\dot{\varphi }(T_\sigma )`$ by $`S_\sigma `$. We have $$p_{i+1}\sigma \mathrm{`}\mathrm{`}\dot{\varphi }(T_\sigma )=S_\sigma \mathrm{"}.$$ When we are building the fusion sequence we can of course make sure that the fusion determines $`\dot{y}`$ as well. Suppose we have that $`p_i\tau _k\mathrm{`}\mathrm{`}t_{\tau _k}\dot{y}\mathrm{"}`$, $`t_{\tau _k}`$ of length $`i+1`$. With lemma 13 we can choose $`q_k`$ strong enough such that for every $`\sigma \mathrm{\Sigma }_{\tau _k}`$ consistent with $`q_k`$ we have a $`t_\sigma `$ of length $`i+2`$ such that $`q_k\sigma \mathrm{`}\mathrm{`}t_\sigma \dot{y}\mathrm{"}`$. So assume we have made sure this is the case and let us define the function $`f`$ in $`V`$ by $`f(b)=t_\sigma `$ for every maximal branch $`bS_\sigma `$ for every $`\sigma :F_i{}_{}{}^{n_i}2`$ consistent with $`p_i`$ for some $`i\omega `$. The function $`f`$ is well-defined by lemma 13 and we have for every $`i\omega `$ and $`\sigma :F_i{}_{}{}^{n_i}2`$ consistent with $`p_i`$ that $`p_i\sigma \mathrm{`}\mathrm{`}f([S_\sigma ])[t_\sigma ]\mathrm{"}`$ and thus $`q\mathrm{`}\mathrm{`}f(\dot{x})=\dot{y}\mathrm{"}`$. ###### Lemma 15. Given an ordinal $`\alpha `$, a $`p_{\alpha +1}`$ and $`_{\alpha +1}`$-names $`\dot{x}`$ and $`\dot{y}`$ such that $`p\mathrm{`}\mathrm{`}\dot{x},\dot{y}V_\alpha \mathrm{"}`$ then there exists a homeomorphism $`f`$, with code in $`V`$, and a $`qp`$ such that $`q\mathrm{`}\mathrm{`}f(\dot{x})=\dot{y}\mathrm{"}`$. Proof By applying remark 9 twice we have an $`rp`$ in $`_{\alpha +1}`$ and $`_{\alpha +1}`$ names $`\dot{\varphi }_x`$, $`\dot{\varphi }_y`$ and $`\dot{T}_x`$, $`\dot{T}_y`$ for maps and perfect trees respectively such that $`r\alpha \mathrm{`}\mathrm{`}\dot{\varphi }_x(p(\alpha ))=\dot{T}_x\text{and}\dot{\varphi }_y(p(\alpha ))=\dot{T}_y\mathrm{"}`$. Without loss of generality we can assume that $`p\alpha =r`$. During the construction of possible finite sub trees $`(T_x)_\sigma `$ for $`\dot{x}`$, when constructing the fusion sequence in the proof of lemma 14 we could of course at the same time also have constructed a similar sequence of finite sub trees $`(T_y)_\sigma `$ for $`\dot{y}`$. Without loss of generality we could also have made sure that in the proof of lemma 14 item 2 is replaced by 1. for every maximal branch $`t`$ in $`T_{\tau _0}`$ there are exactly two different branches of length $`M(\tau _0)`$ in $`T_\sigma `$ extending $`t`$, Following the proof of lemma 14 we have for every $`\sigma :F_i{}_{}{}^{n_i}2`$ consistent with $`p_{i+1}`$ finite sub trees $`S_\sigma ^x`$ and $`S_\sigma ^y`$ such that $$p_{i+1}\sigma \mathrm{`}\mathrm{`}\dot{\varphi }_x((T_x)_\sigma )=S_\sigma ^x\text{and}\dot{\varphi }_y((T_y)_\sigma )=S_\sigma ^y.$$ We are ready to define the homeomorphism $`f`$ in $`V`$ that maps $`x`$ onto $`y`$ in the extension. Suppose $`\tau :F_i{}_{}{}^{n_i}2`$ and $`\sigma :F_{i+1}{}_{}{}^{n_{i+1}}2`$ such that $`\sigma `$ extends $`\tau `$. Every maximal branch in $`(T_x)_\tau `$ corresponds to exactly one maximal branch in $`(T_y)_\tau `$. Let $`f`$ map the splitting point in $`(T_x)_\sigma `$ above any maximal branch in $`(T_x)_\tau `$ onto the splitting point in $`(T_y)_\sigma `$ above the corresponding maximal branch in $`(T_y)_\tau `$. The function $`f`$ thus defined will be a continuous and one-to-one mapping between two Cantor sets, so a homeomorphism. Furthermore the fusion $`q`$ forces that $`f`$ maps $`x`$ onto $`y`$ in the extension. ###### Lemma 16. Suppose that $`\alpha `$ is a limit ordinal of cofinality $`\mathrm{}_0`$. Let $`x`$ be a real in $`V_\alpha `$ such that $`x_{\beta <\alpha }V_\beta `$, and let $`p_\alpha `$ be a witness of this. Also let $`F,H[\mathrm{dom}(p)]^{<\omega }`$ such that $`FH`$ and let $`n`$ and $`m`$ be two integers such that $`m<n`$. If $`\tau :F{}_{}{}^{m}2`$ is consistent with $`p`$, and $`u_\tau {}_{}{}^{<\omega }2`$ is such that $$p\tau \mathrm{`}\mathrm{`}u_\tau \dot{x}\mathrm{"},$$ then there exists a $`(q,j)>_H(p\tau ,n)`$ such that for every $`\sigma :H{}_{}{}^{n}2`$ consistent with $`q`$, we have a $`u_\sigma {}_{}{}^{<\omega }2`$ such that $`q\sigma \mathrm{`}\mathrm{`}u_\sigma \dot{x}\mathrm{"}`$; in addition we have $`l(u_\sigma )=l(u_\varsigma )`$ and $`u_\sigma u_\varsigma `$ whenever $`\sigma `$ and $`\varsigma `$ are distinct and consistent with $`q`$. Before we prove the lemma we need some more notation. We let $`^{}`$ denote forcing in $`V_\delta `$ over $`_{\delta \alpha }`$. Here we use again the same notation as in \[BauLav79\] where for $`\delta <\alpha `$ $`P_{\delta \alpha }=\{p_\alpha :\mathrm{dom}(p)\{\xi :\delta \xi <\alpha \}\}`$, and if $`p_\alpha `$ then $`p^\delta =p(p\delta )_{\delta \alpha }`$. The mapping which carries $`p`$ into $`(p\delta ,p^\delta )`$ is an isomorphism of $`_\alpha `$ onto a dense subset of $`_\delta \times _{\delta \alpha }`$ (see \[BauLav79\]). Proof\[Lemma 16\] Choose a $`\delta `$ such that $`\mathrm{max}(H)<\delta <\alpha `$. Let $`\tau :F{}_{}{}^{m}2`$ be consistent with $`p`$ and let $`\mathrm{\Sigma }_\tau `$ denote all the $`\tau `$ extending functions $`\sigma :H{}_{}{}^{n}2`$. Because $`p`$ forces that $`xV_\delta `$, there is an antichain below $`p^\delta `$ of size $`\|\mathrm{\Sigma }_\tau \|`$ such that all these elements force different interpretations of $`\dot{x}`$ in the extension. In other words there exist a sequence $`\{\dot{f}_\sigma :\sigma \mathrm{\Sigma }_\tau \}`$ of $`_\delta `$ names for elements of $`_{\delta \alpha }`$ and a sequence $`\{\dot{u}_\sigma :\sigma \mathrm{\Sigma }_\tau \}`$ of $`_\delta `$ names for elements of $`{}_{}{}^{<\omega }2`$ such that for all $`\sigma \mathrm{\Sigma }_\tau `$ we have (1) $$(p\tau )\delta \mathrm{`}\mathrm{`}\dot{f}_\sigma p^\delta \text{and}\dot{f}_\sigma ^{}\mathrm{`}\mathrm{`}\dot{u}_\sigma \dot{x}\mathrm{"}\mathrm{"},$$ and if $`\sigma `$ and $`\varsigma `$ are distinct then (2) $$(p\tau )\delta \mathrm{`}\mathrm{`}l(\dot{u}_\sigma )=l(\dot{u}_\varsigma )\text{and}\dot{u}_\sigma \dot{u}_\varsigma \mathrm{"}.$$ Repeatedly using lemma 2.3 of \[BauLav79\] we see that there exist a $`(q^{},j)>_H((p\tau )\delta ,n)`$ and sequences $`\{f_\sigma :\sigma \mathrm{\Sigma }_\tau \}`$, $`\{u_\sigma :\sigma \mathrm{\Sigma }_\tau \}{}_{}{}^{i}2`$ for some integer $`i`$ such that for every $`\sigma \mathrm{\Sigma }_\tau `$ we have (3) $$q^{}_\delta \mathrm{`}\mathrm{`}\dot{f}_\sigma =f_\sigma \text{and}\dot{u}_\sigma =u_\sigma \mathrm{"}.$$ Now let $`q`$ denote the element of $`_\alpha `$ such that $`q\delta =q^{}`$, and $`(q\sigma )\delta \mathrm{`}\mathrm{`}q^\delta =f_\sigma \mathrm{"}`$ for every $`\sigma \mathrm{\Sigma }_\tau `$ consistent with $`q^{}`$. This completes the proof. ###### Theorem 17. For $`\alpha `$ a limit ordinal of cofinality $`\mathrm{}_0`$ and $`x`$ and $`y`$ reals in $`V_\alpha `$ such that $`x_{\beta <\alpha }V_\beta `$, there exist a continuous function $`f`$ defined in $`V`$ such that in $`V_\alpha `$ the equation $`f(x)=y`$ holds. If also $`y_{\beta <\alpha }V_\beta `$ then $`f`$ can be chosen to be a homeomorphism. Proof For the first part of the theorem suppose that we have $`p_\alpha `$ such that $`p\mathrm{`}\mathrm{`}\dot{x}_{\beta <\alpha }V_\beta \text{and}\dot{y}_{\beta <\alpha }V_\beta \mathrm{"}`$. We will construct a fusion sequence below $`p`$ and define a continuous function $`f`$ in $`V`$ such that the fusion of the sequence forces that $`f(x)=y`$ holds in $`V_\alpha `$. Let $`p_0=p_1=p`$, $`n_0=n_1=0`$, $`F_0=\mathrm{}`$, and choose $`F_1[\mathrm{dom}(p)]^{<\omega }`$ in such a way that we are building a fusion sequence. Suppose we have constructed the sequence up to $`i`$, we will construct the next element of the fusion sequence. Let $`\{\tau _k:k<K\}`$ denote an enumeration of all maps from $`F_{i1}`$ into $`{}_{}{}^{n_{i1}}2`$ consistent with $`p_i`$. According to lemma 16 there exists a $`(q_0,j_0)>_{F_i}(p_i\tau _0,n_i)`$ such that for every $`\sigma :F_i{}_{}{}^{n_i}2`$ consistent with $`q_0`$ we have distinct $`u_\sigma `$’s in $`{}_{}{}^{m(\tau _0)}2`$ (where $`m(\tau _0)`$ follows from lemma 16), such that $`q_0\sigma \mathrm{`}\mathrm{`}u_\sigma \dot{x}\mathrm{"}`$. Now use lemma 12 to construct $`r_0_\alpha `$ such that $`r_0q_0`$ and $`r_0\tau _0=q_0`$. We now iteratively consider all the $`\tau _k`$. In the general case if $`\tau _k`$ is not consistent with $`r_{k1}`$ then we make sure that $`r_k=r_{k1}`$ and $`j_k=j_{k1}`$. If $`\tau _k`$ is consistent with $`r_{k1}`$ we find by lemma 16 a $`(q_k,j_k)>_{F_i}(r_{k1}\tau _k,n_i)`$ such that for every $`\sigma :F_i{}_{}{}^{n_i}2`$ consistent with $`q_k`$ we have distinct $`u_\sigma `$’s in $`{}_{}{}^{m(\tau _k)}2`$ such that $`q_k\sigma \mathrm{`}\mathrm{`}u_\sigma \dot{x}\mathrm{"}`$. Now use lemma 12 to construct $`r_k_\alpha `$ such that $`r_kr_{k1}`$ and $`r_k\tau _k=q_k`$. After considering all $`\tau _k`$ we define $`p_{i+1}=r_{K1}`$ and $`n_{i+1}=\mathrm{max}\{j_k:k<K\}`$. If we take a closer look at lemma 16 we can also let the fusion sequence that we just constructed determine $`\dot{y}`$. Because if we have $`p\tau \mathrm{`}\mathrm{`}t_\tau \dot{y}\mathrm{"}`$, following the proof of lemma 16 we can make sure that (by some strengthening of $`q^{}`$ or the $`f_\sigma `$’s, if necessary) there exist $`t_\sigma `$’s in $`{}_{}{}^{<\omega }2`$, not necessarily distinct, extending $`t_\tau `$ such that for $`\sigma :H{}_{}{}^{n}2`$ consistent with $`q`$ we also have $`q\sigma \mathrm{`}\mathrm{`}t_\sigma \dot{y}\mathrm{"}`$. So assume we have done this. We have for every $`\sigma :F_i{}_{}{}^{n_i}2`$ consistent with $`p_{i+1}`$ (4) $$p_{i+1}\sigma \mathrm{`}\mathrm{`}u_\sigma \dot{x}\text{and}t_\sigma \dot{y}\mathrm{"}.$$ Now we are ready to define our function $`f`$ which will map $`x`$ in $`V_\alpha `$ continuously onto $`y`$. Let $`f([u_\sigma ])[t_\sigma ]`$ for all $`\sigma :F_i{}_{}{}^{n_i}2`$ and all $`i\omega `$. Then $`p_i\sigma \mathrm{`}\mathrm{`}f(\dot{x})[t_\sigma ]\mathrm{"}`$ for $`\sigma :F_i{}_{}{}^{n_i}2`$ consistent with $`p_i`$ and $`i\omega `$. It follows that the fusion $`q`$ forces that in $`V_\alpha `$ we have $`f(x)=y`$. Moreover $`f`$ is a continuous function, this follows from lemma 16. For the second part of the theorem suppose that $`p\mathrm{`}\mathrm{`}\dot{x},\dot{y}_{\beta <\alpha }V_\beta \mathrm{"}`$. Just as in lemma 16 we can choose not only the $`u_\sigma `$’s in equation 4 distinct but also the $`t_\sigma `$’s for $`\sigma \mathrm{\Sigma }_\tau `$ and $`\tau :F_i{}_{}{}^{n_i}2`$ for some $`i\omega `$. With this, the constructed continuous function $`f`$ is actually a homeomorphism. As there are no reals added at limit stages of cofinality larger than $`\mathrm{}_0`$ we have as a corollary to theorems 11 and 17 ###### Corollary 18. For every $`\alpha `$ and every $`\dot{x}`$ and $`\dot{y}`$ $`_\alpha `$-names for reals in $`V_\alpha _{\beta <\alpha }V_\beta `$ there exists a homeomorphism $`f`$ defined in $`V`$ such that in $`V_\alpha `$ we have $`f(x)=y`$. ###### Remark 19. It is not the case that the $`𝔱𝔣`$ number is the same for all compact metric spaces, e.g. every Cook continuum $`X`$ has $`𝔱𝔣(X)=𝔠`$ (it only has the identity and constant mappings as self-maps, see \[Coo67\]). On the other hand, in the Sacks model one has $`𝔱𝔣(C)=𝔱𝔣()=𝔱𝔣([0,1])=\mathrm{}_1`$. To see this, observe that our proof produces, given $`x`$ and $`y`$, two copies of the Cantor set $`A`$ and $`B`$ containing $`x`$ and $`y`$ respectively and a continuous map $`f:AB`$, say, with $`f(x)=y`$. One can then extend $`f`$ to a continuous map $`\stackrel{~}{f}:[0,1][0,1]`$ (or $`\stackrel{~}{f}:`$), whose code will still be in $`V`$. ###### Remark 20. If $`\mathrm{cov}(nowheredense)=𝔠`$ for the unit interval $`I`$, then remark 4 shows that $`𝔱𝔣(I)=𝔠`$. Suppose that $`\mathrm{cov}(nowheredense)=\kappa <𝔠`$, for $`I`$, then we can cover $`I`$ by $`\kappa `$ many Cantor sets $`\{C_\alpha \}_{\alpha <\kappa }`$ in such a way that for every two reals $`x`$ and $`y`$ there exists an $`\alpha `$ such that $`x,yC_\alpha `$. For every $`\alpha `$ we have a transitive family of continuous functions $`_\alpha `$ on $`C_\alpha `$ such that $`\|_\alpha \|=𝔱𝔣(C)`$. We can extend every $`f_\alpha `$ to a continuous self map $`\stackrel{~}{f}`$ of $`I`$. So $`=\{\stackrel{~}{f}:\text{there is an}\alpha <\kappa \text{and}f_\alpha \}`$ is a transitive set of continuous functions on $`I`$, and its cardinality is less than or equal to $`\kappa \times 𝔱𝔣(C)=𝔱𝔣(C)`$. So if we can cover the unit interval with less than $`𝔠`$ many nowhere dense sets we have $`𝔱𝔣(I)𝔱𝔣(C)`$. ## 4. The cardinal $`𝔱𝔣`$ and side-by-side Sacks forcing In this paper we showed that after adding $`\mathrm{}_2`$ many Sacks reals iteratively to a model of $`ZFC+CH`$ we end up with a model of $`𝔱𝔣<𝔠`$. Now consider $`𝕊(\kappa )`$, the poset for adding $`\kappa `$ many Sacks reals side-by-side (see \[Bau85\]). We have that $`𝕊(\kappa )`$ has the $`(2^\mathrm{}_0)^+`$-chain condition and preserves $`\mathrm{}_1`$. Suppose that $`\kappa \mathrm{}_1`$ and $`\mathrm{cf}(\kappa )\mathrm{}_1`$. If $`V`$ is a model of $`CH`$ and $`G`$ is $`𝕊(\kappa )`$-generic over $`V`$, we have in $`V[G]`$ that $`2^\mathrm{}_0=\kappa `$ and all cardinals are preserved. A natural question would be if we get a model of $`𝔱𝔣<𝔠`$ when we add $`\mathrm{}_2`$ many Sacks reals side-by-side to a model of $`ZFC+CH`$. The answer to this question is in the negative. Suppose that $`V`$ is a model of $`ZFC`$. Consider the poset $`=𝕊(\{1,2,3,4\})`$ that adds four Sacks reals side-by-side to the model $`V`$. We define $`_1`$ to be the p.o.-set $`𝕊(\{1,2\})`$ and $`_2`$ to be the p.o.-set $`𝕊(\{3,4\})`$. Suppose $`G`$ is $``$ generic over $`V`$ then $`G_{12}=G\{1,2\}`$ is $`_1`$ generic and $`G_{34}=G\{3,4\}`$ is $`_2`$ generic over $`V`$. The following holds. ###### Lemma 21. In $`V[G]`$ we have $`V[G_{12}]V[G_{34}]=V`$. Proof Suppose that $`\dot{X}`$ is a $``$ name and $`q`$ an element of $``$ such that $`q\mathrm{`}\mathrm{`}\dot{X}V[G_{12}]V[G_{34}]\mathrm{"}`$. So there exists a $`_1`$ name $`\dot{Y}`$ and a $`_2`$ name $`\dot{Z}`$ such that $`q\mathrm{`}\mathrm{`}\dot{X}=\dot{Y}=\dot{Z}\mathrm{"}`$. Aiming for a contradiction assume $`\dot{X}`$ is a name for an object not in $`V`$. There exists a $`n\omega `$ such that $`q`$ does not decide $`n\dot{X}`$. Now we have $`q_1=q\{1,2\}`$ does not decide $`n\dot{Y}`$, and $`q_2=q\{3,4\}`$ does not decide $`n\dot{Z}`$. So we can find in $`_1`$ a $`rq_1`$ such that $`r\mathrm{`}\mathrm{`}n\dot{Y}\mathrm{"}`$ and in $`_2`$ a $`tq_2`$ such that $`t\mathrm{`}\mathrm{`}n\dot{Z}\mathrm{"}`$. This gives the contradiction we are looking for because $`rt\mathrm{`}\mathrm{`}\dot{Y}\dot{Z}\mathrm{"}`$ and $`rtq`$. So $`\dot{X}`$ must be a name of an element in $`V`$. Now we can prove that adding $`\mathrm{}_2`$ many Sacks reals to a model of $`ZFC+CH`$ we do not produce a model of $`𝔱𝔣<𝔠`$. ###### Theorem 22. Suppose $`VCH`$ and $`G`$ is a $`𝕊(\kappa )`$-generic filter over $`V`$, where $`\kappa \mathrm{}_1`$ and $`\mathrm{cf}(\kappa )\mathrm{}_1`$, then $`V[G]𝔱𝔣=𝔠`$. Proof For every $`\alpha <\beta <\kappa `$ we have that there exists a function $`f_{\alpha ,\beta }V[G\{\alpha ,\beta \}]`$ mapping $`s_\alpha `$ onto $`s_\beta `$ or vice versa. This function $`f_{\alpha ,\beta }`$ is not a member of $`V`$ for the obvious reason that assuming that $`f_{\alpha ,\beta }`$ maps $`s_\alpha `$ onto $`s_\beta `$ we get $`s_\beta V[G\{\alpha \}]`$, which, of course, is false. Using lemma 21 and the fact that $`2\kappa =\kappa `$ we see that the size of $`𝔱𝔣`$ is at least $`\kappa `$, because $`f_{2\alpha ,2\alpha +1}f_{2\beta ,2\beta +1}`$ for every $`\alpha \beta `$. By theorem 3 we are done.
warning/0004/math0004058.html	ar5iv	text	# Embedding Obstructions and 𝟒-dimensional Thickenings of 𝟐-complexes ## 1. Introduction By general position any $`n`$-dimensional simplicial complex $`K`$ PL embeds into $`^{2n+1}`$, while the image of a generic map of $`K`$ into $`^{2n}`$ has a finite number of double points. By counting double points of an immersion one gets the cohomological obstruction to embeddability of an $`n`$-complex into $`^{2n}`$, introduced by Van Kampen . He also constructed for each $`n`$ examples which do not admit an embedding. An application of Whitney trick shows that this obstruction is complete for $`n>2`$, see , , , . It follows from Kuratowski’s planarity criterion for graphs that this result also holds for $`n=1`$. The remaining case, $`n=2`$, was open until recently when the obstruction was shown in to be incomplete. This paper is centered around the question of embeddability of $`2`$-complexes in $`^4`$, and is motivated by the result of . We define for $`2`$-complexes $`K`$ with $`H_1(K;)=0`$ a sequence of higher embedding obstructions $`\{o_m(K)\}`$, using Massey products on the boundary of a $`4`$-dimensional thickening $`M^4`$ of $`K`$. Roughly, Van Kampen’s obstruction corresponds in this setting to the intersection pairing on $`M`$, modulo the choice of a thickening $`M`$. Since different thickenings may give different Massey products, $`\{o_m(K)\}`$ in general are subsets of the corresponding cohomology groups; $`o_{m+1}(K)`$ is defined if $`o_m(K)`$ contains zero. If $`K`$ embeds into $`^4`$ then $`0o_m(K)`$ for each $`m`$. We prove that these higher obstruction detect non-embeddability of the family of examples introduced in , by showing that $`o_m(K)`$ does not contain zero for some $`m`$. Our proof uses the result of Conway - Gordon and Sachs that any embedding of a complete graph on $`6`$ vertices into $`S^3`$ contains two disjoint linking cycles (, ). In the simplest relative case, for the disjoint union of $`2`$-disks with a prescribed embedding of their boundaries into $`S^3`$, by a result of Turaev Massey products correspond to Milnor’s $`\overline{\mu }`$-invariants of the link in $`S^3`$, so our obstructions may be thought of as an absolute analogue of $`\overline{\mu }`$-invariants. As in the case of $`\overline{\mu }`$-invariants (for example, Whitehead double of the Hopf link is not a slice link, while all $`\overline{\mu }`$-invariants vanish), one does not expect that the entire sequence of obstructions defined here is complete, although no examples are known at this time. The question about $`2`$-complexes has an additional subtlety, being in piecewise-linear category, where embeddings are not necessarily locally flat. The definition of Van Kampen’s obstruction is recalled in section 2. In section 3 we prove its reformulation in the context of thickenings, and we introduce the sequence of higher obstructions $`\{o_m(K)\}`$. We review the examples of $`2`$-complexes in in section 4, and we compute the obstructions for these examples. Section 5 gives a reformulation of Van Kampen’s obstruction in terms of configuration spaces, which suggests another approach to defining higher embedding obstructions. The present study of the embedding problem for $`2`$-complexes in $`^4`$ is motivated, in part, by the $`4`$-dimensional topological surgery conjecture, via its (A,B)-slice reformulation . More precisely, the surgery conjecture is equivalent to the relative embedding question for a certain family of $`4`$-dimensional handlebodies – “thickenings” of $`2`$-complexes in the sense of section 3. However, many interesting examples of these handlebodies have non-trivial first homology, and for this application the obstructions $`\{o_i(K)\}`$ need to be extended to the general case. ## 2. Van Kampen’s obstruction In this section we briefly review the definition of Van Kampen’s obstruction, more details are given in . In 1933 Van Kampen introduced an obstruction $`o(K)H_{/2}^{2n}(K^{};)`$ to piecewise-linear embeddability of an $`n`$-dimensional simplicial complex $`K`$ into $`^{2n}`$. The cohomology in question is $`/2`$-equivariant cohomology where $`/2`$ acts on the deleted product $`K^{}=K\times K\mathrm{\Delta }`$ of a complex $`K`$ by exchanging the factors of $`K^{}`$ and acts on the coefficients by multiplication with $`(1)^n`$. The diagonal $`\mathrm{\Delta }`$ consists of all products $`\sigma \times \tau `$ such that simplices $`\sigma `$ and $`\tau `$ have at least one vertex in common. Note that for $`n`$ even (in particular, in the case of main interest in this paper, $`n=2`$) the action of $`/2`$ on the coefficients is trivial, and $`o(K)`$ is an element of the ordinary cohomology group $`H^{2n}(K^{}/(/2);)`$. Let $`f`$ be any PL immersion of $`K`$ into $`^{2n}`$. The obstruction is defined on the cochain level by counting algebraic intersection numbers of the images of disjoint top-dimensional simplices of $`K`$: $`o_f(\sigma ^n\times \tau ^n)=f(\sigma )f(\tau )`$. Here $`\sigma \times \tau `$ is viewed as an oriented generator of the $`2n`$-th chain group of $`K\times K\mathrm{\Delta }`$. The cohomology class $`o(K)`$ of $`o_f`$ is independent of the chosen immersion $`f`$. Clearly $`o(K)`$ is trivial if $`K`$ embeds into $`^{2n}`$. Shapiro and Wu made this definition precise and proved, using the Whitney trick, the converse in dimension $`n`$ greater than $`2`$. ###### Theorem 2.1 (, , , ). For $`n2`$ an $`n`$-dimensional simplicial complex $`K`$ admits an embedding into $`^{2n}`$ if and only if $`o(K)`$ vanishes. See for a modern exposition of the proof for $`n>2`$. For $`n=1`$ this theorem follows from Kuratowski’s planarity criterion. The obstruction in the remaining case, for $`n=2`$, was shown to be incomplete in . We recall the construction of examples in in section 4. ## 3. Obstructions via $`4`$-dimensional thickenings of $`2`$-complexes In this section we give a rational reformulation of Van Kampen’s obstruction $`o(K)`$ in terms of thickenings of $`K`$, and we introduce a sequence of higher embedding obstructions $`\{o_m(K)\}`$ for $`2`$-complexes whose rational first homology vanishes. Throughout this section all coefficients are $``$, unless stated otherwise, and $`K`$ denotes a simplicial $`2`$-complex. ###### Definition 3.1. A thickening of $`K`$ is a smooth $`4`$-manifold $`M`$ with boundary, obtained by replacing each $`i`$-simplex of $`K`$ with a $`4`$-dimensional $`i`$-handle, $`i=0,1,2`$. The attaching map of each $`2`$-handle is required to be isotopic, within the union of $`0`$\- and $`1`$-handles, to the attaching map of the corresponding $`2`$-dimensional simplex. In general, $`K`$ may have different thickenings depending on the choice of attaching maps of the $`2`$-handles. For example, $`S^2\times S^24`$-cell and the boundary-connected sum $`S^2\times D^2\mathrm{}S^2\times D^2`$ are both thickenings of $`S^2S^2`$. The intersection pairing on $`M`$ defines an element $`\iota Hom(H_2(M)H_2(M),)`$. Let $`\overline{\iota }H^4(K\times K\mathrm{\Delta };)`$ denote its image under the homomorphism $$Hom(H_2(M)H_2(M),)Hom(H_2(K)H_2(K),)$$ $$H^4(K\times K)H^4(K\times K\mathrm{\Delta })$$ where the last map is induced by inclusion. ###### Theorem 3.2. The image of the (rational) Van Kampen’s obstruction $`o(K)`$ under the homomorphism induced by the quotient map $$H_{/2}^4(K\times K\mathrm{\Delta };)H^4(K\times K\mathrm{\Delta };)$$ coincides with $`\overline{\iota }`$. Proof. Suppose a thickening $`M`$ is induced by an immersion $`f:K^4`$, so that $`f`$ extends to an immersion $`M^4`$. By subdividing the complex $`K`$, if necessary, one may assume that $`f(\sigma )f(\tau )=\mathrm{}`$ for all (open) $`2`$-simplices $`\sigma \tau `$ with $`\sigma \times \tau \mathrm{\Delta }`$, and $`f\|_\sigma `$ is an embedding for each $`\sigma `$. Let $`\overline{o}_f:C_4(K\times K)`$ denote the extension by zero on the diagonal of Van Kampen’s cochain $`o_f:C_4(K\times K\mathrm{\Delta })`$. It suffices to prove that $`[\overline{o}_f]`$ and $`\iota `$ define identical elements in $`Hom(H_2(K)H_2(K),)`$. Let $`a`$, $`b`$ be two classes in $`H_2(K)`$ and let $`\alpha =\mathrm{\Sigma }\alpha _i\sigma _i`$, $`\beta =\mathrm{\Sigma }\beta _i\sigma _i`$ be their cycle representatives, where $`\{\sigma _i\}`$ is the set of $`2`$-simplices of $`K`$. In order to compute $`ab`$ in $`M`$, perturb $`\alpha `$ and $`\beta `$ to $`\stackrel{~}{\alpha }`$ and $`\stackrel{~}{\beta }`$ which intersect each other transversely (in a finite number of double points). The intersection number of two cycles $`f(\stackrel{~}{\alpha })`$ and $`f(\stackrel{~}{\beta })`$ in $`^4`$ is trivial. On the other hand, $`f(\stackrel{~}{\alpha })f(\stackrel{~}{\beta })`$ may be computed as the sum of two terms: one is the intersection number of $`\stackrel{~}{\alpha }`$ and $`\stackrel{~}{\beta }`$ in $`M`$, the other is obtained by considering the intersections of $`f(\stackrel{~}{\alpha })`$ and $`f(\stackrel{~}{\beta })`$ in $`^4`$, which are singular points of $`f`$. This last term is equal to $`\overline{o}_f(\alpha \times \beta )`$, and this proves $$[\overline{o}_f]=\iota :H_2(K)H_2(K).$$ The restriction of $`[\overline{o}_f]`$ to $`H^4(K\times K\mathrm{\Delta };)`$ coincides with $`o(K)`$, thus the result is proved for thickenings induced by immersions. In general not every thickening of $`K`$ may be immersed into $`^4`$. Let $`M^4`$ be an arbitrary thickening of $`K`$ and let $`f:K^4`$ be any immersion. Again one may assume that $`f(\sigma )f(\tau )=\mathrm{}`$ if $`\sigma \times \tau \mathrm{\Delta }`$, $`\sigma \tau `$, and $`f\|_\sigma `$ is an embedding for each simplex $`\sigma `$. The immersion $`f`$ extends to an embedding of $`0`$\- and $`1`$-handles of $`M`$. There is an integer obstruction to extending it over each $`2`$-handle, due to a possible difference of the framing of the $`2`$-handle and of the normal bundle of the $`2`$-simplex in $`^4`$. However, each $`2`$-handle may be mapped into $`^4`$ as a bundle over the corresponding $`2`$-simplex, pinched over several points. The proof, given above in the case of an immersion, carries through, if one extends $`o_f`$ to $`\overline{o}_f`$ by setting $`\overline{o}_f(\sigma \times \sigma )`$ to be equal to the difference in framings, discussed above, and setting $`\overline{o}_f(\sigma \times \tau )=0`$ for all $`\sigma \times \tau \mathrm{\Delta }`$, $`\sigma \tau `$. ∎ ###### Remark 3.3. In general the intersection pairing varies within the homotopy type of a $`4`$-manifold $`M`$. In the example above the intersection pairing on $`S^2\times D^2\mathrm{}S^2\times D^2`$ is trivial, while the pairing on $`S^2\times S^24`$-cell is non-degenerate. However, theorem 3.2 shows that the pull-back of the intersection pairing on thickenings to a cohomology class on $`K\times K\mathrm{\Delta }`$ is an invariant of $`K`$, which coincides with the image of the (negative) Van Kampen’s obstruction. As a corollary to the proof of theorem 3.2, we have the following result. ###### Lemma 3.4. Let $`K`$ be a $`2`$-complex such that Van Kampen’s obstruction $`o(K)`$ vanishes. Then there is a $`4`$-dimensional thickening $`M`$ of $`K`$ with the trivial intersection pairing $`\iota =0Hom(H_2(M)H_2(M);)`$. Proof. Any cochain representative of the obstruction $`o(K)`$ is given by $`o_f`$ for some immersion $`f`$, see or . Since $`o(K)`$ vanishes, there exists an immersion $`f:K^4`$, giving rise to the trivial Van Kampen’s cochain $`o_f=0`$. Let $`M`$ denote the thickening induced by $`f`$. It follows from the proof of theorem 3.2 that if one extends $`o_f`$ by zero on the diagonal to a cochain $`\overline{o}_f`$ on $`K\times K`$, then $`\iota =[\overline{o}_f]=0Hom(H_2(M)H_2(M);)`$. ∎ Before introducing the higher embedding obstructions, we recall the definition of Massey products. See for proofs and additional properties. ###### Definition 3.5. Let $`X`$ be a space, and let $`\alpha _1,\mathrm{},\alpha _m`$ be elements in $`H^1(X)`$. Suppose there is a collection of 1-cochains $`S=\{c_{ij}C^1(X)\|1ijm,(i,j)(1,m)\}`$ satisfying $$[c_{ii}]=\alpha _i\mathrm{for}\mathrm{each}i=1,\mathrm{},m,$$ $$\delta c_{ik}=\underset{j=i}{\overset{k1}{}}c_{ij}c_{j+1,k}\mathrm{for}i<k.$$ Then the cochain $`_{j=1}^{m1}c_{1j}c_{j+1,m}`$ is a cocycle, and its cohomology class in $`H^2(X)`$ is called the Massey product of $`\alpha _1,\mathrm{},\alpha _m`$ defined by the system $`S`$. The set of Massey products corresponding to all such defining systems is denoted by $`<\alpha _1,\mathrm{},\alpha _m>H^2(X)`$. Massey product of two elements is just a cup product. Note that given some classes $`\alpha _1,\mathrm{},\alpha _m`$, $`<\alpha _1,\mathrm{},\alpha _m>`$ is not necessarily defined. However, if all Massey products of less than $`m`$ elements vanish, then for any $`\alpha _1,\mathrm{},\alpha _mH^1(X)`$, $`<\alpha _1,\mathrm{},\alpha _m>`$ is a well-defined element. The following lemma justifies our definition of higher embedding obstructions. ###### Lemma 3.6. Let $`M`$ be a $`4`$-manifold with boundary and with $`H_1(M;)=0`$, and suppose $`M`$ admits an embedding into $`^4`$. Then all Massey products on $`H^1(M;)`$ vanish. Proof. Let $`N`$ denote the complement $`^4M`$. By Alexander duality, $`H_2(N)`$ and $`H^2(N)`$ are trivial. The map $`i^{}:H^1(N)H^1(M)`$ in the cohomology sequence of the pair $`(N,M)`$ is an isomorphism, since by assumption and by Poincaré duality $`H^1(N,M)H_3(N)`$ and $`H^2(N,M)H_2(N)`$ are trivial. Assume inductively that all Massey products of length less than $`m`$ vanish for some $`m2`$; then for any $`\alpha _1,\mathrm{},\alpha _mH^1(M)`$ one has $$<\alpha _1,\mathrm{},\alpha _m>=i^{}<(i^{})^1\alpha _1,\mathrm{},(i^{})^1\alpha _m>H^2(M).$$ However, this is the image of an element in $`H^2(N)=0`$, and the result follows. ∎ Let $`K`$ be a $`2`$-complex with $`H_1(K;)=0`$, and assume $`o(K)`$ vanishes. Let $`M`$ be a thickening of $`K`$ with trivial intersection pairing (its existence is given by lemma 3.4.) Note that the map $`H_2(M)H_2(M)`$ is an isomorphism, since by assumption on $`K`$, $`H_3(M,M)H^1(M)=0`$, and the map $`H_2(M)H_2(M,M)`$ is trivial by assumption on the intersection pairing. We now give the definition of higher embedding obstructions. Let $`a_1,a_2,a_3`$ be classes in $`H_2(K)`$, and let $`\alpha _1,\alpha _2,\alpha _3H^1(M)`$ denote their images under the isomorphism $$H_2(K)H_2(M)H_2(M)H^1(M).$$ The triple cup product $`(\alpha _1\alpha _2\alpha _3)[M]`$ defines a homomorphism $`H_2(K)H_2(K)H_2(K)`$, and an element $`o_3(K,M)H^6(K\times K\times K)`$. The cohomology class $`o_3(K,M)`$ depends in general on the choice of a thickening $`M`$, thus we define the third obstruction $`o_3(K)`$ to be the subset $`\{o_3(K,M)\}H^6(K^3;)`$ where $`M`$ is to vary over all thickenings of $`K`$ with trivial intersection pairing. Note that if $`o_3(K)`$ is defined and contains zero, then there is a thickening $`M`$ of $`K`$ such that all cup products on $`H^2(M)`$ vanish. ###### Definition 3.7. Define $`o_2(K)`$ to be the Van Kampen’s obstruction $`o(K)`$. If $`o_2(K)`$ vanishes, then $`o_3(K)H^6(K^3)`$ is defined as above. Assume by induction that for some $`m>3`$ there is a thickening $`M`$ of $`K`$ such that $`o_{m1}(K,M)`$ is defined and is equal to zero (equivalently, the intersection pairing on $`M`$ is trivial, and all Massey products on $`H^1(M)`$ of at most $`(m2)`$ elements vanish.) Let $`a_1,\mathrm{},a_m`$, be classes in $`H_2(K)`$, and let $`\alpha _1,\mathrm{},\alpha _m`$ denote the corresponding elements in $`H^1(M)`$. The class $`o_m(K,M)H^{2m}(K^m;)`$ is defined by the homomorphism $$H_{2m}(K^m)_1^mH_2(K)_1^mH^1(M)$$ which sends $`a_1\mathrm{}a_m`$ to $`(<\alpha _1,\mathrm{},\alpha _{m1}>\alpha _m)[M]`$. Here since all Massey products on $`H^1(M)`$ of less than $`(m1)`$ elements vanish, $`<\alpha _1,\mathrm{},\alpha _{m1}>H^2(M)`$ is a well-defined element. ###### Definition 3.8. The obstruction $`o_m(K)`$ is defined to be the subset $$\{o_m(K,M)\}H^{2m}(K^m),$$ where $`M`$ is to vary over all thickenings such that $`o_{m1}(K,M)=0`$. Note that $`o_m(K)`$ is defined if $`o_{m1}(K)`$ is defined and contains zero. Lemma 3.6 implies the following corollary. ###### Corollary 3.9. Let $`K`$ be a $`2`$-complex with $`H_1(K;)=0`$. If $`K`$ admits an embedding into $`^4`$ then $`o_m(K)`$ is defined and contains zero for each $`m`$. In section 4 we show that $`o_m(K)`$ does not contain zero for some $`m`$ for examples in , thus giving another proof that they do not embed into $`^4`$. The relative embedding problem. Let $`K`$ be a $`2`$-complex with $`H_1(K;)=0`$, and let $`L`$ be a $`1`$-dimensional subcomplex of $`K`$ with a prescribed embedding $`\varphi :LS^3`$. Consider the relative embedding problem: does there exist an embedding $`KB^4`$ which extends $`\varphi `$? Denote $`B^4_\varphi `$(thickening of $`K`$) by $`M`$, where thickening is taken in the sense of Definition 3.1. Let $`K_L^m`$ denote the subset in $`K^m`$ consisting of all $`m`$-tuples $`(x_1,\mathrm{},x_m)`$ such that $`x_iL`$ for some $`i`$. Assume that $`\pi _0(L)\pi _0(K)`$ is injective to have $`H_1(M)=0`$. Analogously to the absolute case, Massey products on $`M`$ define an element, depending on $`M`$, in the relative cohomology group $`H^{2m}(K^m,K_L^m)`$. Let $`o_m(K,L,\varphi )`$ denote the set of these elements in $`H^{2m}(K^m,K_L^m)`$, where $`M`$ is to vary over all thickenings for which the $`(m1)`$-st obstruction is zero. If there is an embedding of $`K`$ into $`B^4`$, extending $`\varphi `$, clearly there is a $`4`$-dimensional thickening $`M`$ which embeds into $`S^4`$, so $`0o_m(K,L,\varphi )`$ for each $`m`$. Consider the simplest relative case, when $`(K,L)=(D^2\mathrm{}D^2,S^1\mathrm{}S^1)`$. By the result of Turaev , the first non-trivial obstruction coincides in this case with the first non-trivial Milnor’s $`\overline{\mu }`$-invariants of the link $`\varphi (L)`$ in $`S^3`$. In this sense the obstructions $`o_m(K)`$ may be thought of as an absolute analogue of $`\overline{\mu }`$-invariants. However, since $`2`$-complexes in general have a more complicated topology, $`\{o_m(K)\}`$ have a larger indeterminacy. ## 4. Examples First we recall the construction of examples in . Let $`C`$ denote the $`2`$-skeleton of the $`6`$-simplex with vertices $`v_1,\mathrm{},v_7`$, with one $`2`$-cell, with vertices $`v_1v_2v_3`$, removed. Take another copy $`C^{}`$ of $`C`$, with vertices $`v_1^{},\mathrm{},v_7^{}`$, and denote by $`\overline{C}`$ the union of $`C`$ and $`C^{}`$, identified along their last vertices, $`v_7=v_7^{}`$. This $`2`$-complex is easily seen to admit an embedding into $`^4`$ (see ). Let $`\gamma `$ (resp. $`\gamma ^{}`$) denote the loop $`v_7v_1v_2v_3v_1v_7`$ (resp. $`v_7v_1^{}v_2^{}v_3^{}v_1^{}v_7`$) in $`\overline{C}`$. Denote by $`F`$ the free group on two generators, and fix a positive integer $`m`$. Let $`\alpha `$ be an element in $`F^m`$, the $`m`$-th term of the lower central series of $`F`$. We identify $`F`$ with $`\pi _1(\gamma \gamma ^{})`$, and we associate to each word in $`F`$ its “standard” representative loop in the wedge of two circles $`\gamma \gamma ^{}`$. Finally, we construct the $`2`$-complex $`K_\alpha `$ by attaching a $`2`$-cell to $`\overline{C}`$ along $`\alpha `$. ###### Theorem 4.1 (). Let $`\alpha `$ be a non-trivial element in $`F^m`$ for some $`m2`$. Then Van Kampen’s obstruction $`o(K_\alpha )`$ vanishes, but the $`2`$-complex $`K_\alpha `$ does not admit an embedding into $`^4`$. We now present a computation of the obstructions $`\{o_i(K_\alpha )\}`$. The class $`m`$ of the commutator $`\alpha `$ is reflected in non-vanishing of the obstruction $`o_{m+1}(K_\alpha )`$. ###### Theorem 4.2. Let $`\alpha `$ be an element in $`F^m`$ for some $`m2`$, and assume $`\alpha F^{m+1}`$. Then $`o_{m+1}(K_\alpha )`$ is defined and does not contain zero. In particular, $`K_\alpha `$ does not admit an embedding into $`^4`$. Proof. First we construct a thickening $`M`$ of $`K_\alpha `$ with trivial intersection pairing and such that $`o_i(K,M)=0`$ for all $`im`$. The complex $`K_\alpha `$ is obtained from $`\overline{C}=CC^{}`$ by attaching a $`2`$-cell along the commutator $`\alpha F^m`$. Van Kampen constructed in an immersion of the $`2`$-skeleton of the $`6`$-simplex with vertices $`v_1,\mathrm{},v_7`$ into $`^4`$ such that the $`2`$-cells with vertices $`v_1v_2v_3`$ and $`v_4v_5v_6`$ intersect in one point, and all other simplices are disjoint and embedded. Consider the corresponding embedding of $`\overline{C}`$, and let $`\overline{M}`$ denote its thickening in $`^4`$. Clearly the intersection pairing on $`\overline{M}`$ is trivial, and all Massey products on $`H^1(\overline{M})`$ vanish. Recall that $`C`$ and $`C^{}`$ have the vertex $`v_7`$ in common. Consider the handle decomposition of $`\overline{M}`$, given by thickenings of simplices of $`K_\alpha `$ in $`^4`$. The union of the handles in $`\overline{M}`$ corresponding to all simplices in $`\overline{C}`$, containing $`v_7`$, is a $`4`$-ball $`B`$. The remaining $`2`$-handles are attached to $`B`$ along a link $`\overline{L}`$ in $`S^3=B`$. Each attaching curve is isotopic, within the union of $`0`$\- and $`1`$-handles, to the boundary curve of the corresponding $`2`$-simplex. There is no $`2`$-cell attached to $`v_1v_2v_3`$, however we introduce in $`S^3`$ a circle, isotopic to it. Because of the choice of the embedding of $`K`$ into $`^4`$, $`\overline{L}`$ is a slice link, and the curves isotopic to $`v_1v_2v_3`$ and $`v_4v_5v_6`$ (respectively $`v_1^{}v_2^{}v_3^{}`$ and $`v_4^{}v_5^{}v_6^{}`$) have linking number one. The remaining $`2`$-cell of $`K_\alpha `$ is attached along the commutator of $`v_1v_2v_3`$ and $`v_1^{}v_2^{}v_3^{}`$. Choosing appropriately the corresponding curve $`l`$ in $`S^3`$, we get the link $`L=\overline{L}l`$ such that all Milnor’s $`\overline{\mu }`$-invariants of $`L`$ of length less than $`m+1`$ vanish, and a $`\overline{\mu }`$-invariant of length $`m+1`$ of the $`3`$-component link $`(v_4v_5v_6,v_4^{}v_5^{}v_6^{},l)`$ is non-trivial. It is a result of Turaev that the first non-vanishing $`\overline{\mu }`$-invariant is equal to the corresponding Massey product on $`M`$, where $`M=\overline{M}_l2`$-handle. (Note that the framings of the components of $`\overline{L}`$ are zero, since the intersection pairing on $`\overline{M}`$ vanishes, and we choose the framing of $`l`$ also to be zero.) This proves that $`o_i(K,M)=0`$ for all $`im`$, and $`o_{m+1}(K_\alpha ,M)0`$. It remains to show that $`o_{m+1}(K_\alpha )`$ does not contain zero. Let $`M`$ be any thickening with trivial intersection pairing and with $`o_m(K,M)=0`$. As above, the union of the handles in $`M`$ corresponding to all simplices, containing $`v_7`$, is a $`4`$-ball $`B`$. By a theorem of Conway - Gordon and Sachs any embedding of the complete graph on $`6`$ vertices in $`S^3`$ contains two disjoint linking cycles. Consider the complete graph on vertices $`v_1,\mathrm{},v_6`$ in $`\overline{C}`$. According to the definition of $`M`$, the attaching curves of the $`2`$-handles in $`S^3=B`$ are isotopic to the attaching maps of simplices of $`K_\alpha `$. As above, we introduce in $`S^3`$ a curve isotopic to $`v_1v_2v_3`$. Now we have in $`S^3`$ a perturbed version of the complete graph on $`6`$ vertices. Since, according to definition 3.1, these perturbations take place in the union of $`0`$\- and $`1`$-handles, at least two of the curves must have a non-trivial linking number. Since the intersection pairing on $`M`$ vanishes, these two circles are the ones isotopic to $`v_1v_2v_3`$ and to $`v_4v_5v_6`$. Similarly we have in $`S^3`$ another, disjoint, copy of a perturbed graph on $`v_1^{},\mathrm{},v_6^{}`$, and two linking circles isotopic to $`v_1^{}v_2^{}v_3^{}`$ and to $`v_4^{}v_5^{}v_6^{}`$. Recall that there are no $`2`$-handles attached to $`v_1v_2v_3`$ or $`v_1^{}v_2^{}v_3^{}`$, however there is a $`2`$-handle whose attaching curve $`l`$ is a commutator of these circles. It is easily seen that the link $`(v_4v_5v_6,v_4^{}v_5^{}v_6^{},l)`$ has a non-trivial $`\overline{\mu }`$-invariant of length $`m+1`$. As above, this is translated into non-vanishing of $`o_{m+1}(K,M)`$. ∎ ###### Remark 4.3. The idea of the proof of the fact that any embedding of the complete graph on $`6`$ vertices in $`S^3`$ contains two linking cycles (, ) is conceptually similar to the proof of Van Kampen that the $`2`$-skeleton of the $`6`$-simplex does not embed into $`^4`$ . In both cases one shows that a certain number is invariant mod $`2`$ for different maps - in one case, the total linking number, in the other case, the total number of singular points of an immersion. In this sense our proof of theorem 4.2 is similar to the proof of theorem 4.1 in . ## 5. A note on configuration spaces In this section we consider an approach to the embedding problem, suggested by obstruction theory and configuration spaces. We give a reformulation of Van Kampen’s obstruction in this context, which suggests another approach to defining higher embedding obstructions. Given a space $`X`$, $`C^m(X)`$ will denote its configuration space of $`m`$ points: $$C^m(X)=\{(x_1,\mathrm{},x_m)X^m\|x_ix_j\mathrm{if}ij\}.$$ In the simplicial category, for a complex $`K`$ we define $$C^m(K)=\{\sigma _1\times \mathrm{}\times \sigma _mK^m\|\mathrm{simplices}\sigma _i,\sigma _j$$ $$\mathrm{have}\mathrm{no}\mathrm{vertices}\mathrm{in}\mathrm{common}\mathrm{for}ij\}.$$ The configuration space of two points $`C^2(X)`$ is sometimes called deleted product and is also denoted by $`X^{}`$. The symmetric groups are denoted by $`S_m`$; $`S_m`$ acts freely on $`C^m(K)`$, and on its $`i`$-skeleton $`(C^m(K))^i`$ for each $`i`$, by exchanging the coordinates. A necessary condition for the existence of an embedding $`K^n^{2n}`$ is the existence, for each $`m`$, of an $`S_m`$-equivariant map $`C^m(K)C^m(^{2n})`$. We will now analyze the first embedding obstruction, corresponding to $`m=2`$, that is, the obstruction to existence of a $`/2`$-equivariant map $$K\times K\mathrm{\Delta }^{2n}\times ^{2n}\mathrm{\Delta }S^{2n1}.$$ The $`/2`$-equivariant homotopy equivalence above is given by the projection of $`^{2n}\times ^{2n}\mathrm{\Delta }`$ onto the unit sphere in the antidiagonal $`\{(x,x)\}^{2n}\times ^{2n}`$. The diagonal $`\mathrm{\Delta }`$ in $`K\times K`$ is the “simplicial” diagonal, as defined in section 2, while $`\mathrm{\Delta }^{2n}\times ^{2n}`$ is the usual set-theoretic diagonal. Recall that the spaces above are denoted in short by $`K^{}`$ and $`(^{2n})^{}`$ respectively. ###### Theorem 5.1. The obstruction to existence of a $`_2`$-equivariant map $`K^{}(^{2n})^{}`$ lies in $`H_{/2}^{2n}(K^{};)`$ and coincides with Van Kampen’s obstruction $`o(K)`$. Proof. Since $`K^{}`$ is a $`(2n)`$-dimensional CW-complex, the only non-trivial obstruction group in this setting is $`H_{/2}^{2n}(K^{};`$ $`\pi _{2n1}(S^{2n1}))H^{2n}_{/2}(K^{};)`$. Let $`f:K^{2n}`$ be any immersion. Since $`f(\sigma )`$ and $`f(\nu )`$ are disjoint for any $`n`$-simplex $`\sigma `$ and any $`(n1)`$-simplex $`\nu `$, $`f\times f`$ restricted to the $`(2n1)`$-skeleton of $`K^{}`$ is a $`/2`$-equivariant embedding into $`(^{2n})^{}`$. Let $`\sigma `$, $`\tau `$ be two $`n`$-dimensional simplices of $`K`$ and consider $`\sigma \times \tau `$ as an oriented generator of $`(2n)`$-dimensional cellular chains on $`K^{}`$. The obstruction cochain $`c_f`$ assigns to $`\sigma \times \tau `$ the element $$c_f(\sigma \times \tau )=[(f\times f)((\sigma \times \tau ))]\pi _{2n1}(S^{2n1}).$$ The map $`f\times f`$ sends $`\sigma \times \tau `$ into $`^{2n}\times ^{2n}`$, and one has $$o_f(\sigma \times \tau )=f(\sigma )f(\tau )=(f\times f)(\sigma \times \tau )\mathrm{\Delta }_{^{2n}}=(f\times f)((\sigma \times \tau ))=c_f(\sigma \times \tau ).$$ This shows that the homotopy-theoretic obstruction coincides with Van Kampen’s obstruction even on the cochain level, when the map of $`(2n1)`$-skeleton of $`K^{}`$ corresponds to the chosen immersion $`f`$. This completes the proof, since the cohomology class $`[c_f]`$ is independent of the choice of a map of $`(2n1)`$-skeleton of $`K^{}`$, being the first non-trivial obstruction. ∎ ###### Remark 5.2. This result is implicitely contained in , , . It is interesting to note that by theorems 2.1 and 5.1, the existence of a $`/2`$ \- equivariant map $`K^{}(^{2n})^{}`$ is equivalent to existence of an embedding $`K^n^{2n}`$ for $`n2`$. We conclude by suggesting the following approach to defining higher embedding obstructions, which will be pursued in a separate paper. Suppose $`o(K)`$ vanishes, so there exists a $`/2`$-equivariant map $`K^{}(^4)^{}`$. One may consider the obstructions to existence of an equivariant, with respect to the free action of the symmetric group $`S_m`$, map $`C^m(K)C^m(^4)`$, $`m=3,4,\mathrm{}`$. Fadell and Neuwirth have determined homotopy types of the symmetric products of Euclidean spaces, thus (rationally) explicitely giving the coefficients of obstruction groups. Note that the examples in (similar to those constructed in ) show that the entire sequence of such obstructions, arising from configuration spaces, is incomplete. Acknowledgements. I would like to thank Michael Freedman, Richard Stong and Peter Teichner for many discussions. This work has been done during my stays at the University of California - San Diego, Michigan State University and the Max-Planck-Institut für Mathematik, and I would like to thank them for their hospitality and support.
warning/0004/hep-ph0004188.html	ar5iv	text	# ASTROPHYSICS AND COSMOLOGY ## 1 GENERAL INTRODUCTION Cosmology (from the Greek: kosmos, universe, world, order, and logos, word, theory) is probably the most ancient body of knowledge, dating from as far back as the predictions of seasons by early civilizations. Yet, until recently, we could only answer to some of its more basic questions with an order of magnitude estimate. This poor state of affairs has dramatically changed in the last few years, thanks to (what else?) raw data, coming from precise measurements of a wide range of cosmological parameters. Furthermore, we are entering a precision era in cosmology, and soon most of our observables will be measured with a few percent accuracy. We are truly living in the Golden Age of Cosmology. It is a very exciting time and I will try to communicate this enthusiasm to you. Important results are coming out almost every month from a large set of experiments, which provide crucial information about the universe origin and evolution; so rapidly that these notes will probably be outdated before they are in print as a CERN report. In fact, some of the results I mentioned during the Summer School have already been improved, specially in the area of the microwave background anisotropies. Nevertheless, most of the new data can be interpreted within a coherent framework known as the standard cosmological model, based on the Big Bang theory of the universe and the inflationary paradigm, which is with us for two decades. I will try to make such a theoretical model accesible to young experimental particle physicists with little or no previous knowledge about general relativity and curved space-time, but with some knowledge of quantum field theory and the standard model of particle physics. ## 2 INTRODUCTION TO BIG BANG COSMOLOGY Our present understanding of the universe is based upon the successful hot Big Bang theory, which explains its evolution from the first fraction of a second to our present age, around 13 billion years later. This theory rests upon four strong pillars, a theoretical framework based on general relativity, as put forward by Albert Einstein and Alexander A. Friedmann in the 1920s, and three robust observational facts: First, the expansion of the universe, discovered by Edwin P. Hubble in the 1930s, as a recession of galaxies at a speed proportional to their distance from us. Second, the relative abundance of light elements, explained by George Gamow in the 1940s, mainly that of helium, deuterium and lithium, which were cooked from the nuclear reactions that took place at around a second to a few minutes after the Big Bang, when the universe was a few times hotter than the core of the sun. Third, the cosmic microwave background (CMB), the afterglow of the Big Bang, discovered in 1965 by Arno A. Penzias and Robert W. Wilson as a very isotropic blackbody radiation at a temperature of about 3 degrees Kelvin, emitted when the universe was cold enough to form neutral atoms, and photons decoupled from matter, approximately 500,000 years after the Big Bang. Today, these observations are confirmed to within a few percent accuracy, and have helped establish the hot Big Bang as the preferred model of the universe. ### 2.1 Friedmann–Robertson–Walker universes Where are we in the universe? During our lectures, of course, we were in Časta Papiernička, in “the heart of Europe”, on planet Earth, rotating (8 light-minutes away) around the Sun, an ordinary star 8.5 kpc<sup>1</sup><sup>1</sup>1One parallax second (1 pc), parsec for short, corresponds to a distance of about 3.26 light-years or $`3\times 10^{18}`$ cm. from the center of our galaxy, the Milky Way, which is part of the local group, within the Virgo cluster of galaxies (of size a few Mpc), itself part of a supercluster (of size $`100`$ Mpc), within the visible universe ($`\mathrm{few}\times 1000`$ Mpc), most probably a tiny homogeneous patch of the infinite global structure of space-time, much beyond our observable universe. Cosmology studies the universe as we see it. Due to our inherent inability to experiment with it, its origin and evolution has always been prone to wild speculation. However, cosmology was born as a science with the advent of general relativity and the realization that the geometry of space-time, and thus the general attraction of matter, is determined by the energy content of the universe , $$G_{\mu \nu }R_{\mu \nu }\frac{1}{2}g_{\mu \nu }R=8\pi GT_{\mu \nu }+\mathrm{\Lambda }g_{\mu \nu }.$$ (1) These non-linear equations are simply too difficult to solve without some insight coming from the symmetries of the problem at hand: the universe itself. At the time (1917-1922) the known (observed) universe extended a few hundreds of parsecs away, to the galaxies in the local group, Andromeda and the Large and Small Magellanic Clouds: The universe looked extremely anisotropic. Nevertheless, both Einstein and Friedmann speculated that the most “reasonable” symmetry for the universe at large should be homogeneity at all points, and thus isotropy. It was not until the detection, a few decades later, of the microwave background by Penzias and Wilson that this important assumption was finally put onto firm experimental ground. So, what is the most general metric satisfying homogeneity and isotropy at large scales? The Friedmann-Robertson-Walker (FRW) metric, written here in terms of the invariant geodesic distance $`ds^2=g_{\mu \nu }dx^\mu dx^\nu `$ in four dimensions, $`\mu =0,1,2,3`$, see Ref. ,<sup>2</sup><sup>2</sup>2I am using $`c=1`$ everywhere, unless specified. $$ds^2=dt^2a^2(t)\left[\frac{dr^2}{1Kr^2}+r^2(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2)\right],$$ (2) characterized by just two quantities, a scale factor $`a(t)`$, which determines the physical size of the universe, and a constant $`K`$, which characterizes the spatial curvature of the universe, $${}_{}{}^{(3)}R=\frac{6K}{a^2(t)}.\{\begin{array}{cc}K=1\hfill & \hfill \mathrm{OPEN}\\ K=0\hfill & \hfill \mathrm{FLAT}\\ K=+1\hfill & \hfill \mathrm{CLOSED}\end{array}$$ (3) Spatially open, flat and closed universes have different geometries. Light geodesics on these universes behave differently, and thus could in principle be distinguished observationally, as we shall discuss later. Apart from the three-dimensional spatial curvature, we can also compute a four-dimensional space-time curvature, $${}_{}{}^{(4)}R=6\frac{\ddot{a}}{a}+6\left(\frac{\dot{a}}{a}\right)^2+6\frac{K}{a^2}.$$ (4) Depending on the dynamics (and thus on the matter/energy content) of the universe, we will have different possible outcomes of its evolution. The universe may expand for ever, recollapse in the future or approach an asymptotic state in between. #### 2.1.1 The expansion of the universe In 1929, Edwin P. Hubble observed a redshift in the spectra of distant galaxies, which indicated that they were receding from us at a velocity proportional to their distance to us . This was correctly interpreted as mainly due to the expansion of the universe, that is, to the fact that the scale factor today is larger than when the photons were emitted by the observed galaxies. For simplicity, consider the metric of a spatially flat universe, $`ds^2=dt^2a^2(t)d\stackrel{}{x}^2`$ (the generalization of the following argument to curved space is straightforward). The scale factor $`a(t)`$ gives physical size to the spatial coordinates $`\stackrel{}{x}`$, and the expansion is nothing but a change of scale (of spatial units) with time. Except for peculiar velocities, i.e. motion due to the local attraction of matter, galaxies do not move in coordinate space, it is the space-time fabric which is stretching between galaxies. Due to this continuous stretching, the observed wavelength of photons coming from distant objects is greater than when they were emitted by a factor precisely equal to the ratio of scale factors, $$\frac{\lambda _{\mathrm{obs}}}{\lambda _{\mathrm{em}}}=\frac{a_0}{a}1+z,$$ (5) where $`a_0`$ is the present value of the scale factor. Since the universe today is larger than in the past, the observed wavelengths will be shifted towards the red, or redshifted, by an amount characterized by $`z`$, the redshift parameter. In the context of a FRW metric, the universe expansion is characterized by a quantity known as the Hubble rate of expansion, $`H(t)=\dot{a}(t)/a(t)`$, whose value today is denoted by $`H_0`$. As I shall deduce later, it is possible to compute the relation between the physical distance $`d_L`$ and the present rate of expansion, in terms of the redshift parameter,<sup>3</sup><sup>3</sup>3The subscript $`L`$ refers to Luminosity, which characterizes the amount of light emitted by an object. See Eq. (69). $$H_0d_L=z+\frac{1}{2}(1q_0)z^2+𝒪(z^3).$$ (6) At small distances from us, i.e. at $`z1`$, we can safely keep only the linear term, and thus the recession velocity becomes proportional to the distance from us, $`v=cz=H_0d_L`$, the proportionality constant being the Hubble rate, $`H_0`$. This expression constitutes the so-called Hubble law, and is spectacularly confirmed by a huge range of data, up to distances of hundreds of megaparsecs. In fact, only recently measurements from very bright and distant supernovae, at $`z1`$, were obtained, and are beginning to probe the second-order term, proportional to the deceleration parameter $`q_0`$, see Eq. (26). I will come back to these measurements in Section 3. One may be puzzled as to why do we see such a stretching of space-time. Indeed, if all spatial distances are scaled with a universal scale factor, our local measuring units (our rulers) should also be stretched, and therefore we should not see the difference when comparing the two distances (e.g. the two wavelengths) at different times. The reason we see the difference is because we live in a gravitationally bound system, decoupled from the expansion of the universe: local spatial units in these systems are not stretched by the expansion.<sup>4</sup><sup>4</sup>4The local space-time of a gravitationally bound system is described by the Schwarzschild metric, which is static . The wavelengths of photons are stretched along their geodesic path from one galaxy to another. In this consistent world picture, galaxies are like point particles, moving as a fluid in an expanding universe. #### 2.1.2 The matter and energy content of the universe So far I have only discussed the geometrical aspects of space-time. Let us now consider the matter and energy content of such a universe. The most general matter fluid consistent with the assumption of homogeneity and isotropy is a perfect fluid, one in which an observer comoving with the fluid would see the universe around it as isotropic. The energy momentum tensor associated with such a fluid can be written as $$T^{\mu \nu }=pg^{\mu \nu }+(p+\rho )U^\mu U^\nu ,$$ (7) where $`p(t)`$ and $`\rho (t)`$ are the pressure and energy density of the fluid at a given time in the expansion, and $`U^\mu `$ is the comoving four-velocity, satisfying $`U^\mu U_\mu =1`$. Let us now write the equations of motion of such a fluid in an expanding universe. According to general relativity, these equations can be deduced from the Einstein equations (1), where we substitute the FRW metric (2) and the perfect fluid tensor (7). The $`\mu =\nu =0`$ component of the Einstein equations constitutes the so-called Friedmann equation $$H^2=\left(\frac{\dot{a}}{a}\right)^2=\frac{8\pi G}{3}\rho +\frac{\mathrm{\Lambda }}{3}\frac{K}{a^2},$$ (8) where I have treated the cosmological constant $`\mathrm{\Lambda }`$ as a different component from matter. In fact, it can be associated with the vacuum energy of quantum field theory, although we still do not understand why should it have such a small value (120 orders of magnitude below that predicted by quantum theory), if it is non-zero. This constitutes today one of the most fundamental problems of physics, let alone cosmology. The conservation of energy ($`T_{;\nu }^{\mu \nu }=0`$), a direct consequence of the general covariance of the theory ($`G_{;\nu }^{\mu \nu }=0`$), can be written in terms of the FRW metric and the perfect fluid tensor (7) as $$\frac{d}{dt}\left(\rho a^3\right)+p\frac{d}{dt}\left(a^3\right)=0,$$ (9) where the energy density and pressure can be split into its matter and radiation components, $`\rho =\rho _\mathrm{M}+\rho _\mathrm{R},p=p_\mathrm{M}+p_\mathrm{R}`$, with corresponding equations of state, $`p_\mathrm{M}=0,p_\mathrm{R}=\rho _\mathrm{R}/3`$. Together, the Friedmann and the energy-conservation equation give the evolution equation for the scale factor, $$\frac{\ddot{a}}{a}=\frac{4\pi G}{3}(\rho +3p)+\frac{\mathrm{\Lambda }}{3},$$ (10) I will now make a few useful definitions. We can write the Hubble parameter today $`H_0`$ in units of 100 km s<sup>-1</sup>Mpc<sup>-1</sup>, in terms of which one can estimate the order of magnitude for the present size and age of the universe, $`H_0`$ $`=`$ $`100h\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1,`$ (11) $`cH_0^1`$ $`=`$ $`3000h^1\mathrm{Mpc},`$ (12) $`H_0^1`$ $`=`$ $`9.773h^1\mathrm{Gyr}.`$ (13) The parameter $`h`$ has been measured to be in the range $`0.4<h<1`$ for decades, and only in the last few years has it been found to lie within 10% of $`h=0.65`$. I will discuss those recent measurements in the next Section. One can also define a critical density $`\rho _c`$, that which in the absence of a cosmological constant would correspond to a flat universe, $`\rho _c{\displaystyle \frac{3H_0^2}{8\pi G}}`$ $`=`$ $`1.88h^2\mathrm{\hspace{0.17em}10}^{29}\mathrm{g}/\mathrm{cm}^3`$ (14) $`=`$ $`2.77h^1\mathrm{\hspace{0.17em}10}^{11}M_{}/(h^1\mathrm{Mpc})^3,`$ (15) where $`M_{}=1.989\times 10^{33}`$ g is a solar mass unit. The critical density $`\rho _c`$ corresponds to approximately 4 protons per cubic meter, certainly a very dilute fluid! In terms of the critical density it is possible to define the ratios $`\mathrm{\Omega }_i\rho _i/\rho _c`$, for matter, radiation, cosmological constant and even curvature, today, $`\mathrm{\Omega }_\mathrm{M}={\displaystyle \frac{8\pi G\rho _\mathrm{M}}{3H_0^2}}\mathrm{\Omega }_\mathrm{R}={\displaystyle \frac{8\pi G\rho _\mathrm{R}}{3H_0^2}}`$ (16) $`\mathrm{\Omega }_\mathrm{\Lambda }={\displaystyle \frac{\mathrm{\Lambda }}{3H_0^2}}\mathrm{\Omega }_K={\displaystyle \frac{K}{a_0^2H_0^2}}.`$ (17) We can evaluate today the radiation component $`\mathrm{\Omega }_\mathrm{R}`$, corresponding to relativistic particles, from the density of microwave background photons, $`\rho _{_{\mathrm{CMB}}}=\frac{\pi ^2}{15}(kT_{_{\mathrm{CMB}}})^4/(\mathrm{}c)^3=4.5\times 10^{34}\mathrm{g}/\mathrm{cm}^3`$, which gives $`\mathrm{\Omega }_{_{\mathrm{CMB}}}=2.4\times 10^5h^2`$. Three massless neutrinos contribute an even smaller amount. Therefore, we can safely neglect the contribution of relativistic particles to the total density of the universe today, which is dominated either by non-relativistic particles (baryons, dark matter or massive neutrinos) or by a cosmological constant, and write the rate of expansion $`H^2`$ in terms of its value today, $$H^2(a)=H_0^2\left(\mathrm{\Omega }_\mathrm{R}\frac{a_0^4}{a^4}+\mathrm{\Omega }_\mathrm{M}\frac{a_0^3}{a^3}+\mathrm{\Omega }_\mathrm{\Lambda }+\mathrm{\Omega }_K\frac{a_0^2}{a^2}\right).$$ (18) An interesting consequence of these redefinitions is that I can now write the Friedmann equation today, $`a=a_0`$, as a cosmic sum rule, $$1=\mathrm{\Omega }_\mathrm{M}+\mathrm{\Omega }_\mathrm{\Lambda }+\mathrm{\Omega }_K,$$ (19) where we have neglected $`\mathrm{\Omega }_\mathrm{R}`$ today. That is, in the context of a FRW universe, the total fraction of matter density, cosmological constant and spatial curvature today must add up to one. For instance, if we measure one of the three components, say the spatial curvature, we can deduce the sum of the other two. Making use of the cosmic sum rule today, we can write the matter and cosmological constant as a function of the scale factor ($`a_01`$) $`\mathrm{\Omega }_\mathrm{M}(a)={\displaystyle \frac{8\pi G\rho _\mathrm{M}}{3H^2(a)}}={\displaystyle \frac{\mathrm{\Omega }_\mathrm{M}}{a+\mathrm{\Omega }_\mathrm{M}(1a)+\mathrm{\Omega }_\mathrm{\Lambda }(a^3a)}}\{\begin{array}{cc}\stackrel{a0}{}\hfill & \hfill 1\\ \stackrel{a\mathrm{}}{}\hfill & \hfill 0\end{array},`$ (22) $`\mathrm{\Omega }_\mathrm{\Lambda }(a)={\displaystyle \frac{\mathrm{\Lambda }}{3H^2(a)}}={\displaystyle \frac{\mathrm{\Omega }_\mathrm{\Lambda }a^3}{a+\mathrm{\Omega }_\mathrm{M}(1a)+\mathrm{\Omega }_\mathrm{\Lambda }(a^3a)}}\{\begin{array}{cc}\stackrel{a0}{}\hfill & \hfill 0\\ \stackrel{a\mathrm{}}{}\hfill & \hfill 1\end{array}.`$ (25) This implies that for sufficiently early times, $`a1`$, all matter-dominated FRW universes can be described by Einstein-de Sitter (EdS) models ($`\mathrm{\Omega }_K=0,\mathrm{\Omega }_\mathrm{\Lambda }=0`$).<sup>5</sup><sup>5</sup>5Note that in the limit $`a0`$ the radiation component starts dominating, see Eq. (18), but we still recover the EdS model. On the other hand, the vacuum energy will always dominate in the future. Another relationship which becomes very useful is that of the cosmological deceleration parameter today, $`q_0`$, in terms of the matter and cosmological constant components of the universe, see Eq. (10), $$q_0\frac{\ddot{a}}{aH^2}\|_0=\frac{1}{2}\mathrm{\Omega }_\mathrm{M}\mathrm{\Omega }_\mathrm{\Lambda },$$ (26) which is independent of the spatial curvature. Uniform expansion corresponds to $`q_0=0`$ and requires a precise cancellation: $`\mathrm{\Omega }_\mathrm{M}=2\mathrm{\Omega }_\mathrm{\Lambda }`$. It represents spatial sections that are expanding at a fixed rate, its scale factor growing by the same amount in equally-spaced time intervals. Accelerated expansion corresponds to $`q_0<0`$ and comes about whenever $`\mathrm{\Omega }_\mathrm{M}<2\mathrm{\Omega }_\mathrm{\Lambda }`$: spatial sections expand at an increasing rate, their scale factor growing at a greater speed with each time interval. Decelerated expansion corresponds to $`q_0>0`$ and occurs whenever $`\mathrm{\Omega }_\mathrm{M}>2\mathrm{\Omega }_\mathrm{\Lambda }`$: spatial sections expand at a decreasing rate, their scale factor growing at a smaller speed with each time interval. #### 2.1.3 Mechanical analogy It is enlightening to work with a mechanical analogy of the Friedmann equation. Let us rewrite Eq. (8) as $$\frac{1}{2}\dot{a}^2\frac{GM}{a}\frac{\mathrm{\Lambda }}{6}a^2=\frac{K}{2}=\mathrm{constant},$$ (27) where $`M\frac{4\pi }{3}\rho a^3`$ is the equivalent of mass for the whole volume of the universe. Equation (27) can be understood as the energy conservation law $`E=T+V`$ for a test particle of unit mass in the central potential $$V(r)=\frac{GM}{r}+\frac{1}{2}kr^2,$$ (28) corresponding to a Newtonian potential plus a harmonic oscillator potential with a negative spring constant $`k\mathrm{\Lambda }/3`$. Note that, in the absence of a cosmological constant ($`\mathrm{\Lambda }=0`$), a critical universe, defined as the borderline between indefinite expansion and recollapse, corresponds, through the Friedmann equations of motion, precisely with a flat universe ($`K=0`$). In that case, and only in that case, a spatially open universe ($`K=1`$) corresponds to an eternally expanding universe, and a spatially closed universe ($`K=+1`$) to a recollapsing universe in the future. Such a well known (textbook) correspondence is incorrect when $`\mathrm{\Omega }_\mathrm{\Lambda }0`$: spatially open universes may recollapse while closed universes can expand forever. One can see in Fig. 1 a range of possible evolutions of the scale factor, for various pairs of values of $`(\mathrm{\Omega }_\mathrm{M},\mathrm{\Omega }_\mathrm{\Lambda })`$. One can show that, for $`\mathrm{\Omega }_\mathrm{\Lambda }0`$, a critical universe ($`H=\dot{H}=0`$) corresponds to those points $`xa_0/a>0`$, for which $`f(x)H^2(a)`$ and $`f^{}(x)`$ vanish, while $`f^{\prime \prime }(x)>0`$, $`f(x)=x^3\mathrm{\Omega }_\mathrm{M}+x^2\mathrm{\Omega }_K+\mathrm{\Omega }_\mathrm{\Lambda }=0,`$ (29) $`f^{}(x)=3x^2\mathrm{\Omega }_\mathrm{M}+2x\mathrm{\Omega }_K=0\{\begin{array}{c}x=0\hfill \\ x=2\mathrm{\Omega }_K/3\mathrm{\Omega }_\mathrm{M}>0\hfill \end{array},`$ (32) $`f^{\prime \prime }(x)=6x\mathrm{\Omega }_\mathrm{M}+2\mathrm{\Omega }_K=\{\begin{array}{cc}+2\mathrm{\Omega }_K>0& x=0\hfill \\ 2\mathrm{\Omega }_K>0& x=2\|\mathrm{\Omega }_K\|/3\mathrm{\Omega }_\mathrm{M}\hfill \end{array}.`$ (35) Using the cosmic sum rule (19), we can write the solutions as $$\mathrm{\Omega }_\mathrm{\Lambda }=\{\begin{array}{cc}0\hfill & \mathrm{\Omega }_\mathrm{M}1\hfill \\ 4\mathrm{\Omega }_\mathrm{M}\mathrm{sin}^3\left[\frac{1}{3}\mathrm{arcsin}(1\mathrm{\Omega }_\mathrm{M}^1)\right]\hfill & \mathrm{\Omega }_\mathrm{M}1\hfill \end{array}.$$ (36) The first solution corresponds to the critical point $`x=0`$ ($`a=\mathrm{}`$), and $`\mathrm{\Omega }_K>0`$, while the second one to $`x=2\|\mathrm{\Omega }_K\|/3\mathrm{\Omega }_\mathrm{M}`$, and $`\mathrm{\Omega }_K<0`$. Expanding around $`\mathrm{\Omega }_\mathrm{M}=1`$, we find $`\mathrm{\Omega }_\mathrm{\Lambda }\frac{4}{27}(\mathrm{\Omega }_\mathrm{M}1)^3/\mathrm{\Omega }_\mathrm{M}^2`$, for $`\mathrm{\Omega }_\mathrm{M}1`$. These critical solutions are asymptotic to the Einstein-de Sitter model ($`\mathrm{\Omega }_\mathrm{M}=1,\mathrm{\Omega }_\mathrm{\Lambda }=0`$), see Fig. 2. #### 2.1.4 Thermodynamical analogy It is also enlightening to find an analogy between the energy conservation equation (9) and the second law of Thermodynamics, $$TdS=dU+pdV,$$ (37) where $`U=\rho V`$ is the total energy of the closed system and $`V=a^3`$ is its physical volume. Equation (9) implies that the expansion of the universe is adiabatic or isoentropic ($`dS=0`$), corresponding to a fluid in thermal equilibrium at a temperature T. For a barotropic fluid, satisfying the equation of state $`p=\omega \rho `$, we can write the energy density evolution as $$\frac{d}{dt}(\rho a^3)=p\frac{d}{dt}(a^3)=3H\omega (\rho a^3).$$ (38) For relativistic particles in thermal equilibrium, the trace of the energy-momentum tensor vanishes (because of conformal invariance) and thus $`p_\mathrm{R}=\rho _\mathrm{R}/3\omega =1/3`$. In that case, the energy density of radiation in thermal equilibrium can be written as $`\rho _\mathrm{R}`$ $`=`$ $`{\displaystyle \frac{\pi ^2}{30}}g_{}T^4,`$ (39) $`g_{}`$ $`=`$ $`{\displaystyle \underset{i=\mathrm{bosons}}{}}g_i\left({\displaystyle \frac{T_i}{T}}\right)^4+{\displaystyle \frac{7}{8}}{\displaystyle \underset{i=\mathrm{fermions}}{}}g_i\left({\displaystyle \frac{T_i}{T}}\right)^4,`$ (40) where $`g_{}`$ is the number of relativistic degrees of freedom, coming from both bosons and fermions. Using the equilibrium expressions for the pressure and density, we can write $`dp=(\rho +p)dT/T`$, and therefore $$dS=\frac{1}{T}d[(\rho +p)V](\rho +p)V\frac{dT}{T^2}=d[\frac{(\rho +p)V}{T}+\mathrm{const}.]$$ (41) That is, up to an additive constant, the entropy per comoving volume is $`S=a^3(\rho +p)V/T`$, which is conserved. The entropy per comoving volume is dominated by the contribution of relativistic particles, so that, to very good approximation, $`S`$ $`=`$ $`{\displaystyle \frac{2\pi ^2}{45}}g_s(aT)^3=\mathrm{constant},`$ (42) $`g_s`$ $`=`$ $`{\displaystyle \underset{i=\mathrm{bosons}}{}}g_i\left({\displaystyle \frac{T_i}{T}}\right)^3+{\displaystyle \frac{7}{8}}{\displaystyle \underset{i=\mathrm{fermions}}{}}g_i\left({\displaystyle \frac{T_i}{T}}\right)^3.`$ (43) A consequence of Eq. (42) is that, during the adiabatic expansion of the universe, the scale factor grows inversely proportional to the temperature of the universe, $`a1/T`$. Therefore, the observational fact that the universe is expanding today implies that in the past the universe must have been much hotter and denser, and that in the future it will become much colder and dilute. Since the ratio of scale factors can be described in terms of the redshift parameter $`z`$, see Eq. (5), we can find the temperature of the universe at an earlier epoch by $$T=T_0(1+z).$$ (44) Such a relation has been spectacularly confirmed with observations of absorption spectra from quasars at large distances, which showed that, indeed, the temperature of the radiation background scaled with redshift in the way predicted by the hot Big Bang model. ### 2.2 Brief thermal history of the universe In this Section, I will briefly summarize the thermal history of the universe, from the Planck era to the present. As we go back in time, the universe becomes hotter and hotter and thus the amount of energy available for particle interactions increases. As a consequence, the nature of interactions goes from those described at low energy by long range gravitational and electromagnetic physics, to atomic physics, nuclear physics, all the way to high energy physics at the electroweak scale, gran unification (perhaps), and finally quantum gravity. The last two are still uncertain since we do not have any experimental evidence for those ultra high energy phenomena, and perhaps Nature has followed a different path. <sup>6</sup><sup>6</sup>6See the recent theoretical developments on large extra dimensions and quantum gravity at the TeV . The way we know about the high energy interactions of matter is via particle accelerators, which are unravelling the details of those fundamental interactions as we increase in energy. However, one should bear in mind that the physical conditions that take place in our high energy colliders are very different from those that occurred in the early universe. These machines could never reproduce the conditions of density and pressure in the rapidly expanding thermal plasma of the early universe. Nevertheless, those experiments are crucial in understanding the nature and rate of the local fundamental interactions available at those energies. What interests cosmologists is the statistical and thermal properties that such a plasma should have, and the role that causal horizons play in the final outcome of the early universe expansion. For instance, of crucial importance is the time at which certain particles decoupled from the plasma, i.e. when their interactions were not quick enough compared with the expansion of the universe, and they were left out of equilibrium with the plasma. One can trace the evolution of the universe from its origin till today. There is still some speculation about the physics that took place in the universe above the energy scales probed by present colliders. Nevertheless, the overall layout presented here is a plausible and hopefully testable proposal. According to the best accepted view, the universe must have originated at the Planck era ($`10^{19}`$ GeV, $`10^{43}`$ s) from a quantum gravity fluctuation. Needless to say, we don’t have any experimental evidence for such a statement: Quantum gravity phenomena are still in the realm of physical speculation. However, it is plausible that a primordial era of cosmological inflation originated then. Its consequences will be discussed below. Soon after, the universe may have reached the Grand Unified Theories (GUT) era ($`10^{16}`$ GeV, $`10^{35}`$ s). Quantum fluctuations of the inflaton field most probably left their imprint then as tiny perturbations in an otherwise very homogenous patch of the universe. At the end of inflation, the huge energy density of the inflaton field was converted into particles, which soon thermalized and became the origin of the hot Big Bang as we know it. Such a process is called reheating of the universe. Since then, the universe became radiation dominated. It is probable (although by no means certain) that the asymmetry between matter and antimatter originated at the same time as the rest of the energy of the universe, from the decay of the inflaton. This process is known under the name of baryogenesis since baryons (mostly quarks at that time) must have originated then, from the leftovers of their annihilation with antibaryons. It is a matter of speculation whether baryogenesis could have occurred at energies as low as the electroweak scale ($`100`$ GeV, $`10^{10}`$ s). Note that although particle physics experiments have reached energies as high as 100 GeV, we still do not have observational evidence that the universe actually went through the EW phase transition. If confirmed, baryogenesis would constitute another “window” into the early universe. As the universe cooled down, it may have gone through the quark-gluon phase transition ($`10^2`$ MeV, $`10^5`$ s), when baryons (mainly protons and neutrons) formed from their constituent quarks. The furthest window we have on the early universe at the moment is that of primordial nucleosynthesis ($`10.1`$ MeV, 1 s – 3 min), when protons and neutrons were cold enough that bound systems could form, giving rise to the lightest elements, soon after neutrino decoupling: It is the realm of nuclear physics. The observed relative abundances of light elements are in agreement with the predictions of the hot Big Bang theory. Immediately afterwards, electron-positron annihilation occurs (0.5 MeV, 1 min) and all their energy goes into photons. Much later, at about (1 eV, $`10^5`$ yr), matter and radiation have equal energy densities. Soon after, electrons become bound to nuclei to form atoms (0.3 eV, $`3\times 10^5`$ yr), in a process known as recombination: It is the realm of atomic physics. Immediately after, photons decouple from the plasma, travelling freely since then. Those are the photons we observe as the cosmic microwave background. Much later ($`110`$ Gyr), the small inhomogeneities generated during inflation have grown, via gravitational collapse, to become galaxies, clusters of galaxies, and superclusters, characterizing the epoch of structure formation. It is the realm of long range gravitational physics, perhaps dominated by a vacuum energy in the form of a cosmological constant. Finally (3K, 13 Gyr), the Sun, the Earth, and biological life originated from previous generations of stars, and from a primordial soup of organic compounds, respectively. I will now review some of the more robust features of the Hot Big Bang theory of which we have precise observational evidence. #### 2.2.1 Primordial nucleosynthesis and light element abundance In this subsection I will briefly review Big Bang nucleosynthesis and give the present observational constraints on the amount of baryons in the universe. In 1920 Eddington suggested that the sun might derive its energy from the fusion of hydrogen into helium. The detailed reactions by which stars burn hydrogen were first laid out by Hans Bethe in 1939. Soon afterwards, in 1946, George Gamow realized that similar processes might have occurred also in the hot and dense early universe and gave rise to the first light elements . These processes could take place when the universe had a temperature of around $`T_{_{\mathrm{NS}}}10.1`$ MeV, which is about 100 times the temperature in the core of the Sun, while the density is $`\rho _{_{\mathrm{NS}}}=\frac{\pi ^2}{30}g_{}T_{_{\mathrm{NS}}}^482`$ g cm<sup>-3</sup>, about the same density as the core of the Sun. Note, however, that although both processes are driven by identical thermonuclear reactions, the physical conditions in star and Big Bang nucleosynthesis are very different. In the former, gravitational collapse heats up the core of the star and reactions last for billions of years (except in supernova explosions, which last a few minutes and creates all the heavier elements beyond iron), while in the latter the universe expansion cools the hot and dense plasma in just a few minutes. Nevertheless, Gamow reasoned that, although the early period of cosmic expansion was much shorter than the lifetime of a star, there was a large number of free neutrons at that time, so that the lighter elements could be built up quickly by succesive neutron captures, starting with the reaction $`n+pD+\gamma `$. The abundances of the light elements would then be correlated with their neutron capture cross sections, in rough agreement with observations . Nowadays, Big Bang nucleosynthesis (BBN) codes compute a chain of around 30 coupled nuclear reactions, to produce all the light elements up to beryllium-7. <sup>7</sup><sup>7</sup>7The rest of nuclei, up to iron (Fe), are produced in heavy stars, and beyond Fe in novae and supernovae explosions. Only the first four or five elements can be computed with accuracy better than 1% and compared with cosmological observations. These light elements are $`H,{}_{}{}^{4}He,D,{}_{}{}^{3}He,{}_{}{}^{7}Li`$, and perhaps also $`{}_{}{}^{6}Li`$. Their observed relative abundance to hydrogen is $`[1:0.25:310^5:210^5:210^{10}]`$ with various errors, mainly systematic. The BBN codes calculate these abundances using the laboratory measured nuclear reaction rates, the decay rate of the neutron, the number of light neutrinos and the homogeneous FRW expansion of the universe, as a function of only one variable, the number density fraction of baryons to photons, $`\eta n_\mathrm{B}/n_\gamma `$. In fact, the present observations are only consistent, see Fig. 3 and Ref. , with a very narrow range of values of $$\eta _{10}10^{10}\eta =4.65.9.$$ (45) Such a small value of $`\eta `$ indicates that there is about one baryon per $`10^9`$ photons in the universe today. Any acceptable theory of baryogenesis should account for such a small number. Furthermore, the present baryon fraction of the critical density can be calculated from $`\eta _{10}`$ as $$\mathrm{\Omega }_\mathrm{B}h^2=3.6271\times 10^3\eta _{10}=0.0190\pm 0.0024(95\%\mathrm{c}.\mathrm{l}.)$$ (46) Clearly, this number is well below closure density, so baryons cannot account for all the matter in the universe, as I shall discuss below. #### 2.2.2 Neutrino decoupling Just before the nucleosynthesis of the lightest elements in the early universe, weak interactions were too slow to keep neutrinos in thermal equilibrium with the plasma, so they decoupled. We can estimate the temperature at which decoupling occurred from the weak interaction cross section, $`\sigma _\mathrm{w}G_F^2T^2`$ at finite temperature $`T`$, where $`G_F=1.2\times 10^5`$ GeV<sup>-2</sup> is the Fermi constant. The neutrino interaction rate, via W boson exchange in $`n+\nu p+e^{}`$ and $`p+\overline{\nu }n+e^+`$, can be written as $$\mathrm{\Gamma }_\nu =n_\nu \sigma _\mathrm{w}\|v\|G_F^2T^5,$$ (47) while the rate of expansion of the universe at that time ($`g_{}=10.75`$) was $`H5.4T^2/M_\mathrm{P}`$, where $`M_\mathrm{P}=1.22\times 10^{19}`$ GeV is the Planck mass. Neutrinos decouple when their interaction rate is slower than the universe expansion, $`\mathrm{\Gamma }_\nu H`$ or, equivalently, at $`T_{\nu \mathrm{dec}}0.8`$ MeV. Below this temperature, neutrinos are no longer in thermal equilibrium with the rest of the plasma, and their temperature continues to decay inversely proportional to the scale factor of the universe. Since neutrinos decoupled before $`e^+e^{}`$ annihilation, the cosmic background of neutrinos has a temperature today lower than that of the microwave background of photons. Let us compute the difference. At temperatures above the the mass of the electron, $`T>m_e=0.511`$ MeV, and below 0.8 MeV, the only particle species contributing to the entropy of the universe are the photons ($`g_{}=2`$) and the electron-positron pairs ($`g_{}=4\times \frac{7}{8}`$); total number of degrees of freedom $`g_{}=\frac{11}{2}`$. At temperatures $`Tm_e`$, electrons and positrons annihilate into photons, heating up the plasma (but not the neutrinos, which had decoupled already). At temperatures $`T<m_e`$, only photons contribute to the entropy of the universe, with $`g_{}=2`$ degrees of freedom. Therefore, from the conservation of entropy, we find that the ratio of $`T_\gamma `$ and $`T_\nu `$ today must be $$\frac{T_\gamma }{T_\nu }=\left(\frac{11}{4}\right)^{1/3}=1.401T_\nu =1.945\mathrm{K},$$ (48) where I have used $`T_{_{\mathrm{CMB}}}=2.725\pm 0.002`$ K. We still have not measured such a relic background of neutrinos, and probably will remain undetected for a long time, since they have an average energy of order $`10^4`$ eV, much below that required for detection by present experiments (of order GeV), precisely because of the relative weakness of the weak interactions. Nevertheless, it would be fascinating if, in the future, ingenious experiments were devised to detect such a background, since it would confirm one of the most robust features of Big Bang cosmology. #### 2.2.3 Matter-radiation equality Relativistic species have energy densities proportional to the quartic power of temperature and therefore scale as $`\rho _\mathrm{R}a^4`$, while non-relativistic particles have essentially zero pressure and scale as $`\rho _\mathrm{M}a^3`$, see Eq. (38). Therefore, there will be a time in the evolution of the universe in which both energy densities are equal $`\rho _\mathrm{R}(t_{\mathrm{eq}})=\rho _\mathrm{M}(t_{\mathrm{eq}})`$. Since then both decay differently, and thus $$1+z_{\mathrm{eq}}=\frac{a_0}{a_{\mathrm{eq}}}=\frac{\mathrm{\Omega }_\mathrm{M}}{\mathrm{\Omega }_\mathrm{R}}=3.1\times 10^4\mathrm{\Omega }_\mathrm{M}h^2,$$ (49) where I have used $`\mathrm{\Omega }_\mathrm{R}h^2=\mathrm{\Omega }_{_{\mathrm{CMB}}}h^2+\mathrm{\Omega }_\nu h^2=3.24\times 10^5`$ for three massless neutrinos at $`T=T_\nu `$. As I will show later, the matter content of the universe today is below critical, $`\mathrm{\Omega }_\mathrm{M}0.3`$, while $`h0.65`$, and therefore $`(1+z_{\mathrm{eq}})3900`$, or about $`t_{\mathrm{eq}}=1.2\times 10^3(\mathrm{\Omega }_\mathrm{M}h^2)^27\times 10^4`$ years after the origin of the universe. Around the time of matter-radiation equality, the rate of expansion (18) can be written as ($`a_01`$) $$H(a)=H_0\left(\mathrm{\Omega }_\mathrm{R}a^4+\mathrm{\Omega }_\mathrm{M}a^3\right)^{1/2}=H_0\mathrm{\Omega }_\mathrm{M}^{1/2}a^{3/2}\left(1+\frac{a_{\mathrm{eq}}}{a}\right)^{1/2}.$$ (50) The horizon size is the coordinate distance travelled by a photon since the beginning of the universe, $`d_HH^1`$, i.e. the size of causally connected regions in the universe. The comoving horizon size is then given by $$d_H=\frac{c}{aH(a)}=cH_0^1\mathrm{\Omega }_\mathrm{M}^{1/2}a^{1/2}\left(1+\frac{a_{\mathrm{eq}}}{a}\right)^{1/2}.$$ (51) Thus the horizon size at matter-radiation equality ($`a=a_{\mathrm{eq}}`$) is $$d_H(a_{\mathrm{eq}})=\frac{cH_0^1}{\sqrt{2}}\mathrm{\Omega }_\mathrm{M}^{1/2}a_{\mathrm{eq}}^{1/2}12(\mathrm{\Omega }_\mathrm{M}h)^1h^1\mathrm{Mpc}.$$ (52) This scale plays a very important role in theories of structure formation. #### 2.2.4 Recombination and photon decoupling As the temperature of the universe decreased, electrons could eventually become bound to protons to form neutral hydrogen. Nevertheless, there is always a non-zero probability that a rare energetic photon ionizes hydrogen and produces a free electron. The ionization fraction of electrons in equilibrium with the plasma at a given temperature is given by $$\frac{1X_e^{\mathrm{eq}}}{X_e^{\mathrm{eq}}}=\frac{4\sqrt{2}\zeta (3)}{\sqrt{\pi }}\eta \left(\frac{T}{m_e}\right)^{3/2}e^{E_{\mathrm{ion}}/T},$$ (53) where $`E_{\mathrm{ion}}=13.6`$ eV is the ionization energy of hydrogen, and $`\eta `$ is the baryon-to-photon ratio (45). If we now use Eq. (44), we can compute the ionization fraction $`X_e^{\mathrm{eq}}`$ as a function of redshift $`z`$, see Fig. 4. Note that the huge number of photons with respect to electrons (in the ratio $`{}_{}{}^{4}He:H:\gamma 1:4:10^{10}`$) implies that even at a very low temperature, the photon distribution will contain a sufficiently large number of high-energy photons to ionize a significant fraction of hydrogen. In fact, defining recombination as the time at which $`X_e^{\mathrm{eq}}0.1`$, one finds that the recombination temperature is $`T_{\mathrm{rec}}=0.3\mathrm{eV}E_{\mathrm{ion}}`$, for $`\eta _{10}5.2`$. Comparing with the present temperature of the microwave background, we deduce the corresponding redshift at recombination, $`(1+z_{\mathrm{rec}})1270`$. Photons remain in thermal equilibrium with the plasma of baryons and electrons through elastic Thomson scattering, with cross section $$\sigma __T=\frac{8\pi \alpha ^2}{3m_e^2}=6.65\times 10^{25}\mathrm{cm}^2=0.665\mathrm{barn},$$ (54) where $`\alpha =1/137.036`$ is the dimensionless electromagnetic coupling constant. The mean free path of photons $`\lambda _\gamma `$ in such a plasma can be estimated from the photon interaction rate, $`\lambda _\gamma ^1\mathrm{\Gamma }_\gamma =n_e\sigma __T`$. For temperatures above a few eV, the mean free path is much smaller that the causal horizon at that time and photons suffer multiple scattering: the plasma is like a dense fog. Photons will decouple from the plasma when their interaction rate cannot keep up with the expansion of the universe and the mean free path becomes larger than the horizon size: the universe becomes transparent. We can estimate this moment by evaluating $`\mathrm{\Gamma }_\gamma =H`$ at photon decoupling. Using $`n_e=X_e\eta n_\gamma `$, one can compute the decoupling temperature as $`T_{\mathrm{dec}}=0.26`$ eV, and the corresponding redshift as $`(1+z_{\mathrm{dec}})1100`$. This redshift defines the so called last scattering surface, when photons last scattered off protons and electrons and travelled freely ever since. This decoupling occurred when the universe was approximately $`t_{\mathrm{dec}}=1.8\times 10^5(\mathrm{\Omega }_\mathrm{M}h^2)^{1/2}5\times 10^5`$ years old. #### 2.2.5 The microwave background One of the most remarkable observations ever made my mankind is the detection of the relic background of photons from the Big Bang. This background was predicted by George Gamow and collaborators in the 1940s, based on the consistency of primordial nucleosynthesis with the observed helium abundance. They estimated a value of about 10 K, although a somewhat more detailed analysis by Alpher and Herman in 1950 predicted $`T_\gamma 5`$ K. Unfortunately, they had doubts whether the radiation would have survived until the present, and this remarkable prediction slipped into obscurity, until Dicke, Peebles, Roll and Wilkinson studied the problem again in 1965. Before they could measure the photon background, they learned that Penzias and Wilson had observed a weak isotropic background signal at a radio wavelength of 7.35 cm, corresponding to a blackbody temperature of $`T_\gamma =3.5\pm 1`$ K. They published their two papers back to back, with that of Dicke et al. explaining the fundamental significance of their measurement . Since then many different experiments have confirmed the existence of the microwave background. The most outstanding one has been the Cosmic Background Explorer (COBE) satellite, whose FIRAS instrument measured the photon background with great accuracy over a wide range of frequencies ($`\nu =197`$ cm<sup>-1</sup>), see Ref. , with a spectral resolution $`\frac{\mathrm{\Delta }\nu }{\nu }=0.0035`$. Nowadays, the photon spectrum is confirmed to be a blackbody spectrum with a temperature given by $$T_{_{\mathrm{CMB}}}=2.725\pm 0.002\mathrm{K}(\mathrm{systematic},95\%\mathrm{c}.\mathrm{l}.)\pm 7\mu \mathrm{K}(1\sigma \mathrm{statistical})$$ (55) In fact, this is the best blackbody spectrum ever measured, see Fig. 5, with spectral distortions below the level of 10 parts per million (ppm). Moreover, the differential microwave radiometer (DMR) instrument on COBE, with a resolution of about $`7^{}`$ in the sky, has also confirmed that it is an extraordinarily isotropic background. The deviations from isotropy, i.e. differences in the temperature of the blackbody spectrum measured in different directions in the sky, are of the order of 20 $`\mu `$K on large scales, or one part in $`10^5`$, see Ref. . There is, in fact, a dipole anisotropy of one part in $`10^3`$, $`\delta T_1=3.372\pm 0.007`$ mK (95% c.l.), in the direction of the Virgo cluster, $`(l,b)=(264.14^{}\pm 0.30,48.26^{}\pm 0.30)`$ (95% c.l.). Under the assumption that a Doppler effect is responsible for the entire CMB dipole, the velocity of the Sun with respect to the CMB rest frame is $`v_{}=371\pm 0.5`$ km/s, see Ref. .<sup>8</sup><sup>8</sup>8COBE even determined the annual variation due to the Earth’s motion around the Sun – the ultimate proof of Copernicus’ hypothesis. When subtracted, we are left with a whole spectrum of anisotropies in the higher multipoles (quadrupole, octupole, etc.), $`\delta T_2=18\pm 2\mu `$K (95% c.l.), see Ref. and Fig. 6. Soon after COBE, other groups quickly confirmed the detection of temperature anisotropies at around 30 $`\mu `$K and above, at higher multipole numbers or smaller angular scales. As I shall discuss below, these anisotropies play a crucial role in the understanding of the origin of structure in the universe. ### 2.3 Large-scale structure formation Although the isotropic microwave background indicates that the universe in the past was extraordinarily homogeneous, we know that the universe today is not exactly homogeneous: we observe galaxies, clusters and superclusters on large scales. These structures are expected to arise from very small primordial inhomogeneities that grow in time via gravitational instability, and that may have originated from tiny ripples in the metric, as matter fell into their troughs. Those ripples must have left some trace as temperature anisotropies in the microwave background, and indeed such anisotropies were finally discovered by the COBE satellite in 1992. The reason why they took so long to be discovered was that they appear as perturbations in temperature of only one part in $`10^5`$. While the predicted anisotropies have finally been seen in the CMB, not all kinds of matter and/or evolution of the universe can give rise to the structure we observe today. If we define the density contrast as $$\delta (\stackrel{}{x},a)\frac{\rho (\stackrel{}{x},a)\overline{\rho }(a)}{\overline{\rho }(a)}=d^3\stackrel{}{k}\delta _k(a)e^{i\stackrel{}{k}\stackrel{}{x}},$$ (56) where $`\overline{\rho }(a)=\rho _0a^3`$ is the average cosmic density, we need a theory that will grow a density contrast with amplitude $`\delta 10^5`$ at the last scattering surface ($`z=1100`$) up to density contrasts of the order of $`\delta 10^2`$ for galaxies at redshifts $`z1`$, i.e. today. This is a necessary requirement for any consistent theory of structure formation . Furthermore, the anisotropies observed by the COBE satellite correspond to a small-amplitude scale-invariant primordial power spectrum of inhomogeneities $$P(k)=\|\delta _k\|^2k^n,\mathrm{with}n=1,$$ (57) where the brackets $``$ represent integration over an ensemble of different universe realizations. These inhomogeneities are like waves in the space-time metric. When matter fell in the troughs of those waves, it created density perturbations that collapsed gravitationally to form galaxies and clusters of galaxies, with a spectrum that is also scale invariant. Such a type of spectrum was proposed in the early 1970s by Edward R. Harrison, and independently by the Russian cosmologist Yakov B. Zel’dovich, see Ref. , to explain the distribution of galaxies and clusters of galaxies on very large scales in our observable universe. Today various telescopes – like the Hubble Space Telescope, the twin Keck telescopes in Hawaii and the European Southern Observatory telescopes in Chile – are exploring the most distant regions of the universe and discovering the first galaxies at large distances. The furthest galaxies observed so far are at redshifts of $`z5`$, or 12 billion light years from the Earth, whose light was emitted when the universe had only about 5% of its present age. Only a few galaxies are known at those redshifts, but there are at present various catalogs like the CfA and APM galaxy catalogs, and more recently the IRAS Point Source redshift Catalog, see Fig. 7, and Las Campanas redshift surveys, that study the spatial distribution of hundreds of thousands of galaxies up to distances of a billion light years, or $`z<0.1`$, that recede from us at speeds of tens of thousands of kilometres per second. These catalogs are telling us about the evolution of clusters of galaxies in the universe, and already put constraints on the theory of structure formation. From these observations one can infer that most galaxies formed at redshifts of the order of $`26`$; clusters of galaxies formed at redshifts of order 1, and superclusters are forming now. That is, cosmic structure formed from the bottom up: from galaxies to clusters to superclusters, and not the other way around. This fundamental difference is an indication of the type of matter that gave rise to structure. The observed power spectrum of the galaxy matter distribution from a selection of deep redshift catalogs can be seen in Fig. 8. We know from Big Bang nucleosynthesis that all the baryons in the universe cannot account for the observed amount of matter, so there must be some extra matter (dark since we don’t see it) to account for its gravitational pull. Whether it is relativistic (hot) or non-relativistic (cold) could be inferred from observations: relativistic particles tend to diffuse from one concentration of matter to another, thus transferring energy among them and preventing the growth of structure on small scales. This is excluded by observations, so we conclude that most of the matter responsible for structure formation must be cold. How much there is is a matter of debate at the moment. Some recent analyses suggest that there is not enough cold dark matter to reach the critical density required to make the universe flat. If we want to make sense of the present observations, we must conclude that some other form of energy permeates the universe. In order to resolve this issue, even deeper galaxy redshift catalogs are underway, looking at millions of galaxies, like the Sloan Digital Sky Survey (SDSS) and the Anglo-Australian two degree field (2dF) Galaxy Redshift Survey, which are at this moment taking data, up to redshifts of $`z<0.5`$, over a large region of the sky. These important observations will help astronomers determine the nature of the dark matter and test the validity of the models of structure formation. Before COBE discovered the anisotropies of the microwave background there were serious doubts whether gravity alone could be responsible for the formation of the structure we observe in the universe today. It seemed that a new force was required to do the job. Fortunately, the anisotropies were found with the right amplitude for structure to be accounted for by gravitational collapse of primordial inhomogeneities under the attraction of a large component of non-relativistic dark matter. Nowadays, the standard theory of structure formation is a cold dark matter model with a non vanishing cosmological constant in a spatially flat universe. Gravitational collapse amplifies the density contrast initially through linear growth and later on via non-linear collapse. In the process, overdense regions decouple from the Hubble expansion to become bound systems, which start attracting eachother to form larger bound structures. In fact, the largest structures, superclusters, have not yet gone non-linear. The primordial spectrum (57) is reprocessed by gravitational instability after the universe becomes matter dominated and inhomogeneities can grow. Linear perturbation theory shows that the growing mode <sup>9</sup><sup>9</sup>9The decaying modes go like $`\delta (t)t^1`$, for all $`\omega `$. of small density contrasts go like $$\delta (a)a^{1+3\omega }=\{\begin{array}{cc}a^2,\hfill & a<a_{\mathrm{eq}}\hfill \\ a,\hfill & a>a_{\mathrm{eq}}\hfill \end{array}$$ (58) in the Einstein-de Sitter limit ($`\omega =p/\rho =1/3`$ and 0, for radiation and matter, respectively). There are slight deviations for $`aa_{\mathrm{eq}}`$, if $`\mathrm{\Omega }_\mathrm{M}1`$ or $`\mathrm{\Omega }_\mathrm{\Lambda }0`$, but we will not be concerned with them here. The important observation is that, since the density contrast at last scattering is of order $`\delta 10^5`$, and the scale factor has grown since then only a factor $`z_{\mathrm{dec}}10^3`$, one would expect a density contrast today of order $`\delta _010^2`$. Instead, we observe structures like galaxies, where $`\delta 10^2`$. So how can this be possible? The microwave background shows anisotropies due to fluctuations in the baryonic matter component only (to which photons couple, electromagnetically). If there is an additional matter component that only couples through very weak interactions, fluctuations in that component could grow as soon as it decoupled from the plasma, well before photons decoupled from baryons. The reason why baryonic inhomogeneities cannot grow is because of photon pressure: as baryons collapse towards denser regions, radiation pressure eventually halts the contraction and sets up acoustic oscillations in the plasma that prevent the growth of perturbations, until photon decoupling. On the other hand, a weakly interacting cold dark matter component could start gravitational collapse much earlier, even before matter-radiation equality, and thus reach the density contrast amplitudes observed today. The resolution of this mismatch is one of the strongest arguments for the existence of a weakly interacting cold dark matter component of the universe. How much dark matter there is in the universe can be deduced from the actual power spectrum (the Fourier transform of the two-point correlation function of density perturbations) of the observed large scale structure. One can decompose the density contrast in Fourier components, see Eq. (56). This is very convenient since in linear perturbation theory individual Fourier components evolve independently. A comoving wavenumber $`k`$ is said to “enter the horizon” when $`k=d_H^1(a)=aH(a)`$. If a certain perturbation, of wavelength $`\lambda =k^1<d_H(a_{\mathrm{eq}})`$, enters the horizon before matter-radiation equality, the fast radiation-driven expansion prevents dark-matter perturbations from collapsing. Since light can only cross regions that are smaller than the horizon, the suppression of growth due to radiation is restricted to scales smaller than the horizon, while large-scale perturbations remain unaffected. This is the reason why the horizon size at equality, Eq. (52), sets an important scale for structure growth, $$k_{\mathrm{eq}}=d_H^1(a_{\mathrm{eq}})0.083(\mathrm{\Omega }_\mathrm{M}h)h\mathrm{Mpc}^1.$$ (59) The suppression factor can be easily computed from (58) as $`f_{\mathrm{sup}}=(a_{\mathrm{enter}}/a_{\mathrm{eq}})^2=(k_{\mathrm{eq}}/k)^2`$. In other words, the processed power spectrum $`P(k)`$ will have the form: $$P(k)\{\begin{array}{cc}k,\hfill & kk_{\mathrm{eq}}\hfill \\ k^3,\hfill & kk_{\mathrm{eq}}\hfill \end{array}$$ (60) This is precisely the shape that large-scale galaxy catalogs are bound to test in the near future, see Fig. 9. Furthermore, since relativistic Hot Dark Matter (HDM) transfer energy between clumps of matter, they will wipe out small scale perturbations, and this should be seen as a distinctive signature in the matter power spectra of future galaxy catalogs. On the other hand, non-relativistic Cold Dark Matter (CDM) allow structure to form on all scales via gravitational collapse. The dark matter will then pull in the baryons, which will later shine and thus allow us to see the galaxies. Naturally, when baryons start to collapse onto dark matter potential wells, they will convert a large fraction of their potential energy into kinetic energy of protons and electrons, ionizing the medium. As a consequence, we expect to see a large fraction of those baryons constituting a hot ionized gas surrounding large clusters of galaxies. This is indeed what is observed, and confirms the general picture of structure formation. ## 3 DETERMINATION OF COSMOLOGICAL PARAMETERS In this Section, I will restrict myself to those recent measurements of the cosmological parameters by means of standard cosmological techniques, together with a few instances of new results from recently applied techniques. We will see that a large host of observations are determining the cosmological parameters with some reliability of the order of 10%. However, the majority of these measurements are dominated by large systematic errors. Most of the recent work in observational cosmology has been the search for virtually systematic-free observables, like those obtained from the microwave background anisotropies, and discussed in Section 4.4. I will devote, however, this Section to the more ‘classical’ measurements of the following cosmological parameters: The rate of expansion $`H_0`$; the matter content $`\mathrm{\Omega }_\mathrm{M}`$; the cosmological constant $`\mathrm{\Omega }_\mathrm{\Lambda }`$; the spatial curvature $`\mathrm{\Omega }_K`$, and the age of the universe $`t_0`$.<sup>10</sup><sup>10</sup>10We will take the baryon fraction as given by observations of light element abundances, in accordance with Big Bang nucleosynthesis, see Eq. (46). These five basic cosmological parameters are not mutually independent. Using the homogeneity and isotropy on large scales observed by COBE, we can infer relationships between the different cosmological parameters through the Einstein-Friedmann equations. In particular, we can deduce the value of the spatial curvature from the Cosmic Sum Rule, $$1=\mathrm{\Omega }_\mathrm{M}+\mathrm{\Omega }_\mathrm{\Lambda }+\mathrm{\Omega }_K,$$ (61) or viceversa, if we determine that the universe is spatially flat from observations of the microwave background, we can be sure that the sum of the matter content plus the cosmological constant must be one. Another relationship between parameters appears for the age of the universe. In a FRW cosmology, the cosmic expansion is determined by the Friedmann equation (8). Defining a new time and normalized scale factor, $$y\frac{a}{a_0}=\frac{1}{1+z},\tau H_0(tt_0),$$ (62) we can write the Friedmann equation with the help of the Cosmic Sum Rule (19) as $$y^{}(\tau )=\left[1+(y^11)\mathrm{\Omega }_\mathrm{M}+(y^21)\mathrm{\Omega }_\mathrm{\Lambda }\right]^{1/2},$$ (63) with initial condition $`y(0)=1,y^{}(0)=1`$. Therefore, the present age $`t_0`$ is a function of the other parameters, $`t_0=f(H_0,\mathrm{\Omega }_\mathrm{M},\mathrm{\Omega }_\mathrm{\Lambda })`$, determined from $$t_0H_0=_0^1𝑑y\left[1+(y^11)\mathrm{\Omega }_\mathrm{M}+(y^21)\mathrm{\Omega }_\mathrm{\Lambda }\right]^{1/2}.$$ (64) We show in Fig. 10 the contour lines for constant $`t_0H_0`$ in parameter space $`(\mathrm{\Omega }_\mathrm{M},\mathrm{\Omega }_\mathrm{\Lambda })`$. There are two specific limits of interest: an open universe with $`\mathrm{\Omega }_\mathrm{\Lambda }=0`$, for which the age is given by $$t_0H_0=\frac{1}{1\mathrm{\Omega }_\mathrm{M}}\frac{\mathrm{\Omega }_\mathrm{M}}{(1\mathrm{\Omega }_\mathrm{M})^{3/2}}\mathrm{ln}\left[\frac{1+(1\mathrm{\Omega }_\mathrm{M})^{1/2}}{\mathrm{\Omega }_\mathrm{M}^{1/2}}\right]=2\underset{n=0}{\overset{\mathrm{}}{}}\frac{(1\mathrm{\Omega }_\mathrm{M})^n}{(2n+1)(2n+3)},$$ (65) and a flat universe with $`\mathrm{\Omega }_\mathrm{\Lambda }=1\mathrm{\Omega }_\mathrm{M}`$, for which the age can also be expressed in compact form, $$t_0H_0=\frac{2}{3(1\mathrm{\Omega }_\mathrm{M})^{1/2}}\mathrm{ln}\left[\frac{1+(1\mathrm{\Omega }_\mathrm{M})^{1/2}}{\mathrm{\Omega }_\mathrm{M}^{1/2}}\right]=\frac{2}{3}\underset{n=0}{\overset{\mathrm{}}{}}\frac{(1\mathrm{\Omega }_\mathrm{M})^n}{2n+1}.$$ (66) We have plotted these functions in Fig. 11. It is clear that in both cases $`t_0H_02/3`$ as $`\mathrm{\Omega }_\mathrm{M}1`$. We can now use these relations as a consistency check between the cosmological observations of $`H_0`$, $`\mathrm{\Omega }_\mathrm{M}`$, $`\mathrm{\Omega }_\mathrm{\Lambda }`$ and $`t_0`$. Of course, we cannot measure the age of the universe directly, but only the age of its constituents: stars, galaxies, globular clusters, etc. Thus we can only find a lower bound on the age of the universe, $`t_0>t_{\mathrm{gal}}+1.5`$ Gyr. As we will see, this is not a trivial bound and, in several occasions, during the progress towards better determinations of the cosmological parameters, the universe seemed to be younger than its constituents, a logical inconsistency, of course, only due to an incorrect assessment of systematic errors . In order to understand those recent measurements, one should also define what is known as the luminosity distance to an object in the universe. Imagine a source that is emitting light at a distance $`d_L`$ from a detector of area $`dA`$. The absolute luminosity $``$ of such a source is nothing but the energy emitted per unit time. A standard candle is a luminous object that can be calibrated with some accuracy and therefore whose absolute luminosity is known, within certain errors. For example, Cepheid variable stars and type Ia supernovae are considered to be reasonable standard candles, i.e. their calibration errors are within bounds. The energy flux $``$ received at the detector is the measured energy per unit time per unit area of the detector coming from that source. The luminosity distance $`d_L`$ is then defined as the radius of the sphere centered on the source for which the absolute luminosity would give the observed flux, $`/4\pi d_L^2`$. In a Friedmann-Robertson-Walker universe, light travels along null geodesics, $`ds^2=0`$, or, see Eq. (2), $$\frac{dr}{\sqrt{1+a_0^2H_0^2r^2\mathrm{\Omega }_K}}=\frac{1}{a_0^2H_0^2}\frac{dz}{\sqrt{(1+z)^2(1+z\mathrm{\Omega }_\mathrm{M})z(2+z)\mathrm{\Omega }_\mathrm{\Lambda }}},$$ (67) which determines the coordinate distance $`r=r(z,H_0,\mathrm{\Omega }_\mathrm{M},\mathrm{\Omega }_\mathrm{\Lambda })`$, as a function of redshift $`z`$ and the other cosmological parameters. Now let us consider the effect of the universe expansion on the observed flux coming from a source at a certain redshift $`z`$ from us. First, the photon energy on its way here will be redshifted, and thus the observed energy $`E_0=E/(1+z)`$. Second, the rate of photon arrival will be time-delayed with respect to that emitted by the source, $`dt_0=(1+z)dt`$. Finally, the fraction of the area of the 2-sphere centered on the source that is covered by the detector is $`dA/4\pi a_0^2r^2(z)`$. Therefore, the total flux detected is $$=\frac{}{4\pi a_0^2r^2(z)}\frac{}{4\pi d_L^2}.$$ (68) The final expression for the luminosity distance $`d_L`$ as a function of redshift is thus given by $$H_0d_L=(1+z)\|\mathrm{\Omega }_K\|^{1/2}\mathrm{sinn}\left[\|\mathrm{\Omega }_K\|^{1/2}_0^z\frac{dz^{}}{\sqrt{(1+z^{})^2(1+z^{}\mathrm{\Omega }_\mathrm{M})z^{}(2+z^{})\mathrm{\Omega }_\mathrm{\Lambda }}}\right],$$ (69) where $`\mathrm{sinn}(x)=x\mathrm{if}K=0;\mathrm{sin}(x)\mathrm{if}K=+1\mathrm{and}\mathrm{sinh}(x)\mathrm{if}K=1`$. Expanding to second order around $`z=0`$, we obtain Eq. (6), $$H_0d_L=z+\frac{1}{2}\left(1\frac{\mathrm{\Omega }_\mathrm{M}}{2}+\mathrm{\Omega }_\mathrm{\Lambda }\right)z^2+𝒪(z^3).$$ (70) This expression goes beyond the leading linear term, corresponding to the Hubble law, into the second order term, which is sensitive to the cosmological parameters $`\mathrm{\Omega }_\mathrm{M}`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }`$. It is only recently that cosmological observations have gone far enough back into the early universe that we can begin to probe the second term, as I will discuss shortly. Higher order terms are not yet probed by cosmological observations, but they would contribute as important consistency checks. Let us now pursue the analysis of the recent determinations of the most important cosmological parameters: the rate of expansion $`H_0`$, the matter content $`\mathrm{\Omega }_\mathrm{M}`$, the cosmological constant $`\mathrm{\Omega }_\mathrm{\Lambda }`$, the spatial curvature $`\mathrm{\Omega }_K`$, and the age of the universe $`t_0`$. ### 3.1 The rate of expansion $`H_0`$ Over most of last century the value of $`H_0`$ has been a constant source of disagreement . Around 1929, Hubble measured the rate of expansion to be $`H_0=500`$ km s<sup>-1</sup>Mpc<sup>-1</sup>, which implied an age of the universe of order $`t_02`$ Gyr, in clear conflict with geology. Hubble’s data was based on Cepheid standard candles that were incorrectly calibrated with those in the Large Magellanic Cloud. Later on, in 1954 Baade recalibrated the Cepheid distance and obtained a lower value, $`H_0=250`$ km s<sup>-1</sup>Mpc<sup>-1</sup>, still in conflict with ratios of certain unstable isotopes. Finally, in 1958 Sandage realized that the brightest stars in galaxies were ionized HII regions, and the Hubble rate dropped down to $`H_0=60`$ km s<sup>-1</sup> Mpc<sup>-1</sup>, still with large (factor of two) systematic errors. Fortunately, in the past 15 years there has been significant progress towards the determination of $`H_0`$, with systematic errors approaching the 10% level. These improvements come from two directions. First, technological, through the replacement of photographic plates (almost exclusively the source of data from the 1920s to 1980s) with charged couple devices (CCDs), i.e. solid state detectors with excellent flux sensitivity per pixel, which were previously used successfully in particle physics detectors. Second, by the refinement of existing methods for measuring extragalactic distances (e.g. parallax, Cepheids, supernovae, etc.). Finally, with the development of completely new methods to determine $`H_0`$, which fall into totally independent and very broad categories: a) Gravitational lensing; b) Sunyaev-Zel’dovich effect; c) Extragalactic distance scale, mainly Cepheid variability and type Ia Supernovae; d) Microwave background anisotropies. I will review here the first three, and leave the last method for Section 4.4, since it involves knowledge about the primordial spectrum of inhomogeneities. #### 3.1.1 Gravitational lensing Imagine a quasi-stellar object (QSO) at large redshift ($`z1`$) whose light is lensed by an intervening galaxy at redshift $`z1`$ and arrives to an observer at $`z=0`$. There will be at least two different images of the same background variable point source. The arrival times of photons from two different gravitationally lensed images of the quasar depend on the different path lengths and the gravitational potential traversed. Therefore, a measurement of the time delay and the angular separation of the different images of a variable quasar can be used to determine $`H_0`$ with great accuracy. This method, proposed in 1964 by Refsdael , offers tremendous potential because it can be applied at great distances and it is based on very solid physical principles . Unfortunately, there are very few systems with both a favourable geometry (i.e. a known mass distribution of the intervening galaxy) and a variable background source with a measurable time delay. That is the reason why it has taken so much time since the original proposal for the first results to come out. Fortunately, there are now very powerful telescopes that can be used for these purposes. The best candidate to-date is the QSO $`0957+561`$, observed with the 10m Keck telescope, for which there is a model of the lensing mass distribution that is consistent with the measured velocity dispersion. Assuming a flat space with $`\mathrm{\Omega }_\mathrm{M}=0.25`$, one can determine $$H_0=72\pm 7(1\sigma \mathrm{statistical})\pm 15\%(\mathrm{systematic})\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1.$$ (71) The main source of systematic error is the degeneracy between the mass distribution of the lens and the value of $`H_0`$. Knowledge of the velocity dispersion within the lens as a function of position helps constrain the mass distribution, but those measurements are very difficult and, in the case of lensing by a cluster of galaxies, the dark matter distribution in those systems is usually unknown, associated with a complicated cluster potential. Nevertheless, the method is just starting to give promising results and, in the near future, with the recent discovery of several systems with optimum properties, the prospects for measuring $`H_0`$ and lowering its uncertainty with this technique are excellent. #### 3.1.2 Sunyaev-Zel’dovich effect As discussed in the previous Section, the gravitational collapse of baryons onto the potential wells generated by dark matter gave rise to the reionization of the plasma, generating an X-ray halo around rich clusters of galaxies, see Fig. 12. The inverse-Compton scattering of microwave background photons off the hot electrons in the X-ray gas results in a measurable distortion of the blackbody spectrum of the microwave background, known as the Sunyaev-Zel’dovich (SZ) effect. Since photons acquire extra energy from the X-ray electrons, we expect a shift towards higher frequencies of the spectrum, $`(\mathrm{\Delta }\nu /\nu )(k_\mathrm{B}T_{\mathrm{gas}}/m_ec^2)10^2`$. This corresponds to a decrement of the microwave background temperature at low frequencies (Rayleigh-Jeans region) and an increment at high frequencies, see Ref. . Measuring the spatial distribution of the SZ effect (3 K spectrum), together with a high resolution X-ray map ($`10^8`$ K spectrum) of the cluster, one can determine the density and temperature distribution of the hot gas. Since the X-ray flux is distance-dependent ($`=/4\pi d_L^2`$), while the SZ decrement is not (because the energy of the CMB photons increases as we go back in redshift, $`\nu =\nu _0(1+z)`$, and exactly compensates the redshift in energy of the photons that reach us), one can determine from there the distance to the cluster, and thus the Hubble rate $`H_0`$. The advantages of this method are that it can be applied to large distances and it is based on clear physical principles. The main systematics come from possible clumpiness of the gas (which would reduce $`H_0`$), projection effects (if the clusters are prolate, $`H_0`$ could be larger), the assumption of hydrostatic equilibrium of the X-ray gas, details of models for the gas and electron densities, and possible contaminations from point sources. Present measurements give the value $$H_0=60\pm 10(1\sigma \mathrm{statistical})\pm 20\%(\mathrm{systematic})\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1,$$ (72) compatible with other determinations. A great advantage of this completely new and independent method is that nowadays more and more clusters are observed in the X-ray, and soon we will have high-resolution 2D maps of the SZ decrement from several balloon flights, as well as from future microwave background satellites, together with precise X-ray maps and spectra from the Chandra X-ray observatory recently launched by NASA, as well as from the European X-ray satellite XMM launched a few months ago by ESA, which will deliver orders of magnitude better resolution than the existing Einstein X-ray satellite. #### 3.1.3 Cepheid variability Cepheids are low-mass variable stars with a period-luminosity relation based on the helium ionization cycles inside the star, as it contracts and expands. This time variability can be measured, and the star’s absolute luminosity determined from the calibrated relationship. From the observed flux one can then deduce the luminosity distance, see Eq. (69), and thus the Hubble rate $`H_0`$. The Hubble Space Telescope (HST) was launched by NASA in 1990 (and repaired in 1993) with the specific project of calibrating the extragalactic distance scale and thus determining the Hubble rate with 10% accuracy. The most recent results from HST are the following $$H_0=71\pm 4(\mathrm{random})\pm 7(\mathrm{systematic})\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1.$$ (73) The main source of systematic error is the distance to the Large Magellanic Cloud, which provides the fiducial comparison for Cepheids in more distant galaxies. Other systematic uncertainties that affect the value of $`H_0`$ are the internal extinction correction method used, a possible metallicity dependence of the Cepheid period-luminosity relation and cluster population incompleteness bias, for a set of 21 galaxies within 25 Mpc, and 23 clusters within $`z<0.03`$. With better telescopes coming up soon, like the Very Large Telescope (VLT) interferometer of the European Southern Observatory (ESO) in the Chilean Atacama desert, with 4 synchronized telescopes by the year 2005, and the Next Generation Space Telescope (NGST) proposed by NASA for 2008, it is expected that much better resolution and therefore accuracy can be obtained for the determination of $`H_0`$. ### 3.2 The matter content $`\mathrm{\Omega }_\mathrm{M}`$ In the 1920s Hubble realized that the so called nebulae were actually distant galaxies very similar to our own. Soon afterwards, in 1933, Zwicky found dynamical evidence that there is possibly ten to a hundred times more mass in the Coma cluster than contributed by the luminous matter in galaxies . However, it was not until the 1970s that the existence of dark matter began to be taken more seriously. At that time there was evidence that rotation curves of galaxies did not fall off with radius and that the dynamical mass was increasing with scale from that of individual galaxies up to clusters of galaxies. Since then, new possible extra sources to the matter content of the universe have been accumulating: $`\mathrm{\Omega }_\mathrm{M}`$ $`=`$ $`\mathrm{\Omega }_{\mathrm{B},\mathrm{lum}}(\mathrm{stars}\mathrm{in}\mathrm{galaxies})`$ (74) $`+`$ $`\mathrm{\Omega }_{\mathrm{B},\mathrm{dark}}(\mathrm{MACHOs}\mathrm{?})`$ (75) $`+`$ $`\mathrm{\Omega }_{\mathrm{CDM}}(\mathrm{weakly}\mathrm{interacting}:\mathrm{axion},\mathrm{neutralino}\mathrm{?})`$ (76) $`+`$ $`\mathrm{\Omega }_{\mathrm{HDM}}(\mathrm{massive}\mathrm{neutrinos}\mathrm{?})`$ (77) The empirical route to the determination of $`\mathrm{\Omega }_\mathrm{M}`$ is nowadays one of the most diversified of all cosmological parameters. The matter content of the universe can be deduced from the mass-to-light ratio of various objects in the universe; from the rotation curves of galaxies; from microlensing and the direct search of Massive Compact Halo Objects (MACHOs); from the cluster velocity dispersion with the use of the Virial theorem; from the baryon fraction in the X-ray gas of clusters; from weak gravitational lensing; from the observed matter distribution of the universe via its power spectrum; from the cluster abundance and its evolution; from direct detection of massive neutrinos at SuperKamiokande; from direct detection of Weakly Interacting Massive Particles (WIMPs) at DAMA and UKDMC, and finally from microwave background anisotropies. I will review here just a few of them. #### 3.2.1 Luminous matter The most straight forward method of estimating $`\mathrm{\Omega }_\mathrm{M}`$ is to measure the luminosity of stars in galaxies and then estimate the mass-to-light ratio, defined as the mass per luminosity density observed from an object, $`\mathrm{{\rm Y}}=/`$. This ratio is usually expressed in solar units, $`_{}/_{}`$, so that for the sun $`\mathrm{{\rm Y}}_{}=1`$. The luminosity of stars depends very sensitively on their mass and stage of evolution. The mass-to-light ratio of stars in the solar neighbourhood is of order $`\mathrm{{\rm Y}}3`$. For globular clusters and spiral galaxies we can determine their mass and luminosity independently and this gives $`\mathrm{{\rm Y}}`$ few. For our galaxy, $$_{\mathrm{gal}}=(1.0\pm 0.3)\times 10^8hL_{}\mathrm{Mpc}^3\mathrm{and}\mathrm{{\rm Y}}_{\mathrm{gal}}=6\pm 3.$$ (78) The contribution of galaxies to the luminosity density of the universe (in the visible-V spectral band, centered at $`5500`$ Å) is $$_V=(1.7\pm 0.6)\times 10^8hL_{}\mathrm{Mpc}^3,$$ (79) which can be translated into a mass density by multiplying by the observed $`\mathrm{{\rm Y}}`$ in that band, $$\mathrm{\Omega }_\mathrm{M}h=(6.1\pm 2.2)\times 10^4\mathrm{{\rm Y}}_V.$$ (80) All the luminous matter in the universe, from galaxies, clusters of galaxies, etc., account for $`\mathrm{{\rm Y}}10`$, and thus $$0.002\mathrm{\Omega }_{\mathrm{lum}}h0.006.$$ (81) As a consequence, the luminous matter alone is far from the critical density. Moreover, comparing with the amount of baryons from Big Bang nucleosynthesis (46), we conclude that $`\mathrm{\Omega }_{\mathrm{lum}}\mathrm{\Omega }_\mathrm{B}`$, so there must be a large fraction of baryons that are dark, perhaps in the form of very dim stars. #### 3.2.2 Rotation curves of spiral galaxies The flat rotation curves of spiral galaxies provide the most direct evidence for the existence of large amounts of dark matter. Spiral galaxies consist of a central bulge and a very thin disk, stabilized against gravitational collapse by angular momentum conservation, and surrounded by an approximately spherical halo of dark matter. One can measure the orbital velocities of objects orbiting around the disk as a function of radius from the Doppler shifts of their spectral lines. The rotation curve of the Andromeda galaxy was first measured by Babcock in 1938, from the stars in the disk. Later it became possible to measure galactic rotation curves far out into the disk, and a trend was found . The orbital velocity rose linearly from the center outward until it reached a typical value of 200 km/s, and then remained flat out to the largest measured radii. This was completely unexpected since the observed surface luminosity of the disk falls off exponentially with radius, $`I(r)=I_0\mathrm{exp}(r/r_D)`$, see Ref. . Therefore, one would expect that most of the galactic mass is concentrated within a few disk lengths $`r_D`$, such that the rotation velocity is determined as in a Keplerian orbit, $`v_{\mathrm{rot}}=(GM/r)^{1/2}r^{1/2}`$. No such behaviour is observed. In fact, the most convincing observations come from radio emission (from the 21 cm line) of neutral hydrogen in the disk, which has been measured to much larger galactic radii than optical tracers. A typical case is that of the spiral galaxy NGC 6503, where $`r_D=1.73`$ kpc, while the furthest measured hydrogen line is at $`r=22.22`$ kpc, about 13 disk lengths away. The measured rotation curve is shown in Fig. 13 together with the relative components associated with the disk, the halo and the gas. Nowadays, thousands of galactic rotation curves are known, and all suggest the existence of about ten times more mass in the halos of spiral galaxies than in the stars of the disk. Recent numerical simulations of galaxy formation in a CDM cosmology suggest that galaxies probably formed by the infall of material in an overdense region of the universe that had decoupled from the overall expansion. The dark matter is supposed to undergo violent relaxation and create a virialized system, i.e. in hydrostatic equilibrium. This picture has led to a simple model of dark-matter halos as isothermal spheres, with density profile $`\rho (r)=\rho _c/(r_c^2+r^2)`$, where $`r_c`$ is a core radius and $`\rho _c=v_{\mathrm{}}^2/4\pi G`$, with $`v_{\mathrm{}}`$ equal to the plateau value of the flat rotation curve. This model is consistent with the universal rotation curve seen in Fig. 13. At large radii the dark matter distribution leads to a flat rotation curve. Adding up all the matter in galactic halos up to maximum radii, one finds $`\mathrm{{\rm Y}}_{\mathrm{halo}}30h`$, and therefore $$\mathrm{\Omega }_{\mathrm{halo}}0.030.05.$$ (82) Of course, it would be extraordinary if we could confirm, through direct detection, the existence of dark matter in our own galaxy. For that purpose, one should measure its rotation curve, which is much more difficult because of obscuration by dust in the disk, as well as problems with the determination of reliable galactocentric distances for the tracers. Nevertheless, the rotation curve of the Milky Way has been measured and conforms to the usual picture, with a plateau value of the rotation velocity of 220 km/s, see Ref. . For dark matter searches, the crucial quantity is the dark matter density in the solar neighbourhood, which turns out to be (within a factor of two uncertainty depending on the halo model) $`\rho _{\mathrm{DM}}=0.3`$ GeV/cm<sup>3</sup>. We will come back to direct searched of dark matter in a later subsection. #### 3.2.3 Microlensing The existence of large amounts of dark matter in the universe, and in our own galaxy in particular, is now established beyond any reasonable doubt, but its nature remains a mystery. We have seen that baryons cannot account for the whole matter content of the universe; however, since the contribution of the halo (82) is comparable in magnitude to the baryon fraction of the universe (46), one may ask whether the galactic halo could be made of purely baryonic material in some non-luminous form, and if so, how one should search for it. In other words, are MACHOs the non-luminous baryons filling the gap between $`\mathrm{\Omega }_{\mathrm{lum}}`$ and $`\mathrm{\Omega }_\mathrm{B}`$? If not, what are they? Let us start a systematic search for possibilities. They cannot be normal stars since they would be luminous; neither hot gas since it would shine; nor cold gas since it would absorb light and reemit in the infrared. Could they be burnt-out stellar remnants? This seems implausible since they would arise from a population of normal stars of which there is no trace in the halo. Neutron stars or black holes would typically arise from Supernova explosions and thus eject heavy elements into the galaxy, while the overproduction of helium in the halo is strongly constrained. They could be white dwarfs, i.e. stars not massive enough to reach supernova phase. Despite some recent arguments, a halo composed by white dwarfs is not rigorously excluded. Are they stars too small to shine? Perhaps M-dwarfs, stars with a mass $`M0.1M_{}`$ which are intrinsically dim; however, very long exposure images of the Hubble Space Telescope restrict the possible M-dwarf contribution to the galaxy to be below 6%. The most plausible alternative is a halo composed of brown dwarfs with mass $`M0.08M_{}`$, which never ignite hydrogen and thus shine only from the residual energy due to gravitational contraction.<sup>11</sup><sup>11</sup>11A sometimes discussed alternative, planet-size Jupiters, can be classified as low-mass brown dwarfs. In fact, the extrapolation of the stellar mass function to small masses predicts a large number of brown dwarfs within normal stellar populations. A final possibility is primordial black holes (PBH), which could have been created in the early universe from early phase transitions , even before baryons were formed, and thus may be classified as non-baryonic. They could make a large contribution towards the total $`\mathrm{\Omega }_\mathrm{M}`$, and still be compatible with Big Bang nucleosynthesis. Whatever the arguments for or against baryonic objects as galactic dark matter, nothing would be more convincing than a direct detection of the various candidates, or their exclusion, in a direct search experiment. Fortunately, in 1986 Paczyński proposed a method for detecting faint stars in the halo of our galaxy . The idea is based on the well known effect that a point-like mass deflector placed between an observer and a light source creates two different images, as shown in Fig. 14. When the source is exactly aligned with the deflector of mass $`M_D`$, the image would be an annulus, an Einstein ring, with radius $$r_\mathrm{E}^2=4GM_Dd,\mathrm{where}d=\frac{d_1d_2}{d_1+d_2}$$ (83) is the reduced distance to the source, see Fig. 14. If the two images cannot be separated because their angular distance $`\alpha `$ is below the resolving power of the observer’s telescope, the only effect will be an apparent brightening of the star, an effect known as gravitational microlensing. The amplification factor is $$A=\frac{2+u^2}{u\sqrt{4+u^2}},\mathrm{where}u\frac{r}{r_\mathrm{E}},$$ (84) with $`r`$ the distance from the line of sight to the deflector. Imagine an observer on Earth watching a distant star in the Large Magellanic Cloud (LMC), 50 kpc away. If the galactic halo is filled with MACHOs, one of them will occasionally pass near the line of sight and thus cause the image of the background star to brighten. If the MACHO moves with velocity $`v`$ transverse to the line of sight, and if its impact parameter, i.e. the minimal distance to the line of sight, is $`b`$, then one expects an apparent lightcurve as shown in Fig. 15 for different values of $`b/r_\mathrm{E}`$. The natural time unit is $`\mathrm{\Delta }t=r_\mathrm{E}/v`$, and the origin corresponds to the time of closest approach to the line of sight. The probability for a target star to be lensed is independent of the mass of the dark matter object . For stars in the LMC one finds a probability, i.e. an optical depth for microlensing of the galactic halo, of approximately $`\tau 10^6`$. Thus, if one looks simultaneously at several millions of stars in the LMC during extended periods of time, one has a good chance of seeing at least a few of them brightened by a dark halo object. In order to be sure one has seen a microlensing event one has to monitor a large sample of stars long enough to identify the characteristic light curve shown in Fig. 15. The unequivocal signatures of such an event are the following: it must be a) unique (non-repetitive in time); b) time-symmetric; and c) achromatic (because of general covariance). These signatures allow one to discriminate against variable stars which constitute the background. The typical duration of the light curve is the time it takes a MACHO to cross an Einstein radius, $`\mathrm{\Delta }t=r_\mathrm{E}/v`$. If the deflector mass is $`1M_{}`$, the average microlensing time will be 3 months, for $`10^2M_{}`$ it is 9 days, for $`10^4M_{}`$ it is 1 day, and for $`10^6M_{}`$ it is 2 hours. A characteristic event, of duration 34 days, is shown in Fig. 16. The first microlensing events towards the LMC were reported by the MACHO and EROS collaborations in 1993 . Nowadays, there are 12 candidates towards the LMC, 2 towards the SMC, around 40 towards the bulge of our own galaxy, and about 2 towards Andromeda, seen by AGAPE , with a slightly different technique based on pixel brightening rather than individual stars. Thus, microlensing is a well established technique with a rather robust future. In particular, it has allowed the MACHO and EROS collaboration to draw exclusion plots for various mass ranges in terms of their maximum allowed halo fraction, see Fig. 17. The MACHO Collaboration conclude in their 5-year analysis, see Ref. , that the spatial distribution of events is consistent with an extended lens distribution such as Milky Way or LMC halo, consisting partially of compact objects. A maximum likelihood analysis gives a MACHO halo fraction of 20% for a typical halo model with a 95% confidence interval of 8% to 50%. A 100% MACHO halo is ruled out at 95% c.l. for all except their most extreme halo model. The most likely MACHO mass is between 0.15 $`M_{}`$ and 0.9 $`M_{}`$, depending on the halo model. The lower mass is characteristic of white dwarfs, but a galactic halo composed primarily of white dwarfs is barely compatible with a range of observational constraints. On the other hand, if one wanted to attribute the observed events to brown dwarfs, one needs to appeal to a very non-standard density and/or velocity distribution of these objects. It is still unclear what sort of objects the microlensing experiments are seeing towards the LMC and where the lenses are. Nevertheless, the field is expanding, with several new experiments already underway, to search for clear signals of parallax, or binary systems, where the degeneracy between mass and distance can be resolved. For a discussion of those new results, see Ref. . #### 3.2.4 Virial theorem and large scale motion Clusters of galaxies are the largest gravitationally bound systems in the universe (superclusters are not yet in equilibrium). We know today several thousand clusters; they have typical radii of $`15`$ Mpc and typical masses of $`29`$ $`\times 10^{14}M_{}`$. Zwicky noted in 1933 that these systems appear to have large amounts of dark matter . He used the virial theorem (for a gravitationally bound system in equilibrium), $`2E_{\mathrm{kin}}=E_{\mathrm{grav}}`$, where $`E_{\mathrm{kin}}=\frac{1}{2}mv^2`$ is the average kinetic energy of one of the bound objects (galaxies) of mass $`m`$ and $`E_{\mathrm{grav}}=mGM/r`$ is the average gravitational potential energy caused by the attraction of the other galaxies. Measuring the velocity dispersion $`v^2`$ from the Doppler shifts of the spectral lines and estimating the geometrical size of the system gives an estimate of its total mass $`M`$. As Zwicky noted, this virial mass of clusters far exceeds their luminous mass, typically leading to a mass-to-light ratio $`\mathrm{{\rm Y}}_{\mathrm{cluster}}=200\pm 70`$. Assuming that the average cluster $`\mathrm{{\rm Y}}`$ is representative of the entire universe <sup>12</sup><sup>12</sup>12Recent observations indicate that $`\mathrm{{\rm Y}}`$ is independent of scale up to supercluster scales $`100h^1`$ Mpc. one finds for the cosmic matter density $$\mathrm{\Omega }_\mathrm{M}=0.24\pm 0.05(1\sigma \mathrm{statistical})\pm 0.09(\mathrm{systematic}).$$ (85) On scales larger than clusters the motion of galaxies is dominated by the overall cosmic expansion. Nevertheless, galaxies exhibit peculiar velocities with respect to the global cosmic flow. For example, our Local Group of galaxies is moving with a speed of $`627\pm 22`$ km/s relative to the cosmic microwave background reference frame, towards the Great Attractor. In the context of the standard gravitational instability theory of structure formation, the peculiar motions of galaxies are attributed to the action of gravity during the universe evolution, caused by the matter density inhomogeneities which give rise to the formation of structure. The observed large-scale velocity fields, together with the observed galaxy distributions, can then be translated into a measure for the mass-to-light ratio required to explain the large-scale flows. An example of the reconstruction of the matter density field in our cosmological vicinity from the observed velocity field is shown in Fig. 18. The cosmic matter density inferred from such analyses is $$\mathrm{\Omega }_\mathrm{M}>0.395\%\mathrm{c}.\mathrm{l}.$$ (86) Related methods that are more model-dependent give even larger estimates. #### 3.2.5 Baryon fraction in clusters Since large clusters of galaxies form through gravitational collapse, they scoop up mass over a large volume of space, and therefore the ratio of baryons over the total matter in the cluster should be representative of the entire universe, at least within a 20% systematic error. Since the 1960s, when X-ray telescopes became available, it is known that galaxy clusters are the most powerful X-ray sources in the sky . The emission extends over the whole cluster and reveals the existence of a hot plasma with temperature $`T10^710^8`$ K, where X-rays are produced by electron bremsstrahlung. Assuming the gas to be in hydrostatic equilibrium and applying the virial theorem one can estimate the total mass in the cluster, giving general agreement (within a factor of 2) with the virial mass estimates. From these estimates one can calculate the baryon fraction of clusters $$f_\mathrm{B}h^{3/2}=0.030.08\frac{\mathrm{\Omega }_\mathrm{B}}{\mathrm{\Omega }_\mathrm{M}}0.15,\mathrm{for}h=0.65,$$ (87) which together with (81) indicates that clusters contain far more baryonic matter in the form of hot gas than in the form of stars in galaxies. Assuming this fraction to be representative of the entire universe, and using the Big Bang nucleosynthesis value of $`\mathrm{\Omega }_\mathrm{B}=0.05\pm 0.01`$, for $`h=0.65`$, we find $$\mathrm{\Omega }_\mathrm{M}=0.3\pm 0.1(\mathrm{statistical})\pm 20\%(\mathrm{systematic}).$$ (88) This value is consistent with previous determinations of $`\mathrm{\Omega }_\mathrm{M}`$. If some baryons are ejected from the cluster during gravitational collapse, or some are actually bound in nonluminous objects like planets, then the actual value of $`\mathrm{\Omega }_\mathrm{M}`$ is smaller than this estimate. #### 3.2.6 Weak gravitational lensing Since the mid 1980s, deep surveys with powerful telescopes have observed huge arc-like features in galaxy clusters, see for instance Fig. 19. The spectroscopic analysis showed that the cluster and the giant arcs were at very different redshifts. The usual interpretation is that the arc is the image of a distant background galaxy which is in the same line of sight as the cluster so that it appears distorted and magnified by the gravitational lens effect: the giant arcs are essentially partial Einstein rings. From a systematic study of the cluster mass distribution one can reconstruct the shear field responsible for the gravitational distortion, see Ref. . This analysis shows that there are large amounts of dark matter in the clusters, in rough agreement with the virial mass estimates, although the lensing masses tend to be systematically larger. At present, the estimates indicate $`\mathrm{\Omega }_\mathrm{M}=0.20.3`$ on scales $`<6h^1`$ Mpc, while $`\mathrm{\Omega }_\mathrm{M}=0.4`$ for the Corona Borealis supercluster, on scales of order 20 Mpc. #### 3.2.7 Structure formation and the matter power spectrum One the most important constraints on the amount of matter in the universe comes from the present distribution of galaxies. As we mentioned in the Section 2.3, gravitational instability increases the primordial density contrast, seen at the last scattering surface as temperature anisotropies, into the present density field responsible for the large and the small scale structure. Since the primordial spectrum is very approximately represented by a scale-invariant Gaussian random field, the best way to present the results of structure formation is by working with the 2-point correlation function in Fourier space (the equivalent to the Green’s function in QFT), the so-called power spectrum. If the reprocessed spectrum of inhomogeneities remains Gaussian, the power spectrum is all we need to describe the galaxy distribution. Non-Gaussian effects are expected to arise from the non-linear gravitational collapse of structure, and may be important at small scales . The power spectrum measures the degree of inhomogeneity in the mass distribution on different scales. It depends upon a few basic ingredientes: a) the primordial spectrum of inhomogeneities, whether they are Gaussian or non-Gaussian, whether adiabatic (perturbations in the energy density) or isocurvature (perturbations in the entropy density), whether the primordial spectrum has tilt (deviations from scale-invariance), etc.; b) the recent creation of inhomogeneities, whether cosmic strings or some other topological defect from an early phase transition are responsible for the formation of structure today; and c) the cosmic evolution of the inhomogeneity, whether the universe has been dominated by cold or hot dark matter or by a cosmological constant since the beginning of structure formation, and also depending on the rate of expansion of the universe. The working tools used for the comparison between the observed power spectrum and the predicted one are very precise N-body numerical simulations and theoretical models that predict the shape but not the amplitude of the present power spectrum. Even though a large amount of work has gone into those analyses, we still have large uncertainties about the nature and amount of matter necessary for structure formation. A model that has become a working paradigm is a flat cold dark matter model with a cosmological constant and $`\mathrm{\Omega }_\mathrm{M}=0.30.4`$. This model will soon be confronted with very precise measurements from SDSS, 2dF, and several other large redshift catalogs, that are already taking data, see Section 4.5. The observational constraints on the power spectrum have a huge lever arm of measurements at very different scales, mainly from the observed cluster abundance, on 10 Mpc scales, to the CMB fluctuations, on 1000 Mpc scales, which determines the normalization of the spectrum. At present, deep redshift surveys are probing scales between 100 and 1000 Mpc, which should begin to see the turnover corresponding to the peak of the power spectrum at $`k_{\mathrm{eq}}`$, see Figs. 8 and 9. The standard CDM model with $`\mathrm{\Omega }_\mathrm{M}=1`$, normalized to the CMB fluctuations on large scales, is inconsistent with the cluster abundance. The power spectra of both a flat model with a cosmological constant or an open universe with $`\mathrm{\Omega }_\mathrm{M}=0.3`$ (defined as $`\mathrm{\Lambda }`$CDM and OCDM, respectively) can be normalized so that they agree with both the CMB and cluster observations. In the near future, galaxy survey observations will greatly improve the power spectrum constraints and will allow a measurement of $`\mathrm{\Omega }_\mathrm{M}`$ from the shape of the spectrum. At present, these measurements suggest a low value of $`\mathrm{\Omega }_\mathrm{M}`$, but with large uncertainties. #### 3.2.8 Cluster abundance and evolution Rich clusters are the most recently formed gravitationally bound systems in the universe. Their number density as a function of time (or redshift) helps determine the amount of dark matter. The observed present ($`z0`$) cluster abundance provides a strong constraint on the normalization of the power spectrum of density perturbations on cluster scales. Both $`\mathrm{\Lambda }`$CDM and OCDM are consistent with the observed cluster abundance at $`z0`$, see Fig. 20, while Standard CDM (Einstein-De Sitter model, with $`\mathrm{\Omega }_\mathrm{M}=1`$), when normalized at COBE scales, produces too many clusters at all redshifts. The evolution of the cluster abundance with redshift breaks the degeneracy among the models at $`z0`$. The low-mass models (Open and $`\mathrm{\Lambda }`$-CDM) predict a relatively small change in the number density of rich clusters as a function of redshift because, due to the low density, hardly any structure growth occurs since $`z1`$. The high-mass models (Tilted and Standard CDM) predict that structure has grown steadily and rich clusters only formed recently: the number density of rich clusters at $`z1`$ is predicted to be exponentially smaller than today. The observation of a single massive cluster is enough to rule out the $`\mathrm{\Omega }_\mathrm{M}=1`$ model. In fact, three clusters have been seen, suggesting a low density universe , $$\mathrm{\Omega }_\mathrm{M}=0.25_{0.10}^{+0.15}(1\sigma \mathrm{statistical})\pm 20\%(\mathrm{systematic}).$$ (89) But one should be cautious. There is the caveat that for this constraint it is assumed that the initial spectrum of density perturbations is Gaussian, as predicted in the simplest models of inflation, but that has not yet been confirmed observationally on cluster scales. #### 3.2.9 Summary of the matter content We can summarize the present situation with Fig. 21, for $`\mathrm{\Omega }_\mathrm{M}`$ as a function of $`H_0`$. There are four bands, the luminous matter $`\mathrm{\Omega }_{\mathrm{lum}}`$; the baryon content $`\mathrm{\Omega }_\mathrm{B}`$, from BBN; the galactic halo component $`\mathrm{\Omega }_{\mathrm{halo}}`$, and the dynamical mass from clusters, $`\mathrm{\Omega }_\mathrm{M}`$. From this figure it is clear that there are in fact three dark matter problems: The first one is where are 90% of the baryons. Between the fraction predicted by BBN and that seen in stars and diffuse gas there is a huge fraction which is in the form of dark baryons. They could be in small clumps of hydrogen that have not started thermonuclear reactions and perhaps constitute the dark matter of spiral galaxies’ halos. Note that although $`\mathrm{\Omega }_\mathrm{B}`$ and $`\mathrm{\Omega }_{\mathrm{halo}}`$ coincide at $`H_070`$ km/s/Mpc, this could be just a coincidence. The second problem is what constitutes 90% of matter, from BBN baryons to the mass inferred from cluster dynamics. This is the standard dark matter problem and could be solved by direct detection of a weakly interacting massive particle in the laboratory. And finally, since we know from observations of the CMB, see Section 4.4, that the universe is flat, what constitutes around 60% of the energy density, from dynamical mass to critical density, $`\mathrm{\Omega }_0=1`$? One possibility could be that the universe is dominated by a diffuse vacuum energy, i.e. a cosmological constant, which only affects the very large scales. Alternatively, the theory of gravity (general relativity) may need to be modified on large scales, e.g. due to quantum gravity effects. The need to introduce an effective cosmological constant on large scales is nowadays the only reason why gravity may need to be modified at the quantum level. Since we still do not have a quantum theory of gravity, such a proposal is still very speculative, and most of the approaches simply consider the inclusion of a cosmological constant as a phenomenological parameter. #### 3.2.10 Massive neutrinos One of the ‘usual suspects’ when addressing the problem of dark matter are neutrinos. They are the only candidates known to exist. If neutrinos have a mass, could they constitute the missing matter? We know from the Big Bang theory, see Section 2.2.2, that there is a cosmic neutrino background at a temperature of approximately 2K. This allows one to compute the present number density in the form of neutrinos, which turns out to be, for massless neutrinos, $`n_\nu (T_\nu )=\frac{3}{11}n_\gamma (T_\gamma )=112\mathrm{cm}^3`$, per species of neutrino. If neutrinos have mass, as recent experiments seem to suggest, see Fig. 22, the cosmic energy density in massive neutrinos would be $`\rho _\nu =n_\nu m_\nu =\frac{3}{11}n_\gamma m_\nu `$, and therefore its contribution today, $$\mathrm{\Omega }_\nu h^2=\frac{m_\nu }{94\mathrm{eV}}.$$ (90) The discussion in the previous Sections suggest that $`\mathrm{\Omega }_\mathrm{M}0.4`$, and thus, for any of the three families of neutrinos, $`m_\nu 40`$ eV. Note that this limit improves by six orders of magnitude the present bound on the tau-neutrino mass . Supposing that the missing mass in non-baryonic cold dark matter arises from a single particle dark matter (PDM) component, its contribution to the critical density is bounded by $`0.05\mathrm{\Omega }_{\mathrm{PDM}}h^20.4`$, see Fig. 21. I will now go through the various logical arguments that exclude neutrinos as the dominant component of the missing dark matter in the universe. Is it possible that neutrinos with a mass $`4\mathrm{eV}m_\nu 40`$ eV be the non-baryonic PDM component? For instance, could massive neutrinos constitute the dark matter halos of galaxies? For neutrinos to be gravitationally bound to galaxies it is necessary that their velocity be less that the escape velocity $`v_{\mathrm{esc}}`$, and thus their maximum momentum is $`p_{\mathrm{max}}=m_\nu v_{\mathrm{esc}}`$. How many neutrinos can be packed in the halo of a galaxy? Due to the Pauli exclusion principle, the maximum number density is given by that of a completely degenerate Fermi gas with momentum $`p_\mathrm{F}=p_{\mathrm{max}}`$, i.e. $`n_{\mathrm{max}}=p_{\mathrm{max}}^3/3\pi ^2`$. Therefore, the maximum local density in dark matter neutrinos is $`\rho _{\mathrm{max}}=n_{\mathrm{max}}m_\nu =m_\nu ^4v_{\mathrm{esc}}^3/3\pi ^2`$, which must be greater than the typical halo density $`\rho _{\mathrm{halo}}=0.3`$ GeV cm<sup>-3</sup>. For a typical spiral galaxy, this constraint, known as the Tremaine-Gunn limit, gives $`m_\nu 40`$ eV, see Ref. . However, this mass, even for a single species, say the tau-neutrino, gives a value for $`\mathrm{\Omega }_\nu h^2=0.5`$, which is far too high for structure formation. Neutrinos of such a low mass would constitute a relativistic hot dark matter component, which would wash-out structure below the supercluster scale, against evidence from present observations, see Fig. 22. Furthermore, applying the same phase-space argument to the neutrinos as dark matter in the halo of dwarf galaxies gives $`m_\nu 100`$ eV, beyond closure density (90). We must conclude that the simple idea that light neutrinos could constitute the particle dark matter on all scales is ruled out. They could, however, still play a role as a sub-dominant hot dark matter component in a flat CDM model. In that case, a neutrino mass of order 1 eV is not cosmological excluded, see Fig. 22. Another possibility is that neutrinos have a large mass, of order a few GeV. In that case, their number density at decoupling, see Section 2.2.2, is suppressed by a Boltzmann factor, $`\mathrm{exp}(m_\nu /T_{\mathrm{dec}})`$. For masses $`m_\nu >T_{\mathrm{dec}}0.8`$ MeV, the present energy density has to be computed as a solution of the corresponding Boltzmann equation. Apart from a logarithmic correction, one finds $`\mathrm{\Omega }_\nu h^20.1(10\mathrm{GeV}/m_\nu )^2`$ for Majorana neutrinos and slightly smaller for Dirac neutrinos. In either case, neutrinos could be the dark matter only if their mass was a few GeV. Laboratory limits for $`\nu _\tau `$ of around 18 MeV , and much more stringent ones for $`\nu _\mu `$ and $`\nu _e`$, exclude the known light neutrinos. However, there is always the possibility of a fourth unknown heavy and stable (perhaps sterile) neutrino. If it couples to the Z boson and has a mass below 45 GeV for Dirac neutrinos (39.5 GeV for Majorana neutrinos), then it is ruled out by measurements at LEP of the invisible width of the Z. There are two logical alternatives, either it is a sterile neutrino (it does not couple to the Z), or it does couple but has a larger mass. In the case of a Majorana neutrino (its own antiparticle), their abundance, for this mass range, is too small for being cosmologically relevant, $`\mathrm{\Omega }_\nu h^20.005`$. If it were a Dirac neutrino there could be a lepton asymmetry, which may provide a higher abundance (similar to the case of baryogenesis). However, neutrinos scatter on nucleons via the weak axial-vector current (spin-dependent) interaction. For the small momentum transfers imparted by galactic WIMPs, such collisions are essentially coherent over an entire nucleus, leading to an enhancement of the effective cross section. The relatively large detection rate in this case allowes one to exclude fourth-generation Dirac neutrinos for the galactic dark matter . Anyway, it would be very implausible to have such a massive neutrino today, since it would have to be stable, with a life-time greater than the age of the universe, and there is no theoretical reason to expect a massive sterile neutrino that does not oscillate into the other neutrinos. Of course, the definitive test to the possible contribution of neutrinos to the overall density of the universe would be to measure directly their mass in laboratory experiments.<sup>13</sup><sup>13</sup>13For a review of Neutrinos, see Bilenky’s contribution to these Proceedings . There are at present two types of experiments: neutrino oscillation experiments, which measure only differences in squared masses, and direct mass-searches experiments, like the tritium $`\beta `$-spectrum and the neutrinoless double-$`\beta `$ decay experiments, which measure directly the mass of the electron neutrino and give a bound $`m_{\nu _e}<`$ 2 eV. Neutrinos with such a mass could very well constitute the HDM component of the universe, $`\mathrm{\Omega }_{\mathrm{HDM}}<0.15`$. The oscillation experiments give a variety of possibilities for $`\mathrm{\Delta }m_\nu ^2=0.33\mathrm{eV}^2`$ from LSND (not yet confirmed), to the atmospheric neutrino oscillations from SuperKamiokande ($`\mathrm{\Delta }m_\nu ^23\times 10^3\mathrm{eV}^2`$) and the solar neutrino oscillations ($`\mathrm{\Delta }m_\nu ^210^5\mathrm{eV}^2`$). Only the first two possibilities would be cosmologically relevant, see Fig. 22. #### 3.2.11 Weakly Interacting Massive Particles Unless we drastically change the theory of gravity on large scales, baryons cannot make up the bulk of the dark matter. Massive neutrinos are the only alternative among the known particles, but they are essentially ruled out as a universal dark matter candidate, even if they may play a subdominant role as a hot dark matter component. There remains the mystery of what is the physical nature of the dominant cold dark matter component. Something like a heavy stable neutrino, a generic Weakly Interacting Massive Particle (WIMP), could be a reasonable candidate because its present abundance could fall within the expected range, $$\mathrm{\Omega }_{\mathrm{PDM}}h^2\frac{G^{3/2}T_0^3h^2}{H_0^2\sigma _{\mathrm{ann}}v_{\mathrm{rel}}}=\frac{3\times 10^{27}\mathrm{cm}^3\mathrm{s}^1}{\sigma _{\mathrm{ann}}v_{\mathrm{rel}}}.$$ (91) Here $`v_{\mathrm{rel}}`$ is the relative velocity of the two incoming dark matter particles and the brackets $`\mathrm{}`$ denote a thermal average at the freeze-out temperature, $`T_\mathrm{f}m_{\mathrm{PDM}}/20`$, when the dark matter particles go out of equilibrium with radiation. The value of $`\sigma _{\mathrm{ann}}v_{\mathrm{rel}}`$ needed for $`\mathrm{\Omega }_{\mathrm{PDM}}1`$ is remarkably close to what one would expect for a WIMP with a mass $`m_{\mathrm{PDM}}=100`$ GeV, $`\sigma _{\mathrm{ann}}v_{\mathrm{rel}}\alpha ^2/8\pi m_{\mathrm{PDM}}3\times 10^{27}\mathrm{cm}^3\mathrm{s}^1`$. We still do not know whether this is just a coincidence or an important hint on the nature of dark matter. There are a few theoretical candidates for WIMPs, like the neutralino, coming from supersymmetric extensions of the standard model of particle physics,<sup>14</sup><sup>14</sup>14For a review of Supersymmetry (SUSY), see Carena’s contribution to these Proceedings. but at present there is no empirical evidence that such extensions are indeed realized in nature. In fact, the non-observation of supersymmetric particles at current accelerators places stringent limits on the neutralino mass and interaction cross section . If WIMPs constitute the dominant component of the halo of our galaxy, it is expected that some may cross the Earth at a reasonable rate to be detected. The direct experimental search for them rely on elastic WIMP collisions with the nuclei of a suitable target. Dark matter WIMPs move at a typical galactic virial velocity of around $`200300`$ km/s, depending on the model. If their mass is in the range $`10100`$ GeV, the recoil energy of the nuclei in the elastic collision would be of order 10 keV. Therefore, one should be able to identify such energy depositions in a macroscopic sample of the target. There are at present three different methods: First, one could search for scintillation light in NaI crystals or in liquid xenon; second, search for an ionization signal in a semiconductor, typically a very pure germanium crystal; and third, use a cryogenic detector at 10 mK and search for a measurable temperature increase of the sample. The main problem with such a type of experiment is the low expected signal rate, with a typical number below 1 event/kg/day. To reduce natural radioactive contamination one must use extremely pure substances, and to reduce the background caused by cosmic rays requires that these experiments be located deeply underground. The best limits on WIMP scattering cross sections come from some germanium experiments , as well as from the NaI scintillation detectors of the UK dark matter collaboration (UKDMC) in the Boulby salt mine in England , and the DAMA experiment in the Gran Sasso laboratory in Italy . Current experiments already touch the parameter space expected from supersymmetric particles, see Fig. 23, and therefore there is a chance that they actually discover the nature of the missing dark matter. The problem, of course, is to attribute a tentative signal unambiguously to galactic WIMPs rather than to some unidentified radioactive background. One specific signature is the annual modulation which arises as the Earth moves around the Sun.<sup>15</sup><sup>15</sup>15The time scale of the Sun’s orbit around the center of the galaxy is too large to be relevant in the analysis. Therefore, the net speed of the Earth relative to the galactic dark matter halo varies, causing a modulation of the expected counting rate. The DAMA/NaI experiment has actually reported such a modulation signal, see Fig. 24, from the combined analysis of their 4-year data , which provides a confidence level of 99.6% for a neutralino mass of $`m_\chi =52_8^{+10}`$ GeV and a proton cross section of $`\xi \sigma _p=7.2_{0.9}^{+0.4}\times 10^6`$ pb, where $`\xi =\rho _\chi /0.3`$ GeV cm<sup>-3</sup> is the local neutralino energy density in units of the galactic halo density. There has been no confirmation yet of this result from other dark matter search groups, but hopefully in the near future we will have much better sensitivity at low masses from the Cryogenic Rare Event Search with Superconducting Thermometers (CRESST) experiment at Gran Sasso as well as at weaker cross sections from the CDMS experiment at Stanford and the Soudan mine, see Fig. 25. The CRESST experiment uses sapphire crystals as targets and a new method to simultaneously measure the phonons and the scintillating light from particle interactions inside the crystal, which allows excellent background discrimination. Very recently there has been the interesting proposal of a completely new method based on a Superheated Droplet Detector (SDD), which claims to have already a similar sensitivity as the more standard methods described above, see Ref. . There exist other indirect methods to search for galactic WIMPs . Such particles could self-annihilate at a certain rate in the galactic halo, producing a potentially detectable background of high energy photons or antiprotons. The absence of such a background in both gamma ray satellites and the Alpha Matter Spectrometer imposes bounds on their density in the halo. Alternatively, WIMPs traversing the solar system may interact with the matter that makes up the Earth or the Sun so that a small fraction of them will lose energy and be trapped in their cores, building up over the age of the universe. Their annihilation in the core would thus produce high energy neutrinos from the center of the Earth or from the Sun which are detectable by neutrino telescopes. In fact, SuperKamiokande already covers a large part of SUSY parameter space. In other words, neutrino telescopes are already competitive with direct search experiments. In particular, the AMANDA experiment at the South Pole , which is expected to have $`10^3`$ Cherenkov detectors 2.3 km deep in very clear ice, over a volume $`1`$ km<sup>3</sup>, is competitive with the best direct searches proposed. The advantages of AMANDA are also directional, since the arrays of Cherenkov detectors will allow one to reconstruct the neutrino trajectory and thus its source, whether it comes from the Earth or the Sun. ### 3.3 The cosmological constant $`\mathrm{\Omega }_\mathrm{\Lambda }`$ A cosmological constant is a term in the Einstein equations, see Eq. (1), that corresponds to the energy density of the vacuum of quantum field theories, $`\mathrm{\Lambda }8\pi G\rho _v`$, see Ref. . These theories predict a value of order $`\rho _vM_\mathrm{P}^45\times 10^{93}`$ g/cm<sup>3</sup>, which is about 123 orders of magnitude larger than the critical density (14). Such a discrepancy is one of the biggest problems of theoretical physics . It has always been assumed that quantum gravity effects, via some as yet unknown symmetry, would exactly cancel the cosmological constant, but this remains a downright speculation. Moreover, one of the difficulties with a non-zero value for $`\mathrm{\Lambda }`$ is that it appears coincidental that we are now living at a special epoch when the cosmological constant starts to dominate the dynamics of the universe, and that it will do so forever after, see Section 2.1.2 and Eq. (22). Nevertheless, ever since Einstein introduced it in 1917, this ethereal constant has been invoked several times in history to explain a number of apparent crises, always to disappear under further scrutiny . In spite of the theoretical prejudice towards $`\mathrm{\Lambda }=0`$, there are new observational arguments for a non-zero value. The most compelling ones are recent evidence that we live in a flat universe, from observations of CMB anisotropies, together with strong indications of a low mass density universe ($`\mathrm{\Omega }_\mathrm{M}<1`$), from the large scale distribution of galaxies, clusters and voids, that indicate that some kind of dark energy must make up the rest of the energy density up to critical, i.e. $`\mathrm{\Omega }_\mathrm{\Lambda }=1\mathrm{\Omega }_\mathrm{M}`$. In addition, the discrepancy between the ages of globular clusters and the expansion age of the universe may be cleanly resolved with $`\mathrm{\Lambda }0`$. Finally, there is growing evidence for an accelerating universe from observations of distant supernovae. I will now discuss the different arguments one by one. The only known way to reconcile a low mass density with a flat universe is if an additional “dark” energy dominates the universe today. It would have to resist gravitational collapse, otherwise it would have been detected already as part of the energy in the halos of galaxies. However, if most of the energy of the universe resists gravitational collapse, it is impossible for structure in the universe to grow. This dilemma can be resolved if the hypothetical dark energy was negligible in the past and only recently became the dominant component. According to general relativity, this requires that the dark energy have negative pressure, since the ratio of dark energy to matter density goes like $`a(t)^{3p/\rho }`$. This argument would rule out almost all of the usual suspects, such as cold dark matter, neutrinos, radiation, and kinetic energy, since they all have zero or positive pressure. Thus, we expect something like a cosmological constant, with negative pressure, $`p\rho `$, to account for the missing energy. This negative pressure would help accelerate the universe and reconcile the expansion age of the universe with the ages of stars in globular clusters, see Fig. 11, where $`t_0H_0`$ is shown as a function of $`\mathrm{\Omega }_\mathrm{M}`$, in a flat universe, $`\mathrm{\Omega }_\mathrm{\Lambda }=1\mathrm{\Omega }_\mathrm{M}`$, and an open one, $`\mathrm{\Omega }_\mathrm{\Lambda }=0`$. For the present age of the universe of $`t_0=13\pm 1`$ Gyr, and the measured rate of expansion, $`H_0=70\pm 7`$ km/s/Mpc, one finds $`t_0H_0=0.93\pm 0.12`$ (adding errors in quadrature), which corresponds to $`\mathrm{\Omega }_\mathrm{M}=0.05_{0.10}^{+0.24}`$ for an open universe, see Fig. 11, marginally consistent with observations of large scale structure. On the other hand, for a flat universe with a cosmological constant, $`t_0H_0=0.93\pm 0.12`$ corresponds to $`\mathrm{\Omega }_\mathrm{M}=0.34_{0.12}^{+0.20}`$, which is perfectly compatible with recent observations. These suggest that we probably live in a flat universe that is accelerating, dominated today by a vacuum energy density. This conclusions have been supported by growingly robust observational evidence from distant supernovae. In their quest for the cosmological parameters, astronomers look for distant astrophysical objects that can serve as standard candles to determine the distance to the object from their observed apparent luminosity. A candidate that has recently been exploited with great success is a certain type of supernova explosions at large redshifts, called SN of type Ia. These are white dwarf stars at the end of their life cycle that accrete matter from a companion until they become unstable and violently explode in a natural thermonuclear explosion that out-shines their progenitor galaxy. The intensity of the distant flash varies in time, it takes about three weeks to reach its maximum brightness and then it declines over a period of months. Although the maximum luminosity varies from one supernova to another, depending on their original mass, their environment, etc., there is a pattern: brighter explosions last longer than fainter ones. By studying the characteristic light curves, see Fig. 26, of a reasonably large statistical sample, cosmologists from two competing groups, the Supernova Cosmology Project and the High-redshift Supernova Project , are confident that they can use this type of supernova as a standard candle. Since the light coming from some of these rare explosions has travelled for a large fraction of the size of the universe, one expects to be able to infer from their distribution the spatial curvature and the rate of expansion of the universe. One of the surprises revealed by these observations is that high redshift type Ia supernovae appear fainter than expected for either an open ($`\mathrm{\Omega }_\mathrm{M}<1`$) or a flat ($`\mathrm{\Omega }_\mathrm{M}=1`$) universe, see Fig. 27. In fact, the universe appears to be accelerating instead of decelerating, as was expected from the general attraction of matter, see Eq. (26); something seems to be acting as a repulsive force on very large scales. The most natural explanation for this is the presence of a cosmological constant, a diffuse vacuum energy that permeates all space and, as explained above, gives the universe an acceleration that tends to separate gravitationally bound systems from each other. The best-fit results from the Supernova Cosmology Project give a linear combination $`0.8\mathrm{\Omega }_\mathrm{M}0.6\mathrm{\Omega }_\mathrm{\Lambda }=0.2\pm 0.1`$ $`(1\sigma )`$, and, for a flat universe ($`\mathrm{\Omega }_\mathrm{M}+\mathrm{\Omega }_\mathrm{\Lambda }=1`$), the best-fit values for the combined analysis of both groups , are $`\mathrm{\Omega }_\mathrm{M}^{\mathrm{flat}}`$ $`=`$ $`0.28_{0.08}^{+0.09}(1\sigma \mathrm{statistical})_{0.04}^{+0.05}(\mathrm{identified}\mathrm{systematics}),`$ (92) $`\mathrm{\Omega }_\mathrm{\Lambda }^{\mathrm{flat}}`$ $`=`$ $`0.72_{0.09}^{+0.08}(1\sigma \mathrm{statistical})_{0.05}^{+0.04}(\mathrm{identified}\mathrm{systematics}).`$ (93) However, one may think that it is still premature to conclude that the universe is indeed accelerating, because of possibly large systematic errors inherent to most cosmological measurements, and in particular to observations of supernovae at large redshifts. There has been attempts to find crucial systematic effects like evolution, chemical composition dependence, reddening by dust, etc. in the supernovae observations that would invalidate the claims, but none of them are now considered as a serious threat. Perhaps the most critical one today seems to be sampling effects, since the luminosities of the high-redshift supernovae ($`z0.51.0`$) are all measured relative to the same set of local supernovae ($`z<0.3`$). Hence, absolute calibrations, completeness levels, and any other systematic effects related to both data sets are critical. For instance, the intense efforts to search for high-redshift objects have led to the peculiar situation where the nearby sample, which is used for calibration, is now smaller than the distant one. Further searches, already underway, for increasing the nearby supernovae sample will provide an important check. Moreover, there are bounds on a cosmological constant that come from the statistics of gravitational lensing, with two different methods. Gravitational lensing can be due to various accumulations of matter along the line of sight to the distant light sources. The first method uses the abundance of multiply imaged sources like quasars, lensed by intervening galaxies. The probability of finding a lensed image is directly proportional to the number of galaxies (lenses) along the path and thus to the distance to the source. This distance, for fixed $`H_0`$, increases dramatically for a large value of the cosmological constant: the age of the universe and the distance to the galaxy become large for $`\mathrm{\Omega }_\mathrm{\Lambda }0`$ because the universe has been expanding for a longer time; therefore, more lenses are predicted for $`\mathrm{\Omega }_\mathrm{\Lambda }>0`$. Using this method, an upper limit of $$\mathrm{\Omega }_\mathrm{\Lambda }<0.75(95\%\mathrm{c}.\mathrm{l}.)$$ (94) has recently been obtained , marginally consistent with the supernovae results, but there are caveats to this powerful method due to uncertainties in the number density and lensing cross section of the lensing galaxies as well as the distant quasars. A second method is lensing by massive clusters of galaxies, which produces widely separated lensed images of quasars and distorted images of background galaxies. The observed statistics, when compared with numerical simulations, rule out the $`\mathrm{\Omega }_\mathrm{M}=1`$ models and set an upper bound on the cosmological constant, $`\mathrm{\Omega }_\mathrm{\Lambda }<0.7`$, see Ref. . However, this limit is very sensitive to the resolution of the numerical simulations, which are currently improving. ### 3.4 The spatial curvature $`\mathrm{\Omega }_K`$ As we will discuss in detail in Section 4.4, observations of the two-point correlation function of temperature anisotropies in the microwave background provide a crucial test for the spatial curvature of the universe. From those observations one can tell whether the photons that left the last scattering surface, at redshift $`z=1100`$, have travelled in straight lines, like in a flat universe, or in curved paths, like in an open one. Very recent observations made by the balloon experiment BOOMERANG suggest that the universe is indeed spatially flat ($`\mathrm{\Omega }_K=0`$) with about 10% accuracy , $$\mathrm{\Omega }_0=\mathrm{\Omega }_\mathrm{M}+\mathrm{\Omega }_\mathrm{\Lambda }=1.0\pm 0.1(95\%\mathrm{c}.\mathrm{l}.)$$ (95) These measuremnts are bound to be improved in the near future, by both balloon experiments and by the Microwave Anisotropy Probe (MAP) satellite, to be launched by NASA at the end of year 2000 . Furthermore, with the launch in 2007 of Planck satellite we will be able to determine $`\mathrm{\Omega }_0`$ with 1% accuracy. ### 3.5 The age of the universe $`t_0`$ The universe must be older than the oldest objects it contains. Those are believed to be the stars in the oldest clusters in the Milky Way, globular clusters. The most reliable ages come from the application of theoretical models of stellar evolution to observations of old stars in globular clusters. For about 30 years, the ages of globular clusters have remained reasonable stable, at about 15 Gyr . However, recently these ages have been revised downward . During the 1980s and 1990s, the globular cluster age estimates have improved as both new observations have been made with CCDs, and since refinements to stellar evolution models, including opacities, consideration of mixing, and different chemical abundances have been incorporated . From the theory side, uncertainties in globular cluster ages come from uncertainties in convection models, opacities, and nuclear reaction rates. From the observational side, uncertainties arise due to corrections for dust and chemical composition. However, the dominant source of systematic errors in the globular cluster age is the uncertainty in the cluster distances. Fortunately, the Hipparcos satellite recently provided geometric parallax measurements for many nearby old stars with low metallicity, typical of glubular clusters, thus allowing for a new calibration of the ages of stars in globular clusters, leading to a downward revision to $`1013`$ Gyr . Moreover, there were very few stars in the Hipparcos catalog with both small parallax erros and low metal abundance. Hence, an increase in the sample size could be critical in reducing the statatistical uncertaintites for the calibration of the globular cluster ages. There are already proposed two new parallax satellites, NASA’s Space Interferometry Mission (SIM) and ESA’s mission, called GAIA, that will give 2 or 3 orders of magnitude more accurate parallaxes than Hipparcos, down to fainter magnitude limits, for several orders of magnitude more stars. Until larger samples are available, however, distance errors are likely to be the largest source of systematic uncertainty to the globular cluster age . The supernovae groups can also determine the age of the universe from their high redshift observations. Figure 28 shows that the confidence regions in the $`(\mathrm{\Omega }_\mathrm{M},\mathrm{\Omega }_\mathrm{\Lambda })`$ plane are almost parallel to the contours of constant age. For any value of the Hubble constant less than $`H_0=70`$ km/s/Mpc, the implied age of the universe is greater than 13 Gyr, allowing enough time for the oldest stars in globular clusters to evolve . Integrating over $`\mathrm{\Omega }_\mathrm{M}`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }`$, the best fit value of the age in Hubble-time units is $`H_0t_0=0.93\pm 0.06`$ or equivalently $`t_0=14.1\pm 1.0(0.65h^1)`$ Gyr . The age would be somewhat larger in a flat universe: $`H_0t_0^{\mathrm{flat}}=0.96_{0.07}^{+0.09}`$ or, equivalently, $$t_0^{\mathrm{flat}}=14.4_{1.1}^{+1.4}(0.65h^1)\mathrm{Gyr}.$$ (96) Furthermore, a combination of 8 independent recent measurements: CMB anisotropies, type Ia SNe, cluster mass-to-light ratios, cluster abundance evolution, cluster baryon fraction, deuterium-to-hidrogen ratios in quasar spectra, double-lobed radio sources and the Hubble constant, can be used to determine the present age of the universe . The result is shown in Fig. 29, compared to other recent determinations. The best fit value for the age of the universe is, according to this analysis, $`t_0=13.4\pm 1.6`$ Gyr, about a billion years younger than other recent estimates . We can summarize this Section by showing the region in parameter space where we stand nowadays, thanks to the recent cosmological observations. We have plotted that region in Fig. 30. One could also superimpose the contour lines corresponding to equal $`t_0H_0`$ lines, as a cross check. It is extraordinary that only in the last few months we have been able to reduce the concordance region to where it stands today, where all the different observations seem to converge. There are still many uncertainties, mainly systematic; however, those are quickly decreasing and becoming predominantly statistical. In the near future, with precise observations of the anisotropies in the microwave background temperature and polarization, to be discussed in Section 4.4, we will be able to reduce those uncertainties to the level of one percent. This is the reason why cosmologists are so excited and why it is claimed that we live in the Golden Age of Cosmology. ## 4 THE INFLATIONARY PARADIGM The hot Big Bang theory is nowadays a very robust edifice, with many independent observational checks: the expansion of the universe; the abundance of light elements; the cosmic microwave background; a predicted age of the universe compatible with the age of the oldest objects in it, and the formation of structure via gravitational collapse of initially small inhomogeneities. Today, these observations are confirmed to within a few percent accuracy, and have helped establish the hot Big Bang as the preferred model of the universe. All the physics involved in the above observations is routinely tested in the laboratory (atomic and nuclear physics experiments) or in the solar system (general relativity). However, this theory leaves a range of crucial questions unanswered, most of which are initial conditions’ problems. There is the reasonable assumption that these cosmological problems will be solved or explained by new physical principles at high energies, in the early universe. This assumption leads to the natural conclusion that accurate observations of the present state of the universe may shed light onto processes and physical laws at energies above those reachable by particle accelerators, present or future. We will see that this is a very optimistic approach indeed, and that there are many unresolved issues related to those problems. However, there might be in the near future reasons to be optimistic. ### 4.1 Shortcomings of Big Bang Cosmology The Big Bang theory could not explain the origin of matter and structure in the universe; that is, the origin of the matter–antimatter asymmetry, without which the universe today would be filled by a uniform radiation continuosly expanding and cooling, with no traces of matter, and thus without the possibility to form gravitationally bound systems like galaxies, stars and planets that could sustain life. Moreover, the standard Big Bang theory assumes, but cannot explain, the origin of the extraordinary smoothness and flatness of the universe on the very large scales seen by the microwave background probes and the largest galaxy catalogs. It cannot explain the origin of the primordial density perturbations that gave rise to cosmic structures like galaxies, clusters and superclusters, via gravitational collapse; the quantity and nature of the dark matter that we believe holds the universe together; nor the origin of the Big Bang itself. A summary of the problems that the Big Bang theory cannot explain is: * The global structure of the universe. \- Why is the universe so close to spatial flatness? \- Why is matter so homogeneously distributed on large scales? * The origin of structure in the universe. \- How did the primordial spectrum of density perturbations originate? * The origin of matter and radiation. \- Where does all the energy in the universe come from? \- What is the nature of the dark matter in the universe? \- How did the matter-antimatter asymmetry arise? * The initial singularity. \- Did the universe have a beginning? \- What is the global structure of the universe beyond our observable patch? Let me discuss one by one the different issues: #### 4.1.1 The Flatness Problem The Big Bang theory assumes but cannot explain the extraordinary spatial flatness of our local patch of the universe. In the general FRW metric (2) the parameter $`K`$ that characterizes spatial curvature is a free parameter. There is nothing in the theory that determines this parameter a priori. However, it is directly related, via the Friedmann equation (8), to the dynamics, and thus the matter content, of the universe, $$K=\frac{8\pi G}{3}\rho a^2H^2a^2=\frac{8\pi G}{3}\rho a^2\left(\frac{\mathrm{\Omega }1}{\mathrm{\Omega }}\right).$$ (97) We can therefore define a new variable, $$x\frac{\mathrm{\Omega }1}{\mathrm{\Omega }}=\frac{\mathrm{const}.}{\rho a^2},$$ (98) whose time evolution is given by $$x^{}=\frac{dx}{dN}=(1+3\omega )x,$$ (99) where $`N=\mathrm{ln}(a/a_i)`$ characterizes the number of $`e`$-folds of universe expansion ($`dN=Hdt`$) and where we have used Eq. (38) for the time evolution of the total energy, $`\rho a^3`$, which only depends on the barotropic ratio $`\omega `$. It is clear from Eq. (99) that the phase-space diagram $`(x,x^{})`$ presents an unstable critical (saddle) point at $`x=0`$ for $`\omega >1/3`$, i.e. for the radiation ($`\omega =1/3`$) and matter ($`\omega =0`$) eras. A small perturbation from $`x=0`$ will drive the system towards $`x=\pm \mathrm{}`$. Since we know the universe went through both the radiation era (because of primordial nucleosynthesis) and the matter era (because of structure formation), tiny deviations from $`\mathrm{\Omega }=1`$ would have grown since then, such that today $$x_0=\frac{\mathrm{\Omega }_01}{\mathrm{\Omega }_0}=x_{\mathrm{in}}\left(\frac{T_{\mathrm{in}}}{T_{\mathrm{eq}}}\right)^2(1+z_{\mathrm{eq}}).$$ (100) In order that today’s value be in the range $`0.1<\mathrm{\Omega }_0<1.2`$, or $`x_0𝒪(1)`$, it is required that at, say, primordial nucleosynthesis ($`T_{_{\mathrm{NS}}}10^6T_{\mathrm{eq}}`$) its value be $$\mathrm{\Omega }(t_{_{\mathrm{NS}}})=1\pm 10^{15},$$ (101) which represents a tremendous finetuning. Perhaps the universe indeed started with such a peculiar initial condition, but it is epistemologically more satisfying if we give a fundamental dynamical reason for the universe to have started so close to spatial flatness. These arguments were first used by Robert Dicke in the 1960s, much before inflation. He argued that the most natural initial condition for the spatial curvature should have been the Planck scale curvature, $`{}_{}{}^{(3)}R=6K/l_\mathrm{P}^2`$, where the Planck length is $`l_\mathrm{P}=(\mathrm{}G/c^3)^{1/2}=1.62\times 10^{33}`$ cm, that is, 60 orders of magnitude smaller than the present size of the universe, $`a_0=1.38\times 10^{28}`$ cm. A universe with this immense curvature would have collapsed within a Planck time, $`t_\mathrm{P}=(\mathrm{}G/c^5)^{1/2}=5.39\times 10^{44}`$ s, again 60 orders of magnitude smaller than the present age of the universe, $`t_0=4.1\times 10^{17}`$ s. Therefore, the flatness problem is also related to the Age Problem, why is it that the universe is so old and flat when, under ordinary circumstances (based on the fundamental scale of gravity) it should have lasted only a Planck time and reached a size of order the Planck length? As we will see, inflation gives a dynamical reason to such a peculiar initial condition. #### 4.1.2 The Homogeneity Problem An expanding universe has particle horizons, that is, spatial regions beyond which causal communication cannot occur. The horizon distance can be defined as the maximum distance that light could have travelled since the origin of the universe , $$d_\mathrm{H}(t)a(t)_0^t\frac{dt^{}}{a(t^{})}H^1(t),$$ (102) which is proportional to the Hubble scale.<sup>16</sup><sup>16</sup>16For the radiation era, the horizon distance is equal to the Hubble scale. For the matter era it is twice the Hubble scale. For instance, at the beginning of nucleosynthesis the horizon distance is a few light-seconds, but grows linearly with time and by the end of nucleosynthesis it is a few light-minutes, i.e. a factor 100 larger, while the scale factor has increased only a factor of 10. The fact that the causal horizon increases faster, $`d_\mathrm{H}t`$, than the scale factor, $`at^{1/2}`$, implies that at any given time the universe contains regions within itself that, according to the Big Bang theory, were never in causal contact before. For instance, the number of causally disconnected regions at a given redshift $`z`$ present in our causal volume today, $`d_\mathrm{H}(t_0)a_0`$, is $$N_{\mathrm{CD}}(z)\left(\frac{a(t)}{d_\mathrm{H}(t)}\right)^3(1+z)^{3/2},$$ (103) which, for the time of decoupling, is of order $`N_{\mathrm{CD}}(z_{\mathrm{dec}})10^51`$. This phenomenon is particularly acute in the case of the observed microwave background. Information cannot travel faster than the speed of light, so the causal region at the time of photon decoupling could not be larger than $`d_\mathrm{H}(t_{\mathrm{dec}})3\times 10^5`$ light years across, or about $`1^{}`$ projected in the sky today. So why should regions that are separated by more than $`1^{}`$ in the sky today have exactly the same temperature, to within 10 ppm, when the photons that come from those two distant regions could not have been in causal contact when they were emitted? This constitutes the so-called horizon problem, see Fig. 31, and was first discussed by Robert Dicke in the 1970s as a profound inconsistency of the Big Bang theory. ### 4.2 Cosmological Inflation In the 1980s, a new paradigm, deeply rooted in fundamental physics, was put forward by Alan H. Guth , Andrei D. Linde and others , to address these fundamental questions. According to the inflationary paradigm, the early universe went through a period of exponential expansion, driven by the approximately constant energy density of a scalar field called the inflaton. In modern physics, elementary particles are represented by quantum fields, which resemble the familiar electric, magnetic and gravitational fields. A field is simply a function of space and time whose quantum oscillations are interpreted as particles. In our case, the inflaton field has, associated with it, a large potential energy density, which drives the exponential expansion during inflation, see Fig. 32. We know from general relativity that the density of matter determines the expansion of the universe, but a constant energy density acts in a very peculiar way: as a repulsive force that makes any two points in space separate at exponentially large speeds. (This does not violate the laws of causality because there is no information carried along in the expansion, it is simply the stretching of space-time.) This superluminal expansion is capable of explaining the large scale homogeneity of our observable universe and, in particular, why the microwave background looks so isotropic: regions separated today by more than $`1^{}`$ in the sky were, in fact, in causal contact before inflation, but were stretched to cosmological distances by the expansion. Any inhomogeneities present before the tremendous expansion would be washed out. This explains why photons from supposedly causally disconneted regions have actually the same spectral distribution with the same temperature, see Fig. 31. Moreover, in the usual Big Bang scenario a flat universe, one in which the gravitational attraction of matter is exactly balanced by the cosmic expansion, is unstable under perturbations: a small deviation from flatness is amplified and soon produces either an empty universe or a collapsed one. As we discussed above, for the universe to be nearly flat today, it must have been extremely flat at nucleosynthesis, deviations not exceeding more than one part in $`10^{15}`$. This extreme fine tuning of initial conditions was also solved by the inflationary paradigm, see Fig. 33. Thus inflation is an extremely elegant hypothesis that explains how a region much, much greater that our own observable universe could have become smooth and flat without recourse to ad hoc initial conditions. Furthermore, inflation dilutes away any “unwanted” relic species that could have remained from early universe phase transitions, like monopoles, cosmic strings, etc., which are predicted in grand unified theories and whose energy density could be so large that the universe would have become unstable, and collapsed, long ago. These relics are diluted by the superluminal expansion, which leaves at most one of these particles per causal horizon, making them harmless to the subsequent evolution of the universe. The only thing we know about this peculiar scalar field, the inflaton, is that it has a mass and a self-interaction potential $`V(\varphi )`$ but we ignore everything else, even the scale at which its dynamics determines the superluminal expansion. In particular, we still do not know the nature of the inflaton field itself, is it some new fundamental scalar field in the electroweak symmetry breaking sector, or is it just some effective description of a more fundamental high energy interaction? Hopefully, in the near future, experiments in particle physics might give us a clue to its nature. Inflation had its original inspiration in the Higgs field, the scalar field supposed to be responsible for the masses of elementary particles (quarks and leptons) and the breaking of the electroweak symmetry. Such a field has not been found yet, and its discovery at the future particle colliders would help understand one of the truly fundamental problems in physics, the origin of masses. If the experiments discover something completely new and unexpected, it would automatically affect the idea of inflation at a fundamental level. #### 4.2.1 Homogeneous scalar field dynamics In this subsection I will describe the theoretical basis for the phenomenon of inflation. Consider a scalar field $`\varphi `$, a singlet under any given interaction, with an effective potential $`V(\varphi )`$. The Lagrangian for such a field in a curved background is $$_{\mathrm{inf}}=\frac{1}{2}g^{\mu \nu }_\mu \varphi _\nu \varphi V(\varphi ),$$ (104) whose evolution equation in a Friedmann-Robertson-Walker metric (2) and for a homogeneous field $`\varphi (t)`$ is given by $$\ddot{\varphi }+3H\dot{\varphi }+V^{}(\varphi )=0,$$ (105) where $`H`$ is the rate of expansion, together with the Einstein equations, $`H^2`$ $`=`$ $`{\displaystyle \frac{\kappa ^2}{3}}\left({\displaystyle \frac{1}{2}}\dot{\varphi }^2+V(\varphi )\right),`$ (106) $`\dot{H}`$ $`=`$ $`{\displaystyle \frac{\kappa ^2}{2}}\dot{\varphi }^2,`$ (107) where $`\kappa ^28\pi G`$. The dynamics of inflation can be described as a perfect fluid (7) with a time dependent pressure and energy density given by $`\rho `$ $`=`$ $`{\displaystyle \frac{1}{2}}\dot{\varphi }^2+V(\varphi ),`$ (108) $`p`$ $`=`$ $`{\displaystyle \frac{1}{2}}\dot{\varphi }^2V(\varphi ).`$ (109) The field evolution equation (105) can then be written as the energy conservation equation, $$\dot{\rho }+3H(\rho +p)=0.$$ (110) If the potential energy density of the scalar field dominates the kinetic energy, $`V(\varphi )\dot{\varphi }^2`$, then we see that $$p\rho \rho \mathrm{const}.H(\varphi )\mathrm{const}.,$$ (111) which leads to the solution $$a(t)\mathrm{exp}(Ht)\frac{\ddot{a}}{a}>0\mathrm{accelerated}\mathrm{expansion}.$$ (112) Using the definition of the number of $`e`$-folds, $`N=\mathrm{ln}(a/a_i)`$, we see that the scale factor grows exponentially, $`a(N)=a_i\mathrm{exp}(N)`$. This solution of the Einstein equations solves immediately the flatness problem. Recall that the problem with the radiation and matter eras is that $`\mathrm{\Omega }=1`$ ($`x=0`$) is an unstable critical point in phase-space. However, during inflation, with $`p\rho \omega 1`$, we have that $`1+3\omega 0`$ and therefore $`x=0`$ is a stable attractor of the equations of motion, see Eq. (99). As a consequence, what seemed an ad hoc initial condition, becomes a natural prediction of inflation. Suppose that during inflation the scale factor increased $`N`$ $`e`$-folds, then $$x_0=x_{\mathrm{in}}e^{2N}\left(\frac{T_{\mathrm{rh}}}{T_{\mathrm{eq}}}\right)^2(1+z_{\mathrm{eq}})e^{2N}\mathrm{\hspace{0.17em}10}^{56}1N65,$$ (113) where we have assumed that inflation ended at the scale $`V_{\mathrm{end}}`$, and the transfer of the inflaton energy density to thermal radiation at reheating occurred almost instantaneously<sup>17</sup><sup>17</sup>17There could be a small delay in thermalization, due to the intrinsic inefficiency of reheating, but this does not change significantly the required number of $`e`$-folds. at the temperature $`T_{\mathrm{rh}}V_{\mathrm{end}}^{1/4}10^{15}`$ GeV. Note that we can now have initial conditions with a large uncertainty, $`x_{\mathrm{in}}1`$, and still have today $`x_01`$, thanks to the inflationary attractor towards $`\mathrm{\Omega }=1`$. This can be understood very easily by realizing that the three curvature evolves during inflation as $$^{(3)}R=\frac{6K}{a^2}=^{(3)}R_{\mathrm{in}}e^{2N}0,\mathrm{for}N1.$$ (114) Therefore, if cosmological inflation lasted over 65 $`e`$-folds, as most models predict, then today the universe (or at least our local patch) should be exactly flat, see Fig. 33, a prediction that can be tested with great accuracy in the near future and for which already seems to be some evidence from observations of the microwave background . Furthermore, inflation also solves the homogeneity problem in a spectacular way. First of all, due to the superluminal expansion, any inhomogeneity existing prior to inflation will be washed out, $$\delta _k\left(\frac{k}{aH}\right)^2\mathrm{\Phi }_ke^{2N}0,\mathrm{for}N1.$$ (115) Moreover, since the scale factor grows exponentially, while the horizon distance remains essentially constant, $`d_H(t)H^1=`$ const., any scale within the horizon during inflation will be stretched by the superluminal expansion to enormous distances, in such a way that at photon decoupling all the causally disconnected regions that encompass our present horizon actually come from a single region during inflation, about 65 $`e`$-folds before the end. This is the reason why two points separated more than $`1^{}`$ in the sky have the same backbody temperature, as observed by the COBE satellite: they were actually in causal contact during inflation. There is at present no other proposal known that could solve the homogeneity problem without invoquing an acausal mechanism like inflation. Finally, any relic particle species (relativistic or not) existing prior to inflation will be diluted by the expansion, $`\rho _\mathrm{M}`$ $``$ $`a^3e^{3N}0,\mathrm{for}N1,`$ (116) $`\rho _\mathrm{R}`$ $``$ $`a^4e^{4N}0,\mathrm{for}N1.`$ (117) Note that the vacuum energy density $`\rho _v`$ remains constant under the expansion, and therefore, very soon it is the only energy density remaining to drive the expansion of the universe. #### 4.2.2 The slow-roll approximation In order to simplify the evolution equations during inflation, we will consider the slow-roll approximation (SRA). Suppose that, during inflation, the scalar field evolves very slowly down its effective potential, then we can define the slow-roll parameters , $`ϵ`$ $``$ $`{\displaystyle \frac{\dot{H}}{H^2}}={\displaystyle \frac{\kappa ^2}{2}}{\displaystyle \frac{\dot{\varphi }^2}{H^2}}1,`$ (118) $`\delta `$ $``$ $`{\displaystyle \frac{\ddot{\varphi }}{H\dot{\varphi }}}1.`$ (119) It is easy to see that the condition $$ϵ<1\frac{\ddot{a}}{a}>0$$ (120) characterizes inflation: it is all you need for superluminal expansion, i.e. for the horizon distance to grow more slowly than the scale factor, in order to solve the homogeneity problem, as well as for the spatial curvature to decay faster than usual, in order to solve the flatness problem. The number of $`e`$-folds during inflation can be written with the help of Eq. (118) as $$N=\mathrm{ln}\frac{a_{\mathrm{end}}}{a_i}=_{t_i}^{t_e}H𝑑t=_{\varphi _i}^{\varphi _e}\frac{\kappa d\varphi }{\sqrt{2ϵ(\varphi )}},$$ (121) which is an exact expression in terms of $`ϵ(\varphi )`$. In the limit given by Eqs. (118), the evolution equations (105) and (106) become $`H^2\left(1{\displaystyle \frac{ϵ}{3}}\right)`$ $``$ $`H^2={\displaystyle \frac{\kappa ^2}{3}}V(\varphi ),`$ (122) $`3H\dot{\varphi }\left(1{\displaystyle \frac{\delta }{3}}\right)`$ $``$ $`3H\dot{\varphi }=V^{}(\varphi ).`$ (123) Note that this corresponds to a reduction of the dimensionality of phase-space from two to one dimensions, $`H(\varphi ,\dot{\varphi })H(\varphi )`$. In fact, it is possible to prove a theorem, for single-field inflation, which states that the slow-roll approximation is an attractor of the equations of motion, and thus we can always evaluate the inflationary trajectory in phase-space within the SRA, therefore reducing the number of initial conditions to just one, the initial value of the scalar field. If $`H(\varphi )`$ only depends on $`\varphi `$, then $`H^{}(\varphi )=\kappa ^2\dot{\varphi }/2`$ and we can rewrite the slow-roll parameters (118) as $`ϵ`$ $`=`$ $`{\displaystyle \frac{2}{\kappa ^2}}\left({\displaystyle \frac{H^{}(\varphi )}{H(\varphi )}}\right)^2{\displaystyle \frac{1}{2\kappa ^2}}\left({\displaystyle \frac{V^{}(\varphi )}{V(\varphi )}}\right)^21,`$ (124) $`\delta `$ $`=`$ $`{\displaystyle \frac{2}{\kappa ^2}}{\displaystyle \frac{H^{\prime \prime }(\varphi )}{H(\varphi )}}{\displaystyle \frac{1}{\kappa ^2}}{\displaystyle \frac{V^{\prime \prime }(\varphi )}{V(\varphi )}}{\displaystyle \frac{1}{2\kappa ^2}}\left({\displaystyle \frac{V^{}(\varphi )}{V(\varphi )}}\right)^2\eta ϵ1.`$ (125) The last expression defines the new slow-roll parameter $`\eta `$, not to be confused with conformal time (see next Section). The number of $`e`$-folds can also be rewritten in this approximation as $$N=\kappa ^2_{\varphi _i}^{\varphi _e}\frac{V(\varphi )d\varphi }{V^{}(\varphi )},$$ (126) a very useful expression for evaluating $`N`$ for a given effective scalar potential $`V(\varphi )`$. ### 4.3 The origin of density perturbations If cosmological inflation made the universe so extremely flat and homogeneous, where did the galaxies and clusters of galaxies come from? One of the most astonishing predictions of inflation, one that was not even expected, is that quantum fluctuations of the inflaton field are stretched by the exponential expansion and generate large-scale perturbations in the metric. Inflaton fluctuations are small wave packets of energy that, according to general relativity, modify the space-time fabric, creating a whole spectrum of curvature perturbations. The use of the word spectrum here is closely related to the case of light waves propagating in a medium: a spectrum characterizes the amplitude of each given wavelength. In the case of inflation, the inflaton fluctuations induce waves in the space-time metric that can be decomposed into different wavelengths, all with approximately the same amplitude, that is, corresponding to a scale-invariant spectrum. These patterns of perturbations in the metric are like fingerprints that unequivocally characterize a period of inflation. When matter fell in the troughs of these waves, it created density perturbations that collapsed gravitationally to form galaxies, clusters and superclusters of galaxies, with a spectrum that is also scale invariant. Such a type of spectrum was proposed in the early 1970s (before inflation) by Harrison and Zel’dovich , to explain the distribution of galaxies and clusters of galaxies on very large scales in our observable universe. Perhaps the most interesting aspect of structure formation is the possibility that the detailed knowledge of what seeded galaxies and clusters of galaxies will allow us to test the idea of inflation. #### 4.3.1 Gauge invariant perturbation theory Until now we have considered only the unperturbed FRW metric described by a scale factor $`a(t)`$ and a homogeneous scalar field $`\varphi (t)`$, $`ds^2`$ $`=`$ $`a^2(\eta )[d\eta ^2\gamma _{ij}dx^idx^j],`$ (127) $`\varphi `$ $`=`$ $`\varphi (\eta ),`$ (128) where $`\eta =𝑑t/a(t)`$ is the conformal time, under which the background equations of motion can be written as $`^2`$ $`=`$ $`{\displaystyle \frac{\kappa ^2}{3}}\left({\displaystyle \frac{1}{2}}\varphi _{}^{}{}_{}{}^{2}+a^2V(\varphi )\right),`$ (129) $`^{}^2`$ $`=`$ $`{\displaystyle \frac{\kappa ^2}{2}}\varphi _{}^{}{}_{}{}^{2},`$ (130) $`\varphi ^{\prime \prime }+2\varphi ^{}+a^2V^{}(\varphi )`$ $`=`$ $`0,`$ (131) where $`=aH`$ and $`\varphi ^{}=a\dot{\varphi }`$. During inflation, the quantum fluctuations of the scalar field will induce metric perturbations which will backreact on the scalar field. Let us consider, in linear perturbation theory, the most general line element with both scalar and tensor metric perturbations ,<sup>18</sup><sup>18</sup>18Note that inflation cannot generate, to linear order, a vector perturbation. together with the scalar field perturbations $`ds^2`$ $`=`$ $`a^2(\eta )\left[(1+2A)d\eta ^22B_{\|i}dx^id\eta \left\{(1+2)\gamma _{ij}+2E_{\|ij}+2h_{ij}\right\}dx^idx^j\right],`$ (132) $`\varphi `$ $`=`$ $`\varphi (\eta )+\delta \varphi (\eta ,x^i).`$ (133) The indices $`\{i,j\}`$ label the three-dimensional spatial coordinates with metric $`\gamma _{ij}`$, and the $`\|i`$ denotes covariant derivative with respect to that metric. The gauge invariant tensor perturbation $`h_{ij}`$ corresponds to a transverse traceless gravitational wave, $`^ih_{ij}=h_i^i=0`$. The four scalar perturbations $`(A,B,,E)`$ are gauge dependent functions of $`(\eta ,x^i)`$. Under a general coordinate (gauge) transformation $`\stackrel{~}{\eta }=\eta +\xi ^0(\eta ,x^i),`$ (134) $`\stackrel{~}{x}^i=x^i+\gamma ^{ij}\xi _{\|j}(\eta ,x^i),`$ (135) with arbitrary functions $`(\xi ^0,\xi )`$, the scalar and tensor perturbations transform, to linear order, as $`\stackrel{~}{A}=A\xi _{}^{0}{}_{}{}^{}\xi ^0,\stackrel{~}{B}=B+\xi ^0\xi ^{},`$ (136) $`\stackrel{~}{}=\xi ^0,\stackrel{~}{E}=E\xi ,`$ (137) $`\stackrel{~}{h}_{ij}=h_{ij},`$ (138) where a prime denotes derivative with respect to conformal time. It is possible to construct, however, two gauge-invariant gravitational potentials , $`\mathrm{\Phi }=A+(BE^{})^{}+(BE^{}),`$ (139) $`\mathrm{\Psi }=+(BE^{}),`$ (140) which are related through the perturbed Einstein equations, $`\mathrm{\Phi }`$ $`=`$ $`\mathrm{\Psi },`$ (141) $`{\displaystyle \frac{k^23K}{a^2}}\mathrm{\Psi }`$ $`=`$ $`{\displaystyle \frac{\kappa ^2}{2}}\delta \rho ,`$ (142) where $`\delta \rho `$ is the gauge-invariant density perturbation, and the latter expression is nothing but the Poisson equation for the gravitational potential, written in relativistic form. During inflation, the energy density is given in terms of a scalar field, and thus the gauge-invariant equations for the perturbations on comoving hypersurfaces (constant energy density hypersurfaces) are $`\mathrm{\Phi }^{\prime \prime }+3\mathrm{\Phi }^{}+(^{}+^2)\mathrm{\Phi }`$ $`=`$ $`{\displaystyle \frac{\kappa ^2}{2}}[\varphi ^{}\delta \varphi ^{}a^2V^{}(\varphi )\delta \varphi ],`$ (143) $`^2\mathrm{\Phi }+3\mathrm{\Phi }^{}+(^{}+^2)\mathrm{\Phi }`$ $`=`$ $`{\displaystyle \frac{\kappa ^2}{2}}[\varphi ^{}\delta \varphi ^{}+a^2V^{}(\varphi )\delta \varphi ],`$ (144) $`\mathrm{\Phi }^{}+\mathrm{\Phi }`$ $`=`$ $`{\displaystyle \frac{\kappa ^2}{2}}\varphi ^{}\delta \varphi ,`$ (145) $`\delta \varphi ^{\prime \prime }+2\delta \varphi ^{}^2\delta \varphi `$ $`=`$ $`4\varphi ^{}\mathrm{\Phi }^{}2a^2V^{}(\varphi )\mathrm{\Phi }a^2V^{\prime \prime }(\varphi )\delta \varphi .`$ (146) This system of equations seem too difficult to solve at first sight. However, there is a gauge invariant combination of variables that allows one to find exact solutions. Let us define $`ua\delta \varphi +z\mathrm{\Phi },`$ (147) $`za{\displaystyle \frac{\varphi ^{}}{}}.`$ (148) Under this redefinition, the above equations simplify enormously to just three independent equations, $`u^{\prime \prime }^2u{\displaystyle \frac{z^{\prime \prime }}{z}}u=0,`$ (149) $`^2\mathrm{\Phi }={\displaystyle \frac{\kappa ^2}{2}}{\displaystyle \frac{}{a^2}}(zu^{}z^{}u),`$ (150) $`\left({\displaystyle \frac{a^2\mathrm{\Phi }}{}}\right)^{}={\displaystyle \frac{\kappa ^2}{2}}zu.`$ (151) From Equation (149) we can find a solution $`u(z)`$, which substituted into (151) can be integrated to give $`\mathrm{\Phi }(z)`$, and together with $`u(z)`$ allow us to obtain $`\delta \varphi (z)`$. #### 4.3.2 Quantum Field Theory in curved space-time Until now we have treated the perturbations as classical, but we should in fact consider the perturbations $`\mathrm{\Phi }`$ and $`\delta \varphi `$ as quantum fields. Note that the perturbed action for the scalar mode $`u`$ can be written as $$\delta S=\frac{1}{2}d^3x𝑑\eta \left[(u^{})^2(u)^2+\frac{z^{\prime \prime }}{z}u^2\right].$$ (152) In order to quantize the field $`u`$ in the curved background defined by the metric (127), we can write the operator $$\widehat{u}(\eta ,𝐱)=\frac{d^3𝐤}{(2\pi )^{3/2}}\left[u_k(\eta )\widehat{a}_𝐤e^{i𝐤𝐱}+u_k^{}(\eta )\widehat{a}_𝐤^{}e^{i𝐤𝐱}\right],$$ (153) where the creation and annihilation operators satisfy the commutation relation of bosonic fields, and the scalar field’s Fock space is defined through the vacuum condition, $`[\widehat{a}_𝐤,\widehat{a}_𝐤^{}^{}]`$ $`=`$ $`\delta ^3(𝐤𝐤^{}),`$ (154) $`\widehat{a}_𝐤\|0`$ $`=`$ $`0.`$ (155) Note that we are not assuming that the inflaton is a fundamental scalar field, but that is can be written as a quantum field with its commutation relations (as much as a pion can be described as a quantum field). The equations of motion for each mode $`u_k(\eta )`$ are decoupled in linear perturbation theory, $$u_k^{\prime \prime }+\left(k^2\frac{z^{\prime \prime }}{z}\right)u_k=0.$$ (156) The ratio $`z^{\prime \prime }/z`$ acts like a time-dependent potential for this Schrödinger like equation. In order to find exact solutions to the mode equation, we will use the slow-roll parameters (118), see Ref. $`ϵ=1{\displaystyle \frac{^{}}{^2}}={\displaystyle \frac{\kappa ^2}{2}}{\displaystyle \frac{z^2}{a^2}},`$ (157) $`\delta =1{\displaystyle \frac{\varphi ^{\prime \prime }}{\varphi ^{}}}=1+ϵ{\displaystyle \frac{z^{}}{z}}.`$ (158) In terms of these parameters, the conformal time and the effective potential for the $`u_k`$ mode can be written as $`\eta ={\displaystyle \frac{1}{}}+{\displaystyle \frac{ϵda}{a}},`$ (159) $`{\displaystyle \frac{z^{\prime \prime }}{z}}=^2[(1+ϵ\delta )(2\delta )+^1(ϵ^{}\delta ^{})].`$ (160) Note that the slow-roll parameters, (157) and (158), can be taken as constant,<sup>19</sup><sup>19</sup>19For instance, there are models of inflation, like power-law inflation, $`a(t)t^p`$, where $`ϵ=\delta =1/p<1`$, that give constant slow-roll parameters. to order $`ϵ^2`$, $`ϵ^{}=2ϵ(ϵ\delta )=𝒪(ϵ^2),`$ (161) $`\delta ^{}=\delta \left(ϵ+\delta +{\displaystyle \frac{\stackrel{\mathrm{}}{\varphi }}{H\ddot{\varphi }}}\right)=𝒪(ϵ^2).`$ (162) In that case, for constant parameters, we can write $`\eta ={\displaystyle \frac{1}{}}{\displaystyle \frac{1}{1ϵ}},`$ (163) $`{\displaystyle \frac{z^{\prime \prime }}{z}}={\displaystyle \frac{1}{\eta ^2}}\left(\nu ^2{\displaystyle \frac{1}{4}}\right),`$ $`\mathrm{where}\nu ={\displaystyle \frac{1+ϵ\delta }{1ϵ}}+{\displaystyle \frac{1}{2}}.`$ (164) We are now going to search for approximate solutions of the mode equation (156), where the effective potential (160) is of order $`z^{\prime \prime }/z2^2`$ in the slow-roll approximation. In quasi-de Sitter there is a characteristic scale given by the (event) horizon size or Hubble scale during inflation, $`H^1`$. There will be modes $`u_k`$ with physical wavelengths much smaller than this scale, $`k/aH`$, that are well within the de Sitter horizon and therefore do not feel the curvature of space-time. On the other hand, there will be modes with physical wavelengths much greater than the Hubble scale, $`k/aH`$. In these two asymptotic regimes, the solutions can be written as $`u_k={\displaystyle \frac{1}{\sqrt{2k}}}e^{ik\eta }kaH,`$ (165) $`u_k=C_1zkaH.`$ (166) In the limit $`kaH`$ the modes behave like ordinary quantum modes in Minkowsky space-time, appropriately normalized, while in the opposite limit, $`u/z`$ becomes constant on superhorizon scales. For approximately constant slow-roll parameters one can find exact solutions to (156), with the effective potential given by (164), that interpolate between the two asymptotic solutions, $$u_k(\eta )=\frac{\sqrt{\pi }}{2}e^{i(\nu +\frac{1}{2})\frac{\pi }{2}}(\eta )^{1/2}H_\nu ^{_{(1)}}(k\eta ),$$ (167) where $`H_\nu ^{_{(1)}}(z)`$ is the Hankel function of the first kind , and $`\nu `$ is given by (164) in terms of the slow-roll parameters. In the limit $`k\eta 0`$, the solution becomes $`\|u_k\|={\displaystyle \frac{2^{\nu \frac{3}{2}}}{\sqrt{2k}}}{\displaystyle \frac{\mathrm{\Gamma }(\nu )}{\mathrm{\Gamma }(\frac{3}{2})}}(k\eta )^{\frac{1}{2}\nu }{\displaystyle \frac{C(\nu )}{\sqrt{2k}}}\left({\displaystyle \frac{k}{aH}}\right)^{\nu \frac{1}{2}},`$ (168) $`C(\nu )=2^{\nu \frac{3}{2}}{\displaystyle \frac{\mathrm{\Gamma }(\nu )}{\mathrm{\Gamma }(\frac{3}{2})}}(1ϵ)^{\nu \frac{1}{2}}1\mathrm{for}ϵ,\delta 1.`$ (169) We can now compute $`\mathrm{\Phi }`$ and $`\delta \varphi `$ from the super-Hubble-scale mode solution (166), for $`kaH`$. Substituting into Eq. (151), we find $`\mathrm{\Phi }=C_1\left(1{\displaystyle \frac{}{a^2}}{\displaystyle a^2𝑑\eta }\right)+C_2{\displaystyle \frac{}{a^2}},`$ (170) $`\delta \varphi ={\displaystyle \frac{C_1}{a^2}}{\displaystyle a^2𝑑\eta }{\displaystyle \frac{C_2}{a^2}}.`$ (171) The term proportional to $`C_1`$ corresponds to the growing solution, while that proportional to $`C_2`$ corresponds to the decaying solution, which can soon be ignored. These quantities are gauge invariant but evolve with time outside the horizon, during inflation, and before entering again the horizon during the radiation or matter eras. We would like to write an expression for a gauge invariant quantity that is also constant for superhorizon modes. Fortunately, in the case of adiabatic perturbations, there is such a quantity: $$\zeta \mathrm{\Phi }+\frac{1}{ϵ}(\mathrm{\Phi }^{}+\mathrm{\Phi })=\frac{u}{z},$$ (172) which is constant, see Eq. (166), for $`kaH`$. In fact, this quantity $`\zeta `$ is identical, for superhorizon modes, to the gauge invariant curvature metric perturbation $`_c`$ on comoving (constant energy density) hypersurfaces, see Ref. , $$\zeta =_c+\frac{1}{ϵ^2}^2\mathrm{\Phi }.$$ (173) Using Eq. (150) we can write the evolution equation for $`\zeta =\frac{u}{z}`$ as $`\zeta ^{}=\frac{1}{ϵ}^2\mathrm{\Phi }`$, which confirms that $`\zeta `$ is constant for (adiabatic<sup>20</sup><sup>20</sup>20This conservation fails for entropy or isocurvature perturbations, see Ref. .) superhorizon modes, $`kaH`$. Therefore, we can evaluate the Newtonian potential $`\mathrm{\Phi }_k`$ when the perturbation reenters the horizon during radiation/matter eras in terms of the curvature perturbation $`_k`$ when it left the Hubble scale during inflation, $$\mathrm{\Phi }_k=\left(1\frac{}{a^2}a^2𝑑\eta \right)_k=\frac{3+3\omega }{5+3\omega }_k=\{\begin{array}{cc}\frac{2}{3}_k& \mathrm{radiation}\mathrm{era},\hfill \\ \frac{3}{5}_k& \mathrm{matter}\mathrm{era}.\hfill \end{array}$$ (174) Let us now compute the tensor or gravitational wave metric perturbations generated during inflation. The perturbed action for the tensor mode can be written as $$\delta S=\frac{1}{2}d^3x𝑑\eta \frac{a^2}{2\kappa ^2}\left[(h_{ij}^{})^2(h_{ij})^2\right],$$ (175) with the tensor field $`h_{ij}`$ considered as a quantum field, $$\widehat{h}_{ij}(\eta ,𝐱)=\frac{d^3𝐤}{(2\pi )^{3/2}}\underset{\lambda =1,2}{}[h_k(\eta )e_{ij}(𝐤,\lambda )\widehat{a}_{𝐤,\lambda }e^{i𝐤𝐱}+h.c.],$$ (176) where $`e_{ij}(𝐤,\lambda )`$ are the two polarization tensors, satisfying symmetric, transverse and traceless conditions $`e_{ij}=e_{ji},k^ie_{ij}=0,e_{ii}=0,`$ (177) $`e_{ij}(𝐤,\lambda )=e_{ij}^{}(𝐤,\lambda ),{\displaystyle \underset{\lambda }{}}e_{ij}^{}(𝐤,\lambda )e^{ij}(𝐤,\lambda )=4,`$ (178) while the creation and annihilation operators satisfy the usual commutation relation of bosonic fields, Eq. (154). We can now redefine our gauge invariant tensor amplitude as $$v_k(\eta )=\frac{a}{\sqrt{2}\kappa }h_k(\eta ),$$ (179) which satisfies the following evolution equation, decoupled for each mode $`v_k(\eta )`$ in linear perturbation theory, $$v_k^{\prime \prime }+\left(k^2\frac{a^{\prime \prime }}{a}\right)v_k=0.$$ (180) The ratio $`a^{\prime \prime }/a`$ acts like a time-dependent potential for this Schrödinger like equation, analogous to the term $`z^{\prime \prime }/z`$ for the scalar metric perturbation. For constant slow-roll parameters, the potential becomes $`{\displaystyle \frac{a^{\prime \prime }}{a}}=2^2\left(1{\displaystyle \frac{ϵ}{2}}\right)={\displaystyle \frac{1}{\eta ^2}}\left(\mu ^2{\displaystyle \frac{1}{4}}\right),`$ (181) $`\mu ={\displaystyle \frac{1}{1ϵ}}+{\displaystyle \frac{1}{2}}.`$ (182) We can solve equation (180) in the two asymptotic regimes, $`v_k={\displaystyle \frac{1}{\sqrt{2k}}}e^{ik\eta }kaH,`$ (183) $`v_k=CakaH.`$ (184) In the limit $`kaH`$ the modes behave like ordinary quantum modes in Minkowsky space-time, appropriately normalized, while in the opposite limit, the metric perturbation $`h_k`$ becomes constant on superhorizon scales. For constant slow-roll parameters one can find exact solutions to (180), with effective potential given by (181), that interpolate between the two asymptotic solutions. These are identical to Eq. (167) except for the substitution $`\nu \mu `$. In the limit $`k\eta 0`$, the solution becomes $$\|v_k\|=\frac{C(\mu )}{\sqrt{2k}}\left(\frac{k}{aH}\right)^{\mu \frac{1}{2}}.$$ (185) Since the mode $`h_k`$ becomes constant on superhorizon scales, we can evaluate the tensor metric perturbation when it reentered during the radiation or matter era directly in terms of its value during inflation. #### 4.3.3 Power spectrum of scalar and tensor metric perturbations Not only do we expect to measure the amplitude of the metric perturbations generated during inflation and responsible for the anisotropies in the CMB and density fluctuations in LSS, but we should also be able to measure its power spectrum, or two-point correlation function in Fourier space. Let us consider first the scalar metric perturbations $`_k`$, which enter the horizon at $`a=k/H`$. Its correlator is given by $`0\|_k^{}_k^{}\|0={\displaystyle \frac{\|u_k\|^2}{z^2}}\delta ^3(𝐤𝐤^{}){\displaystyle \frac{𝒫_{}(k)}{4\pi k^3}}(2\pi )^3\delta ^3(𝐤𝐤^{}),`$ (186) $`𝒫_{}(k)={\displaystyle \frac{k^3}{2\pi ^2}}{\displaystyle \frac{\|u_k\|^2}{z^2}}={\displaystyle \frac{\kappa ^2}{2ϵ}}\left({\displaystyle \frac{H}{2\pi }}\right)^2\left({\displaystyle \frac{k}{aH}}\right)^{32\nu }A_S^2\left({\displaystyle \frac{k}{aH}}\right)^{n1},`$ (187) where we have used $`_k=\zeta _k=\frac{u_k}{z}`$ and Eq. (168). This last equation determines the power spectrum in terms of its amplitude at horizon-crossing, $`A_S`$, and a tilt, $$n1\frac{d\mathrm{ln}𝒫_{}(k)}{d\mathrm{ln}k}=32\nu =2\left(\frac{\delta 2ϵ}{1ϵ}\right)2\eta 6ϵ,$$ (188) see Eqs. (124), (125). Note from this equation that it is possible, in principle, to obtain from inflation a scalar tilt which is either positive ($`n>1`$) or negative ($`n<1`$). Furthermore, depending on the particular inflationary model , we can have significant departures from scale invariance. Let us consider now the tensor (gravitational wave) metric perturbation, which enter the horizon at $`a=k/H`$, $`{\displaystyle \underset{\lambda }{}}0\|h_{k,\lambda }^{}h_{k^{},\lambda }\|0=4{\displaystyle \frac{2\kappa ^2}{a^2}}\|v_k\|^2\delta ^3(𝐤𝐤^{}){\displaystyle \frac{𝒫_g(k)}{4\pi k^3}}(2\pi )^3\delta ^3(𝐤𝐤^{}),`$ (189) $`𝒫_g(k)=8\kappa ^2\left({\displaystyle \frac{H}{2\pi }}\right)^2\left({\displaystyle \frac{k}{aH}}\right)^{32\mu }A_T^2\left({\displaystyle \frac{k}{aH}}\right)^{n_T},`$ (190) where we have used Eqs. (179) and (185). Therefore, the power spectrum can be approximated by a power-law expression, with amplitude $`A_T`$ and tilt $$n_T\frac{d\mathrm{ln}𝒫_g(k)}{d\mathrm{ln}k}=32\mu =\left(\frac{2ϵ}{1ϵ}\right)2ϵ<0,$$ (191) which is always negative. In the slow-roll approximation, $`ϵ1`$, the tensor power spectrum is scale invariant. ### 4.4 The anisotropies of the microwave background The metric fluctuations generated during inflation are not only responsible for the density perturbations that gave rise to galaxies via gravitational collapse, but one should also expect to see such ripples in the metric as temperature anisotropies in the cosmic microwave background, that is, minute deviations in the temperature of the blackbody spectrum when we look at different directions in the sky. Such anisotropies had been looked for ever since Penzias and Wilson’s discovery of the CMB, but had eluded all detection, until COBE satellite discovered them in 1992, see Fig. 6. The reason why they took so long to be discovered was that they appear as perturbations in temperature of only one part in $`10^5`$. Soon after COBE, other groups quickly confirmed the detection of temperature anisotropies at around 30 $`\mu `$K, at higher multipole numbers or smaller angular scales. There are at this moment dozens of ground and balloon-borne experiments analysing the anisotropies in the microwave background with angular resolutions from $`10^{}`$ to a few arc minutes in the sky, see Fig. 34. #### 4.4.1 Acoustic oscillations in the plasma The physics of the CMB anisotropies is relatively simple . The universe just before recombination is a very tightly coupled fluid, due to the large electromagnetic Thomson cross section (54). Photons scatter off charged particles (protons and electrons), and carry energy, so they feel the gravitational potential associated with the perturbations imprinted in the metric during inflation. An overdensity of baryons (protons and neutrons) does not collapse under the effect of gravity until it enters the causal Hubble radius. The perturbation continues to grow until radiation pressure opposes gravity and sets up acoustic oscillations in the plasma, very similar to sound waves. Since overdensities of the same size will enter the Hubble radius at the same time, they will oscillate in phase. Moreover, since photons scatter off these baryons, the acoustic oscillations occur also in the photon field and induces a pattern of peaks in the temperature anisotropies in the sky, at different angular scales, see Fig. 34. There are three different effects that determine the temperature anisotropies we observe in the CMB. First, gravity: photons fall in and escape off gravitational potential wells, characterized by $`\mathrm{\Phi }`$ in the comoving gauge, and as a consequence their frequency is gravitationally blue- or red-shifted, $`\delta \nu /\nu =\mathrm{\Phi }`$. If the gravitational potential is not constant, the photons will escape from a larger or smaller potential well than they fell in, so their frequency is also blue- or red-shifted, a phenomenon known as the Rees-Sciama effect. Second, pressure: photons scatter off baryons which fall into gravitational potential wells and the two competing forces create acoustic waves of compression and rarefaction. Finally, velocity: baryons accelerate as they fall into potential wells. They have minimum velocity at maximum compression and rarefaction. That is, their velocity wave is exactly $`90^{}`$ off-phase with the acoustic waves. These waves induce a Doppler effect on the frequency of the photons. The temperature anisotropy induced by these three effects is therefore given by $$\frac{\delta T}{T}(𝐫)=\mathrm{\Phi }(𝐫,t_{\mathrm{dec}})+2_{t_{\mathrm{dec}}}^{t_0}\dot{\mathrm{\Phi }}(𝐫,t)𝑑t+\frac{1}{3}\frac{\delta \rho }{\rho }\frac{𝐫𝐯}{c}.$$ (192) Metric perturbations of different wavelengths enter the horizon at different times. The largest wavelengths, of size comparable to our present horizon, are entering now. There are perturbations with wavelengths comparable to the size of the horizon at the time of last scattering, of projected size about $`1^{}`$ in the sky today, which entered precisely at decoupling. And there are perturbations with wavelengths much smaller than the size of the horizon at last scattering, that entered much earlier than decoupling, all the way to the time of radiation-matter equality, which have gone through several acoustic oscillations before last scattering. All these perturbations of different wavelengths leave their imprint in the CMB anisotropies. The baryons at the time of decoupling do not feel the gravitational attraction of perturbations with wavelength greater than the size of the horizon at last scattering, because of causality. Perturbations with exactly that wavelength are undergoing their first contraction, or acoustic compression, at decoupling. Those perturbations induce a large peak in the temperature anisotropies power spectrum, see Fig. 34. Perturbations with wavelengths smaller than these will have gone, after they entered the Hubble scale, through a series of acoustic compressions and rarefactions, which can be seen as secondary peaks in the power spectrum. Since the surface of last scattering is not a sharp discontinuity, but a region of $`\mathrm{\Delta }z100`$, see Fig. 4, there will be scales for which photons, travelling from one energy concentration to another, will erase the perturbation on that scale, similarly to what neutrinos or HDM do for structure on small scales. That is the reason why we don’t see all the acoustic oscillations with the same amplitude, but in fact they decay exponentialy towards smaller angular scales, an effect known as Silk damping, due to photon diffusion . #### 4.4.2 The Sachs-Wolfe effect The anisotropies corresponding to large angular scales are only generated via gravitational red-shift and density perturbations through the Einstein equations, $`\delta \rho /\rho =2\mathrm{\Phi }`$ for adiabatic perturbations; we can ignore the Doppler contribution, since the perturbation is non-causal. In that case, the temperature anisotropy in the sky today is given by $$\frac{\delta T}{T}(\theta ,\varphi )=\frac{1}{3}\mathrm{\Phi }(\eta _{\mathrm{LS}})Q(\eta _0,\theta ,\varphi )+2_{\eta _{\mathrm{LS}}}^{\eta _0}𝑑r\mathrm{\Phi }^{}(\eta _0r)Q(r,\theta ,\varphi ),$$ (193) where $`\eta _0`$ is the coordinate distance to the last scattering surface, i.e. the present conformal time, while $`\eta _{\mathrm{LS}}0`$ determines that comoving hypersurface. The above expression is known as the Sachs-Wolfe effect , and contains two parts, the intrinsic and the Integrated Sachs-Wolfe (ISW) effect, due to integration along the line of sight of time variations in the gravitational potential. In linear perturbation theory, the scalar metric perturbations can be separated into $`\mathrm{\Phi }(\eta ,𝐱)\mathrm{\Phi }(\eta )Q(𝐱)`$, where $`Q(𝐱)`$ are the scalar harmonics, eigenfunctions of the Laplacian in three dimensions, $`^2Q_{klm}(r,\theta ,\varphi )=k^2Q_{klm}(r,\theta ,\varphi )`$. These functions have the general form $$Q_{klm}(r,\theta ,\varphi )=\mathrm{\Pi }_{kl}(r)Y_{lm}(\theta ,\varphi ),$$ (194) where $`Y_{lm}(\theta ,\varphi )`$ are the usual spherical harmonics . In order to compute the temperature anisotropy associated with the Sachs-Wolfe effect, we have to know the evolution of the metric perturbation during the matter era, $$\mathrm{\Phi }^{\prime \prime }+3\mathrm{\Phi }^{}+a^2\mathrm{\Lambda }\mathrm{\Phi }2K\mathrm{\Phi }=0.$$ (195) In the case of a flat universe without cosmological constant, the Newtonian potential remains constant during the matter era and only the intrinsic SW effect contributes to $`\delta T/T`$. In case of a non-vanishing $`\mathrm{\Lambda }`$, since its contribution is negligible in the past, see Eq. (25), most of the photon’s trajectory towards us is unperturbed, and the only difference with respect to the $`\mathrm{\Lambda }=0`$ case is an overall factor . We will consider here the approximation $`\mathrm{\Phi }=\mathrm{const}`$. during the matter era and ignore that factor, see Ref. . In a flat universe, the radial part of the eigenfunctions (194) can be written as $$\mathrm{\Pi }_{kl}(r)=\sqrt{\frac{2}{\pi }}kj_l(kr),$$ (196) where $`j_l(z)`$ are the spherical Bessel functions . The growing mode solution of the metric perturbation that left the Hubble scale during inflation contributes to the temperature anisotropies on large scales (193) as $$\frac{\delta T}{T}(\theta ,\varphi )=\frac{1}{3}\mathrm{\Phi }(\eta _{\mathrm{LS}})Q=\frac{1}{5}Q(\eta _0,\theta ,\varphi )\underset{l=2}{\overset{\mathrm{}}{}}\underset{m=l}{\overset{l}{}}a_{lm}Y_{lm}(\theta ,\varphi ),$$ (197) where we have used the fact that at reentry (at the surface of last scattering) the gauge invariant Newtonian potential $`\mathrm{\Phi }`$ is related to the curvature perturbation $``$ at Hubble-crossing during inflation, see Eq. (174); and we have expanded $`\delta T/T`$ in spherical harmonics. We can now compute the two-point correlation function or angular power spectrum, $`C(\theta )`$, of the CMB anisotropies on large scales, defined as an expansion in multipole number, $$C(\theta )=\frac{\delta T}{T}^{}(𝐧)\frac{\delta T}{T}(𝐧^{})_{𝐧𝐧^{}=\mathrm{cos}\theta }=\frac{1}{4\pi }\underset{l=2}{\overset{\mathrm{}}{}}(2l+1)C_lP_l(\mathrm{cos}\theta ),$$ (198) where $`P_l(z)`$ are the Legendre polynomials , and we have averaged over different universe realizations. Since the coefficients $`a_{lm}`$ are isotropic (to first order), we can compute the $`C_l=\|a_{lm}\|^2`$ as $$C_l^{(S)}=\frac{4\pi }{25}_0^{\mathrm{}}\frac{dk}{k}𝒫_{}(k)j_l^2(k\eta _0),$$ (199) where we have used Eqs. (197) and (186). In the case of scalar metric perturbation produced during inflation, the scalar power spectrum at reentry is given by $`𝒫_{}(k)=A_S^2(k\eta _0)^{n1}`$, in the power-law approximation, see Eq. (187). In that case, one can integrate (199) to give $`C_l^{(S)}={\displaystyle \frac{2\pi }{25}}A_S^2{\displaystyle \frac{\mathrm{\Gamma }[\frac{3}{2}]\mathrm{\Gamma }[1\frac{n1}{2}]\mathrm{\Gamma }[l+\frac{n1}{2}]}{\mathrm{\Gamma }[\frac{3}{2}\frac{n1}{2}]\mathrm{\Gamma }[l+2\frac{n1}{2}]}},`$ (200) $`{\displaystyle \frac{l(l+1)C_l^{(S)}}{2\pi }}={\displaystyle \frac{A_S^2}{25}}=\mathrm{constant},\mathrm{for}n=1.`$ (201) This last expression corresponds to what is known as the Sachs-Wolfe plateau, and is the reason why the coefficients $`C_l`$ are always plotted multiplied by $`l(l+1)`$, see Fig. 34. Tensor metric perturbations also contribute with an approximately constant angular power spectrum, $`l(l+1)C_l`$. The Sachs-Wolfe effect for a gauge invariant tensor perturbation is given by $$\frac{\delta T}{T}(\theta ,\varphi )=_{\eta _{\mathrm{LS}}}^{\eta _0}𝑑rh^{}(\eta _0r)Q_{rr}(r,\theta ,\varphi ),$$ (202) where $`Q_{rr}`$ is the $`rr`$-component of the tensor harmonic along the line of sight . The tensor perturbation $`h`$ during the matter era satisfies the following evolution equation $$h_k^{\prime \prime }+3h_k^{}+(k^2+2K)h_k=0,$$ (203) which depends on the wavenumber $`k`$, contrary to what happens with the scalar modes, see Eq. (195). For a flat ($`K=0`$) universe, the solution to this equation is $`h_k(\eta )=hG_k(\eta )`$, where $`h`$ is the constant tensor metric perturbation at horizon crossing and $`G_k(\eta )=3j_1(k\eta )/k\eta `$, normalized so that $`G_k(0)=1`$ at the surface of last scattering. The radial part of the tensor harmonic $`Q_{rr}`$ in a flat universe can be written as $$Q_{kl}^{rr}(r)=\left[\frac{(l1)l(l+1)(l+2)}{\pi k^2}\right]^{1/2}\frac{j_l(kr)}{r^2}.$$ (204) The tensor angular power spectrum can finally be expressed as $`C_l^{(T)}={\displaystyle \frac{9\pi }{4}}(l1)l(l+1)(l+2){\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dk}{k}}𝒫_g(k)I_{kl}^2,`$ (205) $`I_{kl}={\displaystyle _0^{x_0}}𝑑x{\displaystyle \frac{j_2(x_0x)j_l(x)}{(x_0x)x^2}},`$ (206) where $`xk\eta `$, and $`𝒫_g(k)`$ is the primordial tensor spectrum (190). For a scale invariant spectrum, $`n_T=0`$, we can integrate (205) to give $$l(l+1)C_l^{(T)}=\frac{\pi }{36}\left(1+\frac{48\pi ^2}{385}\right)A_T^2B_l,$$ (207) with $`B_l=(1.1184,0.8789,\mathrm{},1.00)`$ for $`l=2,3,\mathrm{},30`$. Therefore, $`l(l+1)C_l^{(T)}`$ also becomes constant for large $`l`$. Beyond $`l30`$, the Sachs-Wolfe expression is not a good approximation and the tensor angular power spectrum decays very quickly at large $`l`$, see Fig.40. #### 4.4.3 The consistency relation In spite of the success of inflation in predicting a homogeneous and isotropic background on which to imprint a scale-invariant spectrum of inhomogeneities, it is difficult to test the idea of inflation. A CMB cosmologist before the 1980s would have argued that ad hoc initial conditions could have been at the origin of the homogeneity and flatness of the universe on large scales, while a LSS cosmologist would have agreed with Harrison and Zel’dovich that the most natural spectrum needed to explain the formation of structure was a scale-invariant spectrum. The surprise was that inflation incorporated an understanding of both the globally homogeneous and spatially flat background, and the approximately scale-invariant spectrum of perturbations in the same formalism. But that could have been a coincidence, and is not epistemologically testable. What is unique to inflation is the fact that inflation determines not just one but two primordial spectra, corresponding to the scalar (density) and tensor (gravitational waves) metric perturbations, from a single continuous function, the inflaton potential $`V(\varphi )`$. In the slow-roll approximation, one determines, from $`V(\varphi )`$, two continuous functions, $`𝒫_{}(k)`$ and $`𝒫_g(k)`$, that in the power-law approximation reduces to two amplitudes, $`A_S`$ and $`A_T`$, and two tilts, $`n`$ and $`n_T`$. It is clear that there must be a relation between the four parameters. Indeed, one can see from Eqs. (207) and (201) that the ratio of the tensor to scalar contribution to the angular power spectrum is proportional to the tensor tilt , $$R\frac{C_l^{(T)}}{C_l^{(S)}}=\frac{25}{9}\left(1+\frac{48\pi ^2}{385}\right)\mathrm{\hspace{0.17em}2}ϵ2\pi n_T.$$ (208) This is a unique prediction of inflation, which could not have been postulated a priori by any cosmologist. If we finally observe a tensor spectrum of anisotropies in the CMB, or a stochastic gravitational wave background in laser interferometers like LIGO or VIRGO , with sufficient accuracy to determine their spectral tilt, one might have some chance to test the idea of inflation, via the consistency relation (208). For the moment, observations of the microwave background anisotropies suggest that the Sachs-Wolfe plateau exists, see Fig. 34, but it is still premature to determine the tensor contribution. Perhaps in the near future, from the analysis of polarization as well as temperature anisotropies, with the CMB satellites MAP and Planck, we might have a chance of determining the validity of the consistency relation. Assuming that the scalar contribution dominates over the tensor on large scales, i.e. $`R1`$, one can actually give a measure of the amplitude of the scalar metric perturbation from the observations of the Sachs-Wolfe plateau in the angular power spectrum , $`\left[{\displaystyle \frac{l(l+1)C_l^{(S)}}{2\pi }}\right]^{1/2}`$ $`=`$ $`{\displaystyle \frac{A_S}{5}}=(1.03\pm 0.07)\times 10^5,`$ (209) $`n`$ $`=`$ $`1.02\pm 0.12.`$ (210) These measurements can be used to normalize the primordial spectrum and determine the parameters of the model of inflation . In the near future these parameters will be determined with much better accuracy, as described in Section 4.4.5. #### 4.4.4 The acoustic peaks The Sachs-Wolfe plateau is a distinctive feature of Fig. 34. These observations confirm the existence of a primordial spectrum of scalar (density) perturbations on all scales, otherwise the power spectrum would have started from zero at $`l=2`$. However, we see that the spectrum starts to rise around $`l=20`$ towards the first acoustic peak, where the SW approximation breaks down and the above formulae are no longer valid. As mentioned above, the first peak in the photon distribution corresponds to overdensities that have undergone half an oscillation, that is, a compression, and appear at a scale associated with the size of the horizon at last scattering, about $`1^{}`$ projected in the sky today. Since photons scatter off baryons, they will also feel the acoustic wave and create a peak in the correlation function. The height of the peak is proportional to the amount of baryons: the larger the baryon content of the universe, the higher the peak. The position of the peak in the power spectrum depends on the geometrical size of the particle horizon at last scattering. Since photons travel along geodesics, the projected size of the causal horizon at decoupling depends on whether the universe is flat, open or closed. In a flat universe the geodesics are straight lines and, by looking at the angular scale of the first acoustic peak, we would be measuring the actual size of the horizon at last scattering. In an open universe, the geodesics are inward-curved trajectories, and therefore the projected size on the sky appears smaller. In this case, the first acoustic peak should occur at higher multipoles or smaller angular scales. On the other hand, for a closed universe, the first peak occurs at smaller multipoles or larger angular scales. The dependence of the position of the first acoustic peak on the spatial curvature can be approximately given by $$l_{\mathrm{peak}}220\mathrm{\Omega }_0^{1/2},$$ (211) where $`\mathrm{\Omega }_0=\mathrm{\Omega }_\mathrm{M}+\mathrm{\Omega }_\mathrm{\Lambda }=1\mathrm{\Omega }_K`$. Present observations, specially the ones of the Mobile Anisotropy Telescope (MAT) in Cerro Tololo, Chile, which produced two data sets, TOCO97 and TOCO98 , and the recent balloon-borne experiment BOOMERANG , suggest that the peak is between $`l=180`$ and 250 at 95% c.l., with an amplitude $`\delta T=80\pm 10\mu `$K, and therefore the universe is most probably flat, see Fig. 36, and Ref. . In particular, these measuremts determine that $$0.85\mathrm{\Omega }_01.25(68\%\mathrm{c}.\mathrm{l}.)$$ (212) That is, the universe is flat, within 10% uncertainty, which is much better than we could ever do before. In the near future we will measure $`\mathrm{\Omega }_0`$ to within 1%, with the new microwave anisotropy satellites. At the moment there is not enough information at small angular scales, or large multipole numbers, to determine the existence or not of the secondary acoustic peaks. These peaks should occur at harmonics of the first one, but are typically much lower because of Silk damping. Since the amplitude and position of the primary and secondary peaks are directly determined by the sound speed (and, hence, the equation of state) and by the geometry and expansion of the universe, they can be used as a powerful test of the density of baryons and dark matter, and other cosmological parameters, see Fig. 35. By looking at these patterns in the anisotropies of the microwave background, cosmologists can determine not only the cosmological parameters, see Fig. 35, but also the primordial spectrum of density perturbations produced during inflation. It turns out that the observed temperature anisotropies are compatible with a scale-invariant spectrum, see Eq. (210), as predicted by inflation. This is remarkable, and gives very strong support to the idea that inflation may indeed be responsible for both the CMB anisotropies and the large-scale structure of the universe. Different models of inflation have different specific predictions for the fine details associated with the spectrum generated during inflation. It is these minute differences that will allow cosmologists to differentiate between alternative models of inflation and discard those that do not agree with observations. However, most importantly, perhaps, the pattern of anisotropies predicted by inflation is completely different from those predicted by alternative models of structure formation, like cosmic defects: strings, vortices, textures, etc. These are complicated networks of energy density concentrations left over from an early universe phase transition, analogous to the defects formed in the laboratory in certain kinds of liquid crystals when they go through a phase transition. The cosmological defects have spectral properties very different from those generated by inflation. That is why it is so important to launch more sensitive instruments, and with better angular resolution, to determine the properties of the CMB anisotropies. #### 4.4.5 The new microwave anisotropy satellites, MAP and Planck The large amount of information encoded in the anisotropies of the microwave background is the reason why both NASA and the European Space Agency have decided to launch two independent satellites to measure the CMB temperature and polarization anisotropies to unprecendented accuracy. The Microwave Anisotropy Probe will be launched by NASA at the end of 2000, and Planck is expected in 2007. As we have emphasized before, the fact that these anisotropies have such a small amplitude allow for an accurate calculation of the predicted anisotropies in linear perturbation theory. A particular cosmological model is characterized by a dozen or so parameters: the rate of expansion, the spatial curvature, the baryon content, the cold dark matter and neutrino contribution, the cosmological constant (vacuum energy), the reionization parameter (optical depth to the last scattering surface), and various primordial spectrum parameters like the amplitude and tilt of the adiabatic and isocurvature spectra, the amount of gravitational waves, non-Gaussian effects, etc. All these parameters can now be fed into a fast code called CMBFAST that computes the predicted temperature and polarization anisotropies to 1% accuracy, and thus can be used to compare with observations. These two satellites will improve both the sensitivity, down to $`\mu `$K, and the resolution, down to arc minutes, with respect to the previous COBE satellite, thanks to large numbers of microwave horns of various sizes, positioned at specific angles, and also thanks to recent advances in detector technology, with high electron mobility transistor amplifiers (HEMTs) for frequencies below 100 GHz and bolometers for higher frequencies. The primary advantage of HEMTs is their ease of use and speed, with a typical sensitivity of 0.5 mKs<sup>1/2</sup>, while the advantage of bolometers is their tremendous sensitivity, better than 0.1 mKs<sup>1/2</sup>, see Ref. . For instance, to appreciate the difference, compare the resolution in the temperature anisotropies that COBE and Planck would observe for the same simulated sky in Fig. 37. This will allow cosmologists to extract information from around 3000 multipoles! Since most of the cosmological parameters have specific signatures in the height and position of the first few acoustic peaks, the higher the resolution, the more peaks one is expected to see, and thus the better the accuracy with which one will be able to measure those parameters, see Table 1. As an example of the kind of data that these two satellites will be able to provide, see Fig. 38, which compares the present observational status with that which will become available around 2008. Although the satellite probes were designed for the accurate measurement of the CMB temperature anisotropies, there are other experiments, like balloon-borne and ground interferometers, which will probably accomplish the same results with similar resolution (in the case of MAP), before the satellites start producing their own results . Probably the most important objective of the future satellites will be the measurement of the CMB polarization anisotropies, yet to be discovered. These anisotropies are predicted by models of structure formation and are expected to arise at the level of microKelvin sensitivities, where the new satellites are aiming at. The complementary information contained in the polarization anisotropies will provide much more stringent constraints on the cosmological parameters than from the temperature anisotropies alone. In particular, the curl-curl component of the polarization power spectra is nowadays the only means we have to determine the tensor (gravitational wave) contribution to the metric perturbations responsible for temperature anisotropies, see Fig. 39. If such a component is found, one could constraint very precisely the model of inflation from its spectral properties, specially the tilt . ### 4.5 From metric perturbations to large scale structure If inflation is responsible for the metric perturbations that gave rise to the temperature anisotropies observed in the microwave background, then the primordial spectrum of density inhomogeneities induced by the same metric perturbations should also be responsible for the present large scale structure . This simple connection allows for more stringent tests on the inflationary paradigm for the generation of metric perturbations, since it relates the large scales (of order the present horizon) with the smallest scales (on galaxy scales). This provides a very large lever arm for the determination of primordial spectra parameters like the tilt, the nature of the perturbations, whether adiabatic or isocurvature, the geometry of the universe, as well as its matter and energy content, whether CDM, HDM or mixed CHDM. #### 4.5.1 The galaxy power spectrum As metric perturbations enter the causal horizon during the radiation or matter era, they create density fluctuations via gravitational attraction of the potential wells. The density contrast $`\delta `$ can be deduced from the Einstein equations in linear perturbation theory, see Eq. (142), $$\delta _k\frac{\delta \rho _k}{\rho }=\left(\frac{k}{aH}\right)^2\frac{2}{3}\mathrm{\Phi }_k=\left(\frac{k}{aH}\right)^2\frac{2+2\omega }{5+3\omega }_k,$$ (213) where we have assumed $`K=0`$, and used Eq. (174). From this expression one can compute the power spectrum, at horizon crossing, of matter density perturbations induced by inflation, see Eq. (186), $$P(k)=\|\delta _k\|^2=A\left(\frac{k}{aH}\right)^n,$$ (214) with $`n`$ given by the scalar tilt (188), $`n=1+2\eta 6ϵ`$. This spectrum reduces to a Harrison-Zel’dovich spectrum (57) in the slow-roll approximation: $`\eta ,ϵ1`$. Since perturbations evolve after entering the horizon, the power spectrum will not remain constant. For scales entering the horizon well after matter domination ($`k^1k_{\mathrm{eq}}^181`$ Mpc), the metric perturbation has not changed significantly, so that $`_k(\mathrm{final})=_k(\mathrm{initial})`$. Then Eq. (213) determines the final density contrast in terms of the initial one. On smaller scales, there is a linear transfer function $`T(k)`$, which may be defined as $$_k(\mathrm{final})=T(k)_k(\mathrm{initial}).$$ (215) To calculate the transfer function one has to specify the initial condition with the relative abundance of photons, neutrinos, baryons and cold dark matter long before horizon crossing. The most natural condition is that the abundances of all particle species are uniform on comoving hypersurfaces (with constant total energy density). This is called the adiabatic condition, because entropy is conserved independently for each particle species $`X`$, i.e. $`\delta \rho _X=\dot{\rho }_X\delta t`$, given a perturbation in time from a comoving hypersurface, so $$\frac{\delta \rho _X}{\rho _X+p_X}=\frac{\delta \rho _Y}{\rho _Y+p_Y},$$ (216) where we have used the energy conservation equation for each species, $`\dot{\rho }_X=3H(\rho _X+p_X)`$, valid to first order in perturbations. It follows that each species of radiation has a common density contrast $`\delta _r`$, and each species of matter has also a common density contrast $`\delta _m`$, with the relation $`\delta _m=\frac{3}{4}\delta _r`$. Within the horizon, the density perturbation amplitude evolves according to the following equation, see Ref. , $$H^2\ddot{\delta }_k+[23(2\omega c_s^2)]H^1\dot{\delta }_k\frac{3}{2}(16c_s^2+8\omega 3\omega ^2)\delta _k=\left(\frac{k}{aH}\right)^2\frac{\delta p_k}{\rho },$$ (217) where $`\omega =p/\rho `$ is the barotropic ratio, and $`c_s^2=\dot{p}/\dot{\rho }`$ is the speed of sound of the fluid. Given the adiabatic condition, the transfer function is determined by the physical processes occuring between horizon entry and matter domination. If the radiation behaves like a perfect fluid, its density perturbation oscillates during this era, with decreasing amplitude. The matter density contrast living in this background does not grow appreciably before matter domination because it has negligible self-gravity. The transfer function is therefore given roughly by, see Eq. (60), $$T(k)=\{\begin{array}{cc}1,\hfill & kk_{\mathrm{eq}}\hfill \\ (k/k_{\mathrm{eq}})^2,\hfill & kk_{\mathrm{eq}}\hfill \end{array}$$ (218) The perfect fluid description of the radiation is far from being correct after horizon entry, because roughly half of the radiation consists of neutrinos whose perturbation rapidly disappears through free streeming. The photons are also not a perfect fluid because they diffuse significantly, for scales below the Silk scale, $`k_S^11`$ Mpc. One might then consider the opposite assumption, that the radiation has zero perturbation after horizon entry. Then the matter density perturbation evolves according to Eq. (217), with $`\delta `$ and $`\rho `$ now referring to the matter alone, $$\ddot{\delta }_k+2H\dot{\delta }_k+(c_s^2k_{\mathrm{ph}}^24\pi G\rho )\delta _k=0,$$ (219) which corresponds to the equation of a damped harmonic oscillator. The zero-frequency oscillator defines the Jeans wavenumber, $`k_J=\sqrt{4\pi G\rho /c_s^2}`$. For $`kk_J`$, $`\delta _k`$ grows exponentially on the dynamical timescale, $`\tau _{\mathrm{dyn}}=\mathrm{Im}\omega ^1=(4\pi G\rho )^{1/2}=\tau _{\mathrm{grav}}`$, which is the time scale for gravitational collapse. One can also define the Jeans length, $$\lambda _J=\frac{2\pi }{k_J}=c_s\sqrt{\frac{\pi }{G\rho }},$$ (220) which separates gravitationally stable from unstable modes. If we define the pressure response timescale as the size of the perturbation over the sound speed, $`\tau _{\mathrm{pres}}\lambda /c_s`$, then, if $`\tau _{\mathrm{pres}}>\tau _{\mathrm{grav}}`$, gravitational collapse of a perturbation can occur before pressure forces can response to restore hydrostatic equilibrium (this occurs for $`\lambda >\lambda _J`$). On the other hand, if $`\tau _{\mathrm{pres}}<\tau _{\mathrm{grav}}`$, radiation pressure prevents gravitational collapse and there are damped acoustic oscillations (for $`\lambda <\lambda _J`$). We will consider now the behaviour of modes within the horizon during the transition from the radiation ($`c_s^2=1/3`$) to the matter era ($`c_s^2=0`$). The growing and the decaying solutions of Eq. (219) are $`\delta =A\left(1+{\displaystyle \frac{3}{2}}y\right),`$ (222) $`\delta =B\left[\left(1+{\displaystyle \frac{3}{2}}y\right)\mathrm{ln}{\displaystyle \frac{\sqrt{1+y}+1}{\sqrt{1+y}1}}3\sqrt{1+y}\right],`$ where $`A`$ and $`B`$ are constants, and $`y=a/a_{\mathrm{eq}}`$. The growing mode solution (222) increases only by a factor of 2 between horizon entry and the epoch when matter starts to dominate, i.e. $`y=1`$. The transfer function is therefore again roughly given by Eq. (218). Since the radiation consists roughly half of neutrinos, which free streem, and half of photons, which either form a perfect fluid or just diffuse, neither the perfect fluid nor the free-streeming approximation looks very sensible. A more precise calculation is needed, including: neutrino free streeming around the epoch of horizon entry; the diffusion of photons around the same time, for scales below Silk scale; the diffusion of baryons along with the photons, and the establishment after matter domination of a common matter density contrast, as the baryons fall into the potential wells of cold dark matter. All these effects apply separately, to first order in the perturbations, to each Fourier component, so that a linear transfer function is produced. There are several parametrizations in the literature, but the one which is more widely used is that of Ref. , $`T(k)=\left[1+\left(ak+(bk)^{3/2}+(ck)^2\right)^\nu \right]^{1/\nu },\nu =1.13,`$ (223) $`a=6.4(\mathrm{\Omega }_\mathrm{M}h)^1h^1\mathrm{Mpc},`$ (224) $`b=3.0(\mathrm{\Omega }_\mathrm{M}h)^1h^1\mathrm{Mpc},`$ (225) $`c=1.7(\mathrm{\Omega }_\mathrm{M}h)^1h^1\mathrm{Mpc}.`$ (226) We see that the behaviour estimated in Eq. (218) is roughly correct, although the break at $`k=k_{\mathrm{eq}}`$ is not at all sharp, see Fig. 40. The transfer function, which encodes the soltion to linear equations, ceases to be valid when the density contrast becomes of order 1. After that, the highly nonlinear phenomenon of gravitational collapse takes place, see Fig. 40. #### 4.5.2 The new redshift catalogs, 2dF and Sloan Digital Sky Survey Our view of the large-scale distribution of luminous objects in the universe has changed dramatically during the last 25 years : from the simple pre-1975 picture of a distribution of field and cluster galaxies, to the discovery of the first single superstructures and voids, to the most recent results showing an almost regular web-like network of interconnected clusters, filaments and walls, separating huge nearly empty volumes. The increased efficiency of redshift surveys, made possible by the development of spectrographs and – specially in the last decade – by an enormous increase in multiplexing gain (i.e. the ability to collect spectra of several galaxies at once, thanks to fibre-optic spectrographs), has allowed us not only to do cartography of the nearby universe, but also to statistically characterize some of its properties, see Ref. . At the same time, advances in theoretical modeling of the development of structure, with large high-resolution gravitational simulations coupled to a deeper yet limited understanding of how to form galaxies within the dark matter halos, have provided a more realistic connection of the models to the observable quantities . Despite the large uncertainties that still exist, this has transformed the study of cosmology and large-scale structure into a truly quantitative science, where theory and observations can progress side by side. For a review of the variety and details about the different existing redshift catalogs, see Ref. , and Fig. 41. Here I will concentrate on two of the new catalogs, which are taking data at the moment and which will revolutionize the field, the 2-degree-Field (2dF) Catalog and the Sloan Digital Sky Survey (SDSS). The advantages of multi-object fibre spectroscopy have been pushed to the extreme with the construction of the 2dF spectrograph for the prime focus of the Anglo-Australian Telescope . This instrument is able to accommodate 400 automatically positioned fibres over a 2 degree in diameter field. This implies a density of fibres on the sky of approximately 130 deg<sup>-2</sup>, and an optimal match to the galaxy counts for a magnitude $`b_J19.5`$, similar to that of previous surveys like the ESP, with the difference that with such an area yield, the same number of redshifts as in the ESP survey can be collected in about 10 exposures, or slightly more than one night of telescope time with typical 1 hour exposures. This is the basis of the 2dF galaxy redshift survey. Its goal is to measure redshifts for more than 250,000 galaxies with $`b_J<19.5`$. In addition, a faint redshift survey of 10,000 galaxies brighter than $`R=21`$ will be done over selected fields within the two main strips of the South and North Galactic Caps. The survey is steadily collecting redshifts, and there were about 93,000 galaxies measured by January 2000. See also Ref. , where the survey is continuously updated. The most ambitious and comprehensive galaxy survey currently in progress is without any doubt the Sloan Digital Sky Survey . The aim of the project is first of all to observe photometrically the whole Northern Galactic Cap, 30 away from the galactic plane (about $`10^4`$ deg<sup>2</sup>) in five bands, at limiting magnitudes from 20.8 to 23.3. The expectation is to detect around 50 million galaxies and around $`10^8`$ star-like sources. This has already led to the discovery of several high-redshift ($`z>4`$) quasars, including the highest-redshift quasar known, at $`z=5.0`$, see Ref. . Using two fibre spectrographs carrying 320 fibres each, the spectroscopic part of the survey will then collect spectra from about $`10^6`$ galaxies with $`r^{}<18`$ and $`10^5`$ AGNs with $`r^{}<19`$. It will also select a sample of about $`10^5`$ red luminous galaxies with $`r^{}<19.5`$, which will be observed spectroscopically, providing a nearly volume-limited sample of early-type galaxies with a median redshift of $`z0.5`$, that will be extremely valuable to study the evolution of clustering. The data expected to arise from these new catalogs is so outstanding that already cosmologists are making simulations and predicting what will be the scientific outcome of these surveys, together with the future CMB anisotropy probes, for the determination of the cosmological parameters of the standard model of cosmology, see Figs. 41 and 42. As often happens in particle physics, not always are observations from a single experiment sufficient to isolate and determine the precise value of the parameters of the standard model. We mentioned in the previous Section that some of the cosmological parameters created similar effects in the temperature anisotropies of the microwave background. We say that these parameters are degenerate with respect to the observations. However, often one finds combinations of various experiments/observations which break the degeneracy, for example by depending on a different combination of parameters. This is precisely the case with the cosmological parameters, as measured by a combination of large-scale structure observations, microwave background anisotropies, Supernovae Ia observations and Hubble Space Telescope measurements, a feature named somewhat idiosyncratically as “cosmic complementarity”, see Ref. . It is expected that in the near future we will be able to determine the parameters of the standard cosmological model with great precision from a combination of several different experiments, as shown in Fig. 42. ## 5 CONCLUSION We have entered a new era in cosmology, were a host of high-precision measurements are already posing challenges to our understanding of the universe: the density of ordinary matter and the total amount of energy in the universe; the microwave background anisotropies on a fine-scale resolution; primordial deuterium abundance from quasar absorption lines; the acceleration parameter of the universe from high-redshift supernovae observations; the rate of expansion from gravitational lensing; large scale structure measurements of the distribution of galaxies and their evolution; and many more, which already put constraints on the parameter space of cosmological models, see Fig. 30. However, these are only the forerunners of the precision era in cosmology that will dominate the new millennium, and will make cosmology a phenomenological science. It is important to bear in mind that all physical theories are approximations of reality that can fail if pushed too far. Physical science advances by incorporating earlier theories that are experimentally supported into larger, more encompassing frameworks. The standard Big Bang theory is supported by a wealth of evidence, nobody really doubts its validity anymore. However, in the last decade it has been incorporated into the larger picture of cosmological inflation, which has become the new standard cosmological model. All cosmological issues are now formulated in the context of the inflationary paradigm. It is the best explanation we have at the moment for the increasing set of cosmological observations. In the next few years we will have an even larger set of high-quality observations that will test inflation and the cold dark matter paradigm of structure formation, and determine most of the 12 or more parameters of the standard cosmological model to a few percent accuracy (see table 1). It may seem that with such a large number of parameters one can fit almost anything. However, that is not the case when there is enough quantity and quality of data. An illustrative example is the standard model of particle physics, with around 21 parameters and a host of precise measurements from particle accelerators all over the world. This model is, nowadays, rigurously tested, and its parameters measured to a precision of better than 1% in some cases. It is clear that high-precision measurements will make the standard model of cosmology as robust as that of particle physics. In fact, it has been the technological advances of particle physics detectors that are mainly responsible for the burst of new data coming from cosmological observations. This is definitely a very healthy field, but there is still a lot to do. With the advent of better and larger precision experiments, cosmology is becoming a mature science, where speculation has given way to phenomenology. There are still many unanswered fundamental questions in this emerging picture of cosmology. For instance, we still do not know the nature of the inflaton field, is it some new fundamental scalar field in the electroweak symmetry breaking sector, or is it just some effective description of a more fundamental high energy interaction? Hopefully, in the near future, experiments in particle physics might give us a clue to its nature. Inflation had its original inspiration in the Higgs field, the scalar field supposed to be responsible for the masses of elementary particles (quarks and leptons) and the breaking of the electroweak symmetry. Such a field has not been found yet, and its discovery at the future particle colliders would help understand one of the truly fundamental problems in physics, the origin of masses. If the experiments discover something completely new and unexpected, it would automatically affect inflation at a fundamental level. One of the most difficult challenges that the new cosmology will have to face is understanding the origin of the cosmological constant, if indeed it is confirmed by independent sets of observations. Ever since Einstein introduced it as a way to counteract gravitational attraction, it has haunted cosmologists and particle physicists for decades. We still do not have a mechanism to explain its extraordinarily small value, 120 orders of magnitude below what is predicted by quantum physics. For several decades there has been the reasonable speculation that this fundamental problem may be related to the quantization of gravity. General relativity is a classical theory of space-time, and it has proved particularly difficult to construct a consistent quantum theory of gravity, since it involves fundamental issues like causality and the nature of space-time itself. The value of the cosmological constant predicted by quantum physics is related to our lack of understanding of gravity at the microscopic level. However, its effect is dominant at the very largest scales of clusters or superclusters of galaxies, on truly macroscopic scales. This hints at what is known in quantum theory as an anomaly, a quantum phenomenon relating both ultraviolet (microscopic) and infrared (macroscopic) divergences. We can speculate that perhaps general relativity is not the correct description of gravity on the very largest scales. In fact, it is only in the last few billion years that the observable universe has become large enough that these global effects could be noticeable. In its infancy, the universe was much smaller than it is now, and, presumably, general relativity gave a correct description of its evolution, as confirmed by the successes of the standard Big Bang theory. As it expanded, larger and larger regions were encompassed, and, therefore, deviations from general relativity would slowly become important. It may well be that the recent determination of a cosmological constant from observations of supernovae at high redshifts is hinting at a fundamental misunderstanding of gravity on the very large scales. If this were indeed the case, we should expect that the new generation of precise cosmological observations will not only affect our cosmological model of the universe but also a more fundamental description of nature. ## ACKNOWLEDGEMENTS I thank the organizers of the CERN-JINR European School of High Energy Physics for a very warm and friendly atmosphere. I also would like to thank my friends and collaborators Andrei Linde, Andrew Liddle, David Wands, David Lyth, Jaume Garriga, Xavier Montes, Enrique Gaztañaga, Elena Pierpaoli, Stefano Borgani, and many others, for sharing with me their insight about this fascinating science of cosmology. This work was supported by the Royal Society.
warning/0004/math0004001.html	ar5iv	text	# An asymptotic expansion for a ratio of products of gamma functions ## 1 Introduction Our starting point is the Gaussian hypergeometric function $`F(a,b;c;z)`$ and its series representation $$\frac{1}{\mathrm{\Gamma }(c)}F(a,b;c;z)=\underset{n=0}{\overset{\mathrm{}}{}}\frac{(a)_n(b)_n}{\mathrm{\Gamma }(c+n)n!}z^n,\|z\|<1,$$ which here is written in terms of Pochhammer symbols $$(x)_n=x(x+1)\mathrm{}(x+n1)=\mathrm{\Gamma }(x+n)/\mathrm{\Gamma }(x).$$ The hypergeometric series appears as one solution of the Gaussian (or hypergeometric) differential equation, which is characterized by its three regular singular points at $`z=0,1,\mathrm{}`$. The local series solutions at $`0`$ and $`1`$ of this differential equation are connected by the continuation formula $$\frac{1}{\mathrm{\Gamma }(c)}F(a,b;c;z)=\frac{\mathrm{\Gamma }(cab)}{\mathrm{\Gamma }(ca)\mathrm{\Gamma }(cb)}F(a,b;1+a+bc;1z)$$ $$+\frac{\mathrm{\Gamma }(a+bc)}{\mathrm{\Gamma }(a)\mathrm{\Gamma }(b)}(1z)^{cab}F(ca,cb;1+cab;1z),$$ (1) $$(\|\mathrm{arg}(1z)\|<\pi ).$$ Here we want to show that Eq. (1) implies an interesting asymptotic expansion for a ratio of products of gamma functions, of which only a special case was known before. By applying the method of Darboux to (1), we derive in Sec. 2 the formula in question. The behaviour of this and a related formula is discussed in Sec. 3 and illustrated by a few numerical examples. ## 2 Derivation of an asymptotic expansion for a ratio of products of gamma functions It is well-known that the late coefficients of a Taylor series expansion contain information about the nearest singular point of the expanded function . In this respect we want to analyze the continuation formula (1), in which then only the second, at $`z=1`$ singular term $`R`$ is relevant, which may be written as $$R=\frac{\mathrm{\Gamma }(a+bc)\mathrm{\Gamma }(1+cab)}{\mathrm{\Gamma }(a)\mathrm{\Gamma }(b)}\underset{m=0}{\overset{\mathrm{}}{}}\frac{(ca)_m(cb)_m}{\mathrm{\Gamma }(1+cab+m)m!}(1z)^{cab+m}.$$ By means of the binomial theorem in its hypergeometric-series-form , we may expand the power factor $$(1z)^{cab+m}=\underset{n=0}{\overset{\mathrm{}}{}}\frac{\mathrm{\Gamma }(a+bcm+n)}{\mathrm{\Gamma }(a+bcm)n!}z^n.$$ Interchanging then the order of the summations and simplifying by means of the reflection formula of the gamma function, we arrive at $$R=\frac{1}{\mathrm{\Gamma }(a)\mathrm{\Gamma }(b)}\underset{n=0}{\overset{\mathrm{}}{}}\underset{m=0}{\overset{\mathrm{}}{}}(1)^m\frac{(ca)_m(cb)_m}{m!}\frac{\mathrm{\Gamma }(a+bcm+n)}{n!}z^n.$$ This is to be compared with the left-hand side $`L`$ of (1), which is $$L=\frac{1}{\mathrm{\Gamma }(a)\mathrm{\Gamma }(b)}\underset{n=0}{\overset{\mathrm{}}{}}\frac{\mathrm{\Gamma }(a+n)\mathrm{\Gamma }(b+n)}{\mathrm{\Gamma }(c+n)n!}z^n.$$ Comparison of the coefficients of these two power series, which according to Darboux and Schäfke and Schmidt should agree asymptotically as $`n\mathrm{}`$, then yields $$\frac{\mathrm{\Gamma }(a+n)\mathrm{\Gamma }(b+n)}{\mathrm{\Gamma }(c+n)}=\underset{m=0}{\overset{M}{}}(1)^m\frac{(ca)_m(cb)_m}{m!}\mathrm{\Gamma }(a+bcm+n)$$ (2) $$+O(\mathrm{\Gamma }(a+bcM1+n)).$$ By means of $$O(\mathrm{\Gamma }(a+bcM1+n))=\mathrm{\Gamma }(a+bc+n)O(n^{M1})$$ and the reflection formula of the gamma function, the relevant formula (2) may also be written as $$\frac{\mathrm{\Gamma }(a+n)\mathrm{\Gamma }(b+n)}{\mathrm{\Gamma }(c+n)\mathrm{\Gamma }(a+bc+n)}=1+\underset{m=1}{\overset{M}{}}\frac{(ca)_m(cb)_m}{m!(1+cabn)_m}+O(n^{M1}).$$ (3) The asymptotic expansion for a ratio of products of gamma functions in this form (3) or the other (2) seems to be new. It is only the special case when $`c=1`$ which is known. This special case was stated by Dingle, first proved by Paris, and reconsidered recently by Olver, who has found a simple direct proof. His proof, as well as the proof of Paris, can be adapted easily to the more general case when $`c`$ is different from $`1`$ . Still another proof is available which includes an integral representation of the remainder term. Our derivation of Eq. (2) or (3) is significantly different from all the earlier proofs of the case when $`c=1`$. ## 3 Discussion and numerical examples We now want to discuss our result in the form (3). First we observe that the substitution $`ca+bc`$ leads to the related formula $$\frac{\mathrm{\Gamma }(a+n)\mathrm{\Gamma }(b+n)}{\mathrm{\Gamma }(c+n)\mathrm{\Gamma }(a+bc+n)}=1+\underset{m=1}{\overset{M}{}}\frac{(ac)_m(bc)_m}{m!(1cn)_m}+O(n^{M1}).$$ (4) Which of (3) or (4) is more advantageous numerically depends on the values of the parameters, and in this respect the two formulas complement each other. Table 1 shows an example with a set of parameters for which (3) gives more accurate values than (4), while Table 2 contains an example for which (4) is superior to (3). For finite $`n`$ and $`M\mathrm{}`$ the series on the right-hand side of (3) converges if $`\text{Re}(1cn)>0`$. The same is true for (4) if $`\text{Re}(1+cabn)>0`$ . Then, in both cases, the Gaussian summation formula yields $$\frac{\mathrm{\Gamma }(1cn)\mathrm{\Gamma }(1+cabn)}{\mathrm{\Gamma }1an)\mathrm{\Gamma }(1bn)},$$ which, by means of the reflection formula of the gamma function, is seen to be equal to $$\frac{\mathrm{\Gamma }(a+n)\mathrm{\Gamma }(b+n)}{\mathrm{\Gamma }(c+n)\mathrm{\Gamma }(a+bc+n)}\frac{\mathrm{sin}(\pi [a+n])\mathrm{sin}(\pi [b+n])}{\mathrm{sin}(\pi [c+n])\mathrm{sin}(\pi [a+bc+n])}.$$ (5) Otherwise (2) – (4) are divergent asymptotic expansions as $`n\mathrm{}`$. Although in our derivation $`n`$ is a sufficiently large positive integer, the asymptotic expansions (2) – (4) are expected to be valid in a certain sector of the complex $`n`$ -plane, and in fact, the proofs of Paris and of Olver apply to complex values of $`n`$. If the series in (3) or (4) converge, their sums are equal to (5), which generally (if neither $`ca`$ nor $`cb`$ is equal to an integer ) is different from the left-hand side of (3) or (4). Therefore (3) and (4) can be valid only in the half-planes in which the series do not converge. This means that (3) is an asymptotic expansion as $`n\mathrm{}`$ in the half-plane $`\text{Re}(c1+n)0`$, and (4) is an asymptotic expansion as $`n\mathrm{}`$ in the half-plane $`\text{Re}(a+bc1+n)0`$. Otherwise the series on the right-hand sides represent a different function, namely (5). A few numerical examples may serve for demonstration of these facts. In Table 3 , the series converge to (5) for $`n=10`$ , and therefore (3) and (4) are not valid. For $`n=20`$, on the other hand, the series diverge and so (3) and (4) hold. The transition between the two regions is at the line $`\text{Re}(n)=12.4`$ in case of (3) or $`\text{Re}(n)=12.5`$ in case of (4). In Table 4, we see convergence for $`n=15`$ and divergence for $`n=5`$, the transition between the two regions being at the line $`\text{Re}(n)=10.4`$ in case of (3) or $`\text{Re}(n)=10.5`$ in case of (4) .
warning/0004/cond-mat0004349.html	ar5iv	text	# Evidence for an energy scale for quasiparticle dispersion in 𝐵⁢𝑖₂⁢𝑆⁢𝑟₂⁢𝐶⁢𝑎⁢𝐶⁢𝑢₂⁢𝑂₈ ## Abstract Quasiparticle dispersion in $`Bi_2Sr_2CaCu_2O_8`$ is investigated with improved angular resolution as a function of temperature and doping. Unlike the linear dispersion predicted by the band calculation, the data show a sharp break in dispersion at $`50\pm 10`$ $`meV`$ binding energy where the velocity changes by a factor of two or more. This change provides an energy scale in the quasiparticle self-energy. This break in dispersion is evident at and away from the d-wave node line, but the magnitude of the dispersion change decreases with temperature and with increasing doping. In a conventional metal the observation of an energy scale often provides significant insight into the physical process in the material. The most noted example is the observation of the phonon anomalies in strong coupling superconductors such as lead, which had a far-reaching impact on the understanding of the superconductivity mechanism . For the high-temperature superconductors, a peculiar normal state property is the fact that there appears to be no energy scale, which is often referred to as the marginal Fermi liquid behavior . This behavior is highly anomalous as one would expect certain energy scales in the problem, say phonons which are obviously present in the crystal. In the theoretical context this lack of an energy scale is believed to be a key feature of a near-by quantum critical point . In the superconducting state, on the other hand, there are energy scales observed in the cuprate superconductors. One of them is the superconducting gap and the other is the so-called 41 meV magnetic resonance . The latter has been attributed to the superconducting gap, or the $`\pi `$-resonance of the SO(5) theory . With its ability to measure both the real and imaginary parts of the self-energy, $`\mathrm{\Sigma }(\omega ,k)`$, angle-resolved photoemission (ARPES) experiments provide a unique opportunity to further explore this issue as any relevant energy scale present will manifest itself in the quasiparticle dynamics. In the known case of electron-phonon interaction the coupling causes a kink in the dispersion and also a change in quasiparticle lifetime near the phonon energy . These canonical changes reveal effects in the real and imaginary parts of the self-energy due to the electron-phonon interaction, an effect which is experimentally observed recently . In this letter, we present high-resolution ARPES data from $`Bi_2Sr_2CaCu_2O_8`$ superconductors as a function of doping and temperature. We have observed a clear break in the quasiparticle dispersion near $`50\pm 10`$ $`meV`$ binding energy (BE), that results in a change in the quasiparticle velocity up to a factor of two or more. This effect is enhanced in the underdoped sample, and appears to persists above $`T_C`$ where the break becomes rather broad. Because the electronic structure calculation predicts a linear dispersion in this range, this result represents an important effect in the real part of the self-energy with a scale near $`50\pm 10`$ $`meV`$. Further, we found that this effect is present at various points of the momentum space. We believe the doping, temperature and $`\stackrel{}{k}`$-dependent information presented here will put a constrain on microscopic theory. Angle-resolved photoemission data have been recorded at beamline $`\mathrm{10.0.1.1}`$ of the Advanced Light Source utilizing $`22`$ $`eV`$, $`33`$ $`eV`$ and $`55`$ $`eV`$ photon energies, in a similar set-up as we have reported recently . The momentum resolution was $`\pm 0.1`$ degrees, which is about an order of magnitude better than our previous study of this material, making the results reported in this letter possible. The energy resolution was $`14`$ $`meV`$. The vacuum during the measurement was better than $`410^{}{}_{}{}^{1}^1`$ $`torr`$. The underdoped (UD) $`Bi_2Sr_2CaCu_2O_8`$ ($`Tc=84K`$) and the slightly overdoped (OD) $`Bi_2Sr_2CaCu_2O_8`$ sample ($`Tc=91K`$) were grown using floating-zone method. The single crystalline samples were oriented and cleaved in situ at low temperature. Fig.1a) shows raw ARPES data collected along the $`(0,0)`$ to $`(\pi ,\pi )`$ (nodal) direction of the Brillouin zone from the OD sample at 30K. In panel 1b) we plot the dispersion determined from the fits to the momentum distribution curves (MDCs) - angle scans at a constant binding energy . MDC plots show a peak on a constant background that can be fitted very well with a simple Lorentzian, as illustrated in the inset b2). Error bars in $`k_{}`$ and energy are determined from the fit uncertainty and energy resolution respectively. The data clearly show a feature dispersing towards the Fermi energy with an obvious break in the slope near $`50`$ $`meV`$ BE. A similar break in the dispersion was also observed at photon energies $`22`$ $`eV`$ and $`55`$ $`eV`$. Data for all three photon energies is plotted in the inset b1) in panel b). To describe the dispersion in the range of ($`200`$ $`meV`$ to $`0`$ $`meV`$) one needs only two straight lines intersecting near $`50`$ $`meV`$ BE. This behavior is clearly different from what one expects from the LDA or any other electronic structure calculation where a linear dispersion in this energy range is predicted. Raw MDCs are plotted in panel 1c), while raw energy distribution curves (EDCs) are plotted in panel 1 d) for reference. We present in Fig. 2 (a-c) the dispersions obtained from different cuts parallel to the $`(0,0)`$ to $`(\pi ,\pi )`$ direction across the Fermi surface for the UD sample at 20K. Within the error bars, the data are again well described by two straight lines with a break near $`50`$ $`meV`$ BE. The energy position of the break is constant throughout the BZ within the experimental uncertainty, despite the opening of the gap. Fig. 3 shows the locations in the two-dimensional zone where the break is experimentally observed. It demonstrates that the effect is present in all directions. We have investigated this effect as a function of doping and temperature. The effect appears to be stronger in the underdoped sample. The change of the quasiparticle velocity at the break is different, which can be illustrated by data along the $`(0,0)`$ to $`(\pi ,\pi )`$ direction. For the underdoped sample, the quasiparticle velocity determined from the MDC fits shows a break from $`3.6`$ $`eV\AA `$ at higher binding energies to $`1.5`$ $`eV\AA `$ near the Fermi level. The respective velocities for the optimally doped samples are $`2.6`$ $`eV\AA `$ and $`1.6`$ $`eV\AA `$. The error in velocities from the fits is $`\pm .1eV\AA `$. Main source of error is the possible misalignment , thus causing some uncertainty in quantitative results. However, the general trend discussed above is robust. In general one expects to see complementary effects in dispersion and EDC and MDC peak widths as they reflect the quasiparticle self-energy. The self-energy can be easily extracted from an ARPES experiment if $`Im\mathrm{\Sigma }(\stackrel{}{k},E)`$ is much smaller than the energy. In this case MDC and EDC methods give the same result for the peak position and for the peak width interpretable as $`Re`$ and $`Im`$ parts of the self energy respectively. In high Tc’s extracting the self energy from ARPES is harder because the EDC peak energy is comparable to the peak width for $`E30meVBE`$. However, assuming weak k-dependence of the $`Im\mathrm{\Sigma }(\stackrel{}{k},E)`$ , the deviation of the MDC dispersion from the LDA calculation gives real part of the self energy and the MDC peak widths represent imaginary part . Fig.2 (a1-c1) shows MDC widths in momentum space along various cuts. The corresponding energy width is given by the momentum width of the MDC peak multiplied by the velocity if the scan direction is along the energy gradient direction. In our geometry this condition is satisfied only along the nodal direction. In Figure 2 d) we plot the energy width from MDC together with EDC width. The step effect in MDC energy width is due to linear approximation to the dispersion in determining the velocity, smoother transition is expected for a less dramatic behavior. EDC peak widths do not simply give $`Im\mathrm{\Sigma }(\stackrel{}{k},E)`$ in the case of broad peaks. Furthermore, EDC data is complicated by energy-dependent background and Fermi cut-off. Nevertheless, the EDC data still indicate a more abrupt change in the width at the energies corresponding to the kink in the dispersion, as shown in Fig. 2 (d-f). We feel that the qualitative consistency between MDC and EDC results is sufficient to make the case for the strong self energy effect in the data. The lack of quantitative agreement between EDC and MDC is a manifestation of the subtle lineshape issues discussed above. The dispersions determined from the OD sample above $`T_C`$ along the $`(0,0)`$ to $`(\pi ,\pi )`$ ($`\mathrm{\Gamma }Y`$) direction are shown in Fig. 4a), while the low temperature dispersions ($`T<T_c`$) are reported in 4b) . The dispersions exhibit the same break structure as contrasted to the straight line. The change of dispersion is more difficult to see in the high temperature data compared to low temperature data, but a weak residual effect still appears to be present. In Fig. 4c)we show the temperature dependence of the EDC width. We see a clear change in $`2\mathrm{\Gamma }`$ around $`50\pm 10meV`$ in the low temperature data, but the effect is harder to see above Tc. We now discuss the origin of the strong self-energy effect near $`50\pm 10`$ $`meV`$. The first possibility that comes to mind is the electron-phonon interaction as there are phonons of this energy scale in the compound . This would explain the persistence of the feature throughout the Brillouin zone and the persistence of the feature above the superconducting transition temperature, since phonons are $`\stackrel{}{k}`$ and T independent. However, as shown recently , the dispersion tends to recover to the one-electron result when the energy is well above the typical phonon energies, with the total range of the perturbed dispersion below $`E_F`$ equal to half the Debye temperature. In our case, there is no indication that the dispersion will recover to LDA behavior and the high energy part of the data cannot be fit with a line passing through the Fermi surface crossing. Another problem with the phonon scenario is that it is not a natural explanation on why the effect is stronger in underdoped sample. The second and in our view most likely possibility is the electron coupling to collective magnetic excitations . The neutron mode is found below $`T_C`$ near $`(\pi ,\pi )`$ at $`41`$ $`meV`$ with a width of $`0.6\pi `$ in $`Bi_2Sr_2CaCu_2O_8`$ system . The energy scale of the neutron mode is consistent with the $`50\pm 10`$ $`meV`$ feature seen in our experiments. This picture is also consistent with the fact that the underdoped sample shows stronger effect than the overdoped one . The broadness of the neutron peak in this compound makes it possible to explain the persistence of the effect throughout the Brillouin zone. Mode’s energy decrease in overdoped sample and the fact that the mode is mainly seen below $`T_C`$ also supports this interpretation of our data. However, we caution that the exact temperature where this mode turns on is a sensitive function of doping and impurities . Of course, we cannot rule out the possibility of a combination of effects due to phonons and magnetic excitations because both have similar energy scales. We should note that the coupling of a quasiparticle to collective excitations was previously discussed, but in a very different context . The break in quasiparticle dispersion of optimally doped sample along the nodal direction is also present in the data of Valla $`etal.`$. However, these authors did not elaborate on this issue and they suggested the absence of energy scale in the problem. The third possibility is that we see the effect related to the opening of the superconducting gap which is of the order of $`50\pm 10`$ $`meV`$ in this compound. Because the anti-nodal direction near $`(\pi ,0)`$ has a very high density of states, we expect an effect when this energy scale is reached . The down side of this scenario is that we do not have theoretical calculations on specifics of the effect. At this point, we do not feel that we can rule out possible interpretation of the data based on the stripe scenario where the 50meV kink in the dispersion reflects the characteristic frequency of the fluctuating stripes. In summary, we have studied doping, momentum and temperature dependence of the quasiparticle dispersion using ARPES with very high momentum resolution. We have uncovered an energy scale of $`50\pm 10`$ $`meV`$ where the quasiparticle dispersion shows a strong change. This effect is seen in all directions, and is stronger in the underdoped sample. The effect is strongest in data below $`T_C`$. We expect that the data presented here is an important part of the puzzle related to high$`T_C`$ superconductivity. We would like to thank J. D. Denlinger for the help with data analisys software. We would like to thank P. D. Johnson,B. O. Wells, A. Fedorov, T. Valla, S.C. Zhang, D.J. Scalapino, Steve Kivelson, D.H. Lee, Bob Laughlin, and P.A. Lee for discussion. The experiment was performed at the Advanced Light Source of Lawrence Berkekely National Laboratory. The Stanford work was supported by NSF grant through the Stanford MRSEC grant and NSF grant DMR-9705210. The work at ALS was supported by DOE’ Office of Basic Energy Science, Division of Materials Science with contract DE-AC03-76SF00098. The SSRL’s work was also supported by the Office’s Division of Materials Science.
warning/0004/quant-ph0004001.html	ar5iv	text	# Interference-induced gain in Autler-Townes doublet of a V-type atom in a cavity ## Abstract We study the Autler-Townes spectrum of a V-type atom coupled to a single-mode, frequency-tunable cavity field at finite termperature, with a pre-selected polarization in the bad cavity limit, and show that, when the mean number of thermal photons $`N1`$ and the excited sublevel splitting is very large (the same order as the cavity linewidth), the probe gain may occur at either sideband of the doublet, depending on the cavity frequency, due to the cavity-induced interference. Within recent years, there has been a resurgence of interest in the phenomenon of quantum interference . The principal reason is that it lies at the heart of many new effects and applications of quantum optics, such as lasing without population inversion , electromagnetically-induced transparency , enhancement of the index of refraction without absorption , fluorescence quenching and spectral line narrowing . The basic system consists of a singlet state connected to a closely-spaced excited doublet by a single-mode laser. Cardimona et al. studied the effect of quantum interference on the resonance fluorescence of such a system, and found that it can be driven into a dark state in which quantum interference prevents any fluorescence from the excited sublevels, regardless of the intensity of the exciting laser. We have recently shown that quantum interference can also lead to narrow resonances, transparency and gain without population inversion in the probe absorption spectrum of such an atomic system . Harris and co-workers generalized the V-type atom to systems where the excited doublets decay to an additional continuum or to a single auxiliary level, in addition to the ground state. They found that at a certain frequency, the absorption rate goes to zero due to destructive interference whereas the emission rate remains finite. It is possible to amplify a laser field at this frequency without population inversion being present. In the case of a single auxiliary level, quantum interference can lead to the elimination of the spectral line at the driving laser frequency in the spontaneous emission spectrum and transparency in the absorption spectrum . It is important for these effects that the dipole moments of the transitions involved are parallel, so that the cross-decay terms are maximal. From the experimental point view, however, it is difficult to find isolated atomic systems which have parallel moments . Various alternative proposals have been made for generating quantum interference effects. For example, if the two upper levels of a V-type atom are coupled by a microwave field or an applied laser, the excited doublet becomes a superposition, so that as the atom decays from one of the excited sublevels it drives the other. For such systems, the cross-decay terms are evident in the atomic dressed picture . A four-level atom with two closely-spaced intermediate states coupled to a two-mode cavity can also show the effect of quantum interference . In fact, the experimental observation of the interference-induced suppression of spontaneous emission was carried out in sodium dimers where the excited sublevels are superpositions of singlet and triplet states that are mixed by a spin-orbit interaction . We have recently also proposed a scheme for engineering of quantum interference (parallel or anti-parallel dipole moments) in a V-type atom coupled to a frequency tunable, single-mode cavity field with a pre-selected polarization at zero temperature . We have found that the effects of the cavity-induced interference are pronounced only when the cavity detuning $`\delta `$ and the excited doublet splitting $`\omega _{21}`$ are much less than the cavity linewidth $`2\kappa `$. Here we shall extend the study to a cavity damped by a thermal reservoir at finite temperature, so that the mean number of thermal photons, $`N`$, in the cavity mode is nonzero. We show that, even in the case of $`\delta `$ and $`\omega _{21}`$ being the same order of the cavity linewidth $`2\kappa `$, the cavity-induced interference is still significant when $`N1`$, and that interference-assisted gain may occur in one component of the Autler-Townes doublet for certain cavity resonant frequency. Such interference-related gain in the Autler-Townes doublet is also reported in free space . Our model consists of a V-type atom with the ground state $`\|0`$ coupled by the single-mode cavity field to the excited doublet $`\|1,\|2`$. Direct transitions between the excited sublevels $`\|1`$ and $`\|2`$ are dipole forbidden. The master equation for the total density matrix operator $`\rho _T`$ in the frame rotating with the average atomic transition frequency $`\omega _0=(\omega _{10}+\omega _{20})/2`$ takes the form $$\dot{\rho }_T=i[H_A+H_C+H_I,\rho _T]+\rho _T,$$ (1) where $`H_C`$ $`=`$ $`\delta a^{}a,`$ (3) $`H_A`$ $`=`$ $`{\displaystyle \frac{1}{2}}\omega _{21}\left(A_{22}A_{11}\right),`$ (4) $`H_I`$ $`=`$ $`i\left(g_1A_{01}+g_2A_{02}\right)a^{}h.c.,`$ (5) $`\rho _T`$ $`=`$ $`\kappa (N+1)\left(2a\rho _Ta^{}a^{}a\rho _T\rho _Ta^{}a\right)`$ (7) $`+\kappa N\left(2a^{}\rho _Taaa^{}\rho _T\rho _Taa^{}\right),`$ with $$\delta =\omega _C\omega _0,\omega _{21}=E_2E_1,g_i=𝐞_\lambda 𝐝_{0i}\sqrt{\frac{\mathrm{}\omega _C}{2ϵ_0V}},(i=1,\mathrm{\hspace{0.17em}2}).$$ (8) Here $`H_C`$, $`H_A`$ and $`H_I`$ are the unperturbed cavity, the unperturbed atom and the cavity-atom interaction Hamiltonians respectively, while $`\rho _T`$ describes damping of the cavity field by the continuum electromagnetic modes at finite temperature, characterized by the decay constant $`\kappa `$ and the mean number of thermal photons $`N`$; $`a`$ and $`a^{}`$ are the photon annihilation and creation operators of the cavity mode, and $`A_{ij}=\|ij\|`$ is the atomic population (the dipole transition) operator for $`i=j`$ $`(ij)`$; $`\delta `$ is the cavity detuning from the average atomic transition frequency, $`\omega _{21}`$ is the splitting of the excited doublet of the atom, and $`g_i`$ is the atom-cavity coupling constant, expressed in terms of $`𝐝_{ij},`$ the dipole moment of the atomic transition from $`\|j`$ to $`\|i,`$ $`𝐞_\lambda `$, the polarization of the cavity mode, and $`V,`$ the volume of the system. In the remainder of this work we assume that the polarization of the cavity field is pre-selected, i.e., the polarization index $`\lambda `$ is fixed to one of two possible directions. In this paper we are interested in the bad cavity limit: $`\kappa g_i`$, that is the atom-cavity coupling is weak, and the cavity has a low $`Q`$ so that the cavity field decay dominates. The cavity field response to the continuum modes is much faster than that produced by its interaction with the atom, so that the atom always experiences the cavity mode in the state induced by the thermal reservoir. Thus one can adiabatically eliminate the cavity-mode variables, giving rise to a master equation for the atomic variables only , which takes the form, $`\dot{\rho }`$ $`=`$ $`i[H_A,\rho ]`$ (14) $`+\{F(\omega _{21})(N+1)[\|g_1\|^2(A_{01}\rho A_{10}A_{11}\rho )+g_1g_2^{}(A_{01}\rho A_{20}A_{21}\rho )]`$ $`+F(\omega _{21})(N+1)\left[\|g_2\|^2\left(A_{02}\rho A_{20}A_{22}\rho \right)+g_1^{}g_2\left(A_{02}\rho A_{10}A_{12}\rho \right)\right]`$ $`+F(\omega _{21})N\left[\|g_1\|^2\left(A_{10}\rho A_{01}\rho A_{00}\right)+g_1g_2^{}A_{20}\rho A_{01}\right]`$ $`+F(\omega _{21})N\left[\|g_2\|^2\left(A_{20}\rho A_{02}\rho A_{00}\right)+g_1^{}g_2A_{10}\rho A_{02}\right]`$ $`+h.c.\}`$ where $`F(\pm \omega _{21})=\left[\kappa +i(\delta \pm \omega _{21}/2)\right]^1`$. Obviously, the equation (14) describes the cavity-induced atomic decay into the cavity mode. The real part of $`F(\pm \omega _{21})\|g_j\|^2`$ represents the cavity-induced decay rate of the atomic excited level $`j(=1,\mathrm{\hspace{0.17em}2})`$, while the imaginary part is associated with the frequency shift of the atomic level resulting from the interaction with the vacuum field in the detuned cavity. The other terms, $`F(\pm \omega _{21})g_ig_j^{},(ij)`$, however, represent the cavity-induced correlated transitions of the atom, i.e., an emission followed by an absorption of the same photon on a different transition, ($`\|1\|0\|2`$ or $`\|2\|0\|1`$), which give rise to the effect of quantum interference. The effect of quantum interference is very sensitive to the orientations of the atomic dipoles and the polarization of the cavity mode. For instance, if the cavity-field polarization is not pre-selected, as in free space, one must replace $`g_ig_j^{}`$ by the sum over the two possible polarization directions, giving $`\mathrm{\Sigma }_\lambda g_ig_j^{}𝐝_{0i}𝐝_{0j}^{}`$ . Therefore, only non-orthogonal dipole transitions lead to nonzero contributions, and the maximal interference effect occurs with the two dipoles parallel. As pointed out in Refs. however, it is questionable whether there is a isolated atomic system with parallel dipoles. Otherwise, if the polarization of the cavity mode is fixed, say $`𝐞_\lambda =𝐞_x`$, the polarization direction along the $`x`$-quantization axis, then $`g_ig_j^{}\left(𝐝_{0i}\right)_x\left(𝐝_{0j}^{}\right)_x`$, which is nonvanishing, regardless of the orientation of the atomic dipole matrix elements. It is apparent that if $`\kappa \delta ,\omega _{21}`$, the frequency shifts are negligibly small , and this equation (14) reduces to that of a V-atom with two parallel transition matrix elements in free space . In the following we shall discuss the effect of quantum interference in the situation of $`\omega _{21}\kappa `$ and $`N1`$, by examining the steady-state absorption spectrum of such a system, which is defined as $$A(\omega )=\mathrm{}e_0^{\mathrm{}}\underset{t\mathrm{}}{lim}[P(t+\tau ),P^{}(t)]e^{i\omega \tau }d\tau ,$$ (15) where $`\omega =\omega _p\omega _0`$, and $`\omega _p`$ is the frequency of the probe field and $`P(t)=d_1A_{01}+d_2A_{02}`$ is the component of the atomic polarization operator in the direction of the probe field polarization vector $`𝐞_p`$, with $`d_i=𝐞_p𝐝_{0i}`$. With the help of the quantum regression theorem, one can calculate the spectrum from the Bloch equations, $`\dot{A}_{11}`$ $`=`$ $`\left[F(\omega _{21})+F^{}(\omega _{21})\right]\|g_1\|^2\left[(N+1)A_{11}NA_{00}\right]`$ (17) $`F(\omega _{21})g_1^{}g_2(N+1)A_{12}F^{}(\omega _{21})g_1g_2^{}(N+1)A_{21},`$ $`\dot{A}_{22}`$ $`=`$ $`\left[F(\omega _{21})+F^{}(\omega _{21})\right]\|g_2\|^2\left[(N+1)A_{22}NA_{00}\right]`$ (19) $`F^{}(\omega _{21})g_1^{}g_2(N+1)A_{12}F(\omega _{21})g_1g_2^{}(N+1)A_{21},`$ $`\dot{A}_{12}`$ $`=`$ $`F(\omega _{21})g_1g_2^{}(N+1)A_{11}F^{}(\omega _{21})g_1g_2^{}(N+1)A_{22}+\left[F(\omega _{21})+F^{}(\omega _{21})\right]g_1g_2^{}NA_{00}`$ (21) $`\left[F^{}(\omega _{21})\|g_1\|^2(N+1)+F(\omega _{21})\|g_2\|^2(N+1)+i\omega _{21}\right]A_{12},`$ $`\dot{A}_{01}`$ $`=`$ $`\left[F(\omega _{21})\|g_1\|^2(2N+1)+F(\omega _{21})\|g_2\|^2Ni{\displaystyle \frac{\omega _{21}}{2}}\right]A_{01}F(\omega _{21})g_1^{}g_2(N+1)A_{02},`$ (22) $`\dot{A}_{02}`$ $`=`$ $`\left[F(\omega _{21})\|g_1\|^2N+F(\omega _{21})\|g_2\|^2(2N+1)+i{\displaystyle \frac{\omega _{21}}{2}}\right]A_{02}F(\omega _{21})g_1g_2^{}(N+1)A_{01}.`$ (23) To monitor quantum interference, we insert a factor $`\eta (=0,\mathrm{\hspace{0.17em}1})`$ in the cross transition terms $`g_ig_j^{}`$. When $`\eta =0`$, the cross transitions are switched off, so no quantum interference is present. Otherwise, the effect of quantum interference is maximal. Figure 1 shows the Autler-Townes spectra for $`g_1=g_2=10`$, $`\kappa =\omega _{21}=100,N=10`$, and different cavity detunings. In the absence of the interference ($`\eta =0`$), two transition paths, $`\|0\|1`$ and $`\|0\|2`$, are independent, which lead to the lower and higher frequency sidebands of the absorption doublet, respectively. It is not difficult to see that the spectral heights and linewidths are mainly determined by the cavity-induced decay constants $`\gamma _i(i=1,\mathrm{\hspace{0.17em}2})`$ of the excited states, which have the forms $$\gamma _1=\frac{\kappa \|g_1\|^2}{\kappa ^2+(\delta +\omega _{21}/2)^2},\gamma _2=\frac{\kappa \|g_2\|^2}{\kappa ^2+(\delta \omega _{21}/2)^2},$$ (24) which vary with the cavity frequency. It is evident that $`\gamma _1<\gamma _2`$ when $`\delta >0`$, and both $`\gamma _1`$ and $`\gamma _2`$ decrease as $`\delta `$ increases. Noting that the lower and higher frequency peaks have respective liewidths $`\mathrm{\Gamma }_l=\gamma _1(2N+1)+\gamma _2N`$ and $`\mathrm{\Gamma }_h=\gamma _1N+\gamma _2(2N+1)`$, and are proportional to $`\mathrm{\Gamma }_{l,h}^1`$, the lower frequency sideband is slightly higher than the higher frequency one in the case of $`\delta >0`$ and both the sidebands can be narrowed by increasing the cavity detuning, see for example, the dashed lines in the following three figures. Whereas, the spectral features are dramatically modified in the presence of the cavity induced interference ($`\eta =1`$). When the cavity is resonant with the average frequency of the atomic transitions, $`\delta =0`$, the doublet is symmetric, and its sidebands are higher and wider than that for $`\eta =0`$, as shown in the frames 1(a), 2(a) and 3(a). Otherwise, it is asymmetric. Either sideband of the doublet can be suppressed, depending upon the cavity frequency, e.g., the higher frequency sideband is suppressed for $`\delta =10`$, $`50`$ and $`100`$, see in Figs. 1(b)–1(d), while the sideband is enhanced for $`\delta =200`$, shown in Fig. 1(e). When the cavity frequency is far off resonant with the atomic transition frequencies, say $`\delta =500`$ in Fig. 1(f), the absorption spectra for $`\eta =0`$ and $`1`$ are virtually same, that is the effect of the cavity induced interference is negligible small. Rather surprisingly, the frame 1(c) shows probe gain in the higher frequency sideband, without the help of any coherent pumping. Moreover, increasing the mean number of thermal photons $`N`$ may enhance the probe gain, see for instance, in Fig. 2 for $`N=20`$, in which the higher-frequency probe gain even occurs for a relative small cavity detuning, say $`\delta =10`$ in the frame 2(b). Contrastively, when the detuning is very large, the probe beam can be amplified at the lower-frequency sideband, rather than at the higher-frequency one, as shown in the frame 2(e) for $`\delta =200`$ for example. Fig. 2 also exhibits that the linewidths are broadened for large number of thermal photons. We present the Autler-Townes spectrum for a large excited level-splitting, $`\omega _{21}=200`$, and a large number of thermal photons, $`N=20`$, in Fig. 3, in which the more pronounced gain, comparing with that for $`\omega _{21}=100`$, is displayed at either the lower-frequency sideband for $`\delta =10`$, $`50`$ and $`100`$, or the higher-frequency sideband for $`\delta =200`$. One can also find that for the large level-splitting, the effect of the cavity-induced interference is still significant when $`\delta =500`$, as shown in Fig. 3(f), where the lower frequency peak is almost suppressed while the other is greatly enhanced. However, when $`\delta \omega _{21}`$, say $`\delta =1000`$ for instance, the effect of the interference disappears (we have exhibited no figure here). In what follows, we shall see that the probe gain is a direct consequence of the cavity-induced quantum interference between the two transition paths, $`\|0\|1`$ and $`\|0\|2`$. The gain at different sidebands has different origin. To show this, we first plot the steady-state population differences between the excited sublevels and the ground level, $`A_{11}A_{00}`$ and $`A_{22}A_{00}`$, and the coherence between the excited sublevels, $`A_{12}`$, against the cavity detuning $`\delta `$ in Fig. 4 for $`g_1=g_2=10`$, $`\kappa =100`$, $`\omega _{21}=200`$ and $`N=20`$. It is clearly that the steady-state populations and coherence are highly dependent on the cavity frequency. The coherence is symmetric with the cavity detuning and reaches the maximum value at $`\delta =0`$, while the population differences are asymmetric. Furthermore, the population inversion may be achieved for certain cavity frequency, for example, if $`143.8<\delta <650`$, then $`A_{11}A_{00}>0`$, while $`A_{22}>A_{00}`$ in the region of $`650<\delta <143.8`$. Therefore, the gain in the region of $`143.8<\delta <143.8`$ must stem from the cavity-induced steady-state coherence between the two dipole-forbidden excited sublevels, rather than from the population inversion between the two dipole transition levels. Whereas, the population inversions may result in the probe gain when the cavity detuning is in the regions of $`650<\delta <143.8`$ and $`143.8<\delta <650`$. We thus conclude that, in the case of $`\delta >0,`$ as shown in Figs. 1-3, the gain at the lower-frequency sideband comes from the contribution of the steady-state atomic coherence $`A_{12}`$, while the gain at the other sideband is attributed to the steady-state population inversion $`\left(A_{11}>A_{00}\right)`$. Noting that, in the absence of the interference ($`\eta =0`$), $`A_{11}=A_{22}=N/(3N+1)`$, $`A_{00}=(N+1)/(3N+1)`$, and $`A_{12}=0`$ are independent of the cavity detuning, the cavity frequency dependence of the steady-state populations and coherence manifests the cavity-induced quantum interference. To further explore the origin of the probe gain, we separate the Autler-Townes spectrum into two parts, in which one corresponds to the contribution of the populations, while the other results from the coherence, in Fig. 5 for $`g_1=g_2=10`$, $`\kappa =100`$, $`\omega _{21}=200`$, $`N=20`$, and various cavity frequencies. It is obvious that when $`\delta =0`$, $`50`$ and $`100`$, the contributions of the coherence to the spectrum are negative (i.e., probe gain), whereas the populations make positive contributions, see, for example, in frames 6(a)–6(c). One can also see that the spectral component resulting from the populations is symmetric only when $`\delta =0`$, otherwise, it has different values at the lower and higher frequency sidebands, which are proportional to $`\left(A_{00}A_{11}\right)`$ and $`\left(A_{00}A_{22}\right)`$, respectively. As shown in Fig. 4, if the cavity detuning is zero, then $`\left(A_{00}A_{11}\right)=\left(A_{00}A_{22}\right)`$, whereas, $`\left(A_{00}A_{11}\right)>\left(A_{00}A_{22}\right)`$ for $`\delta =50`$ and $`100`$. As a result, the lower frequency peak is higher than that of the other one in the cases of $`\delta =50`$ and $`100`$. Therefore, the total spectrum may exhibit the probe gain at the higher frequency sideband at these cavity frequencies, see, for example, in Figs. 3(c) and 3(d). The gain is purely attributed to the cavity-induced steady-state atomic coherence. However, when $`\delta =200`$, the situation is reverse: the coherence gives rise to the probe absorption, while the populations lead to the gain at the lower frequency sideband, due to the population inversion between the levels $`\|0`$ and $`\|1`$, as illustrated in Fig. 4. In summary, we have shown that maximal quantum interference can be achieved in a V-type atom coupled to a single-mode, frequency-tunable cavity field at finite temperature, with a pre-selected polarization in the bad cavity limit. There are no special restrictions on the atomic dipole moments, as long as the polarization of the cavity field is pre-selected. We have investigated the cavity modification of the Autler-Townes spectrum of such a system, and predicted the probe gain at either sideband of the doublet, depending upon the cavity resonant frequency, when the excited sublevel splitting is very large (the same order as the cavity linewidth) and the mean number of thermal photons $`N1`$. The gain occurring at different sidebands has the various origin: in the case of $`\delta >0`$, the lower frequency gain is due to the nonzero steady-state coherence, while the higher frequency one is attributed to the steady-state population inversion. Both the nonzero coherence and population inversion originate from the cavity-induced quantum interference. ###### Acknowledgements. This work is supported by the United Kingdom EPSRC. We gratefully acknowledge conversations with G. S. Agarwal and Z. Ficek. We would also like to thank S. Menon for bringing their paper to our attention.
warning/0004/math-ph0004009.html	ar5iv	text	# Tri-valent graphs and solitons (September 20, 1999) It is shown that a real self-adoint operator of order 4 on the tri-valent tree $`\mathrm{\Gamma }_3`$ has $`(L,A,B)`$-triple deformations that preserve one energy level. Laplace type discrete symmetries of such operators are constructed. Until recently, nonlinear integrable systems have been known for lattices $`Z`$ and $`Z^2`$, only. That are $`(L,A)`$-pairs for Toda type systems for $`Z`$ and $`(L,A,B)`$-triples in the case of $`Z^2`$). Discrete Eiler-Darboux and Laplace type spectral symmetries for linear second order operators $`L`$ also have been known only for these lattices (see ). Note, that tri-valent tree $`\mathrm{\Gamma }_3`$ is a discrete model of the hyperbolic geometry of Lobachevski plane, unlike $`Z^2`$ which is a discrete model of the Euclidian plane. All previous attempts to find isospectral deformation of the second order operators $`L`$ on $`\mathrm{\Gamma }_3`$, even in the form of the $`(L,A,B)`$-triple, $`\dot{L}=LABL`$, which preserve only one spectral level $`L\mathrm{\Psi }=0`$, failed (see ). Let $`(L\mathrm{\Psi })_P=_Qb_{P,Q}\mathrm{\Psi }_Q`$ be a linear operator on a graph. Then the maximum diameter $`max_Pd(Q_1,Q_2)`$, where $`b_{P,Q_1}0,b_{P,Q_2}0`$, or $`b_{Q_1,Q_2}0`$, is called the order of the equation $`L\mathrm{\Psi }=0`$. The metric on the graph is defined by the condition that the length of an edge equals 1. Here $`\mathrm{\Psi }_P`$ is a function on vertices $`P`$. Let us consider a graph such that each edge has two vertices and at each vertex three edges come together. ###### Theorem 1 Let $`L`$ be a generic real self-adjoint operator of order 4 on the tree $`\mathrm{\Gamma }_3`$. Then there exist isospectral deformations of one energy level $`L\mathrm{\Psi }=0`$ of the form $`(L,A,B)`$-triple: $$\dot{L}=LABL$$ where $$(L\mathrm{\Psi })_P=b_{PP^{\prime \prime }}\mathrm{\Psi }_{P^{\prime \prime }}+b_{PP^{}}\mathrm{\Psi }_P^{}+w_P\mathrm{\Psi }_P,$$ $`P,P^{},P^{\prime \prime }`$ are vertices, $`d(P,P^{\prime \prime })=2,d(P,P^{})=1`$, and we assume that $`b_{P,P^{\prime \prime }}>0.`$ Here $`B=A^t`$, $`(A\mathrm{\Psi })_P=c_{PP^{}}\mathrm{\Psi }_P^{}.`$ The coefficients $`c_{P,P^{}}`$ for the nearest neighbours $`P,P^{}`$ are defined as follows. Let us fix a vertex $`P_0`$ on $`\mathrm{\Gamma }_3`$, and let $`R_i\gamma `$ be edges of the shortest path $`\gamma `$ which connects $`P_0`$ and $`P`$ and oriented from $`P_0`$ towards $`P`$. Let edges $`R_{i_1}^{},R_{i_2}^{}`$ meet at the initial vertex of the edge $`R_i`$, and edges $`R_{i_1}^{\prime \prime },R_{i_2}^{\prime \prime }`$ get off the second vertex of $`R_i`$. Then the formula $$\chi (R_i)=\frac{\left(b_{R_{i_1}^{\prime \prime }R_i}b_{R_{i_2}^{\prime \prime }R_i}\right)}{\left(b_{R_{i_1}^{}R_i}b_{R_{i_2}^{}R_i}\right)}.$$ define multiplicative one-cocycle on $`\mathrm{\Gamma }_3`$. The coefficients of the operator $`A`$ are equal $$c_R=\frac{1}{b_{R_1^{}R_2^{}}}\left(\underset{R_i\gamma }{}\chi (R_i)\right),R=PP^{}.$$ These formulas imply that the operator $`LA+A^tL`$ has order not greater than $`4`$. Therefore, the dynamical systems $`\dot{L}=LA+A^tL`$ is well-defined. It has the form: $`\dot{b}_{PP^{\prime \prime }}`$ $`=`$ $`b_{P^{}P^{\prime \prime }}c_{P^{}P}+c_{P^{}P}b_{P^{}P};`$ $`\dot{b}_{PP^{}}`$ $`=`$ $`b_{P^{}P_i^{\prime \prime }}c_{P_i^{\prime \prime }P^{}}+c_{P_\alpha ^{}P}b_{P_\alpha ^{}P^{}}+w_Pc_{PP^{}}+w_P^{}c_{P^{}P};`$ $`\dot{w}_P`$ $`=`$ $`2b_{PP^{}}c_{P^{}P},i,\alpha =1,2.`$ (0.1) Here $`P_\alpha ^{}PP^{}P_i^{\prime \prime }`$ are shortest paths of length $`d=3`$, that contain the edge $`PP^{}=R`$. Remark 1. For any tri-valent graph $`\mathrm{\Gamma }`$ the coefficients $`c_{PP^{}}`$ of the operator $`A`$ are defined on the abelian cover of $`\mathrm{\Gamma }`$, defined by the the cocycle $`\chi `$. ###### Theorem 2 A generic self-adjoint real fourth order operator $`L`$ on the tree $`\mathrm{\Gamma }_3`$ admits one-parametric family of factorizations of the form $$L=Q^tQ+u_P,\mathrm{where}(Q\psi )_P=\underset{Q}{}d_{PQ}\psi _Q+v_P\psi _P,$$ where $$b_{PP^{\prime \prime }}=d_{P^{}P}d_{P^{}P^{\prime \prime }};b_{PP^{}}=d_{P^{}P}v_P^{}+d_{PP^{}}v_P,$$ $$w_P=v_P^2+\underset{P^{}}{}d_{P^{}P}^2+u_P(\mathrm{let}d_{PQ}>0.)$$ The coefficients $`d_{PQ}`$ are defined uniquely, the coefficient $`v_P`$ is defined by one parameter, that is its evaluation at the initial point $`P_0\mathrm{\Gamma }_3`$. The factorization define the Laplace type transform $$\stackrel{~}{L}=Qu_P^1Q^t+1,\stackrel{~}{\psi }=Q\psi ,$$ where $`\stackrel{~}{L}\stackrel{~}{\psi }=0`$, if $`L\psi =0.`$ The self-adjoint operator $`\stackrel{~}{L}`$ is defined uniqully up to the transformation $$\stackrel{~}{L}f_P^1\stackrel{~}{L}f_P,\stackrel{~}{\psi }f_P^1\stackrel{~}{\psi }.$$ Let us take$`f_P=u_P^{1/2}.`$ Then we get $`\stackrel{~}{L}=\stackrel{~}{Q}^t\stackrel{~}{Q}+u_P,`$ where $$\stackrel{~}{Q}=u_P^{1/2}Q^tu_P^{1/2},\stackrel{~}{\psi }=u_P^{1/2}Q\psi .$$ (cf. for $`Z^2`$). Remark 2. The factorization of $`L`$ depends on the existence of a solution of the linear equation $`b_{PQ}=d_{QP}v_Q+d_{PQ}v_P`$. Note, that this operator has non-trivial (one-dimensional) kernel iff the cocycle $`\chi `$ is cohomological to zero on $`\mathrm{\Gamma }`$.
warning/0004/astro-ph0004353.html	ar5iv	text	# Galactic Gamma-Ray Background Radiation from Supernova Remnants ## 1 Introduction Observations of the diffuse Galactic $`\gamma `$-ray emission give information about the Galactic Cosmic Rays (GCRs), the interstellar gas and diffuse photon fields, and about the interactions between them. The observational results obtained with the Energetic Gamma Ray Experiment Telescope (EGRET) on the Compton Gamma Ray Observatory can be described fairly well by a suitable model for the diffuse interstellar gas, Cosmic Ray (CR) and photon distributions (e.g. Hunter et al. 1997a; Hunter et al. 1997b). However, above 1 GeV the observed average diffuse $`\gamma `$-ray intensity in the inner Galaxy, $`300^{}<l<60^{}`$, $`\|b\|10^{}`$, exceeds the model prediction significantly. There are at least two possible explanations for this discrepancy (e.g. Hunter et al. 1997b; Weekes et al. 1997). The high-energy $`\gamma `$-ray excess may indicate that the GCR spectrum observed in the local neighborhood is not representative of the diffuse CR population in the Galactic disk; a harder average diffuse proton spectrum is required to explain the $`\gamma `$-ray excess if it is due to $`\pi ^0`$-decay. An unresolved distribution of CR sources is the other possibility. The physical picture which we consider in this paper corresponds to the second possibility. The idea that CRs, after leaving their sources, could in principle produce $`\gamma `$-rays in ambient dense clouds with a harder spectrum than those produced by the average GCRs, was proposed by Aharonian and Atoyan (1996). Another proposed possibility, also invoking transport effects, is the local hardening of the CR energy spectrum in the direction perpendicular to the Galactic disk above strong CR sources, especially above the inner region of the Galaxy, due to faster CR convection in a faster Galactic Wind (Völk 1999). This contributes a principally observable harder than average $`\gamma `$-ray component in such regions. In contrast, we consider here the accelerating particles inside their sources, where they are much more strongly scattered than in the ISM, neglecting the contributions invoked by Aharonian and Atoyan, and by Völk (see also below). We assume that SNRs are the dominant sources of the Galactic CRs. On this premise we find that CRs, accelerated and confined in SNRs, give an important contribution to the high-energy $`\gamma `$-ray emission from the Galactic disk. Since the CR energy spectrum inside SNRs is much harder than on average in the Galaxy — the average spectrum being softened by rigidity-dependent escape from the Galaxy in the diffusion region above the disk — the relative SNR contribution increases with energy and becomes in fact dominant at $`\gamma `$-ray energies $`ϵ_\gamma >\mathrm{\hspace{0.17em}100}`$ GeV. It may substantially increase the diffuse TeV $`\gamma `$-ray emission from the Galactic disk so as to constitute a significant and spatially variable observational background which must also be taken into account in the search for spatially extended Galactic CR sources in this energy region. A physically analogous problem is the contribution of CR electrons in SNRs to the radio synchrotron spectrum of normal galaxies, without an Active Galactic Nucleus. It has been recently discussed by Lisenfeld & Völk (1999). In this paper we shall investigate pion-decay $`\gamma `$-ray emission from CR nuclei as well as Inverse Compton (IC) radiation and Bremsstrahlung due to CR electrons. It will be shown that the average IC $`\gamma `$-ray background from SNRs is comparable in magnitude and spectral form to the pion-decay background at high energies, whereas the corresponding Bremsstrahlung component is negligible. ## 2 Gamma-ray luminosity of old SNRs The majority of the GCRs, at least up to kinetic energies $`ϵ10^{14}`$ eV, is presumably accelerated in SNRs. According to modern theory a significant part of the hydrodynamic Supernova (SN) explosion energy $`E_{SN}10^{51}`$ erg is converted into CRs already in the early Sedov phase of the evolution, due to diffusive shock acceleration (e.g. Berezhko et al. 1996; Berezhko & Völk 1997). Later on, the CR energy content and the high-energy $`\gamma `$-ray production slowly decrease with time. This is at least true as long as the progenitor star is not so massive as to have a strong wind which significantly modifies the circumstellar medium (Berezhko & Völk, 2000). The total number of SNRs $`N_{SN}=\nu _{SN}T_{SN}`$ is an increasing function of their assumed life time $`T_{SN}`$, i.e. the time until which they can confine the accelerated particles; here $`\nu _{SN}`$ is the Galactic SN rate. Therefore we conclude that the population of the oldest SNRs dominates the total $`\gamma `$-ray luminosity of the ensemble of Galactic SNRs. Thus we consider only old SNRs which nevertheless are still strong enough to confine most of the CRs produced during the prior evolutionary stages. We then have the situation that the CRs in the Galaxy are represented by two basically different populations. The first one consists of the ordinary GCRs and presumably occupies a large Galactic residence volume quasi-uniformly. This residence volume exceeds that occupied by the CR sources by far (e.g. Ptuskin et al. 1997; for an earlier review, see Berezinsky et al. 1990). The second CR population, which we call Source Cosmic Rays (SCRs), is represented by shock accelerated CRs that are still confined in the localized SNRs. During the initial, active period of SNR evolution of about $`t<\mathrm{\hspace{0.17em}10}^5`$ yr when the SN shock is relatively strong, the volume occupied by the accelerated CRs practically coincides with the shock volume. In later stages the shock becomes weak and CRs begin to leave the SNR acceleration region. After some period of time $`T_{SN}`$ the escaping SCRs become very well mixed with the “sea” of GCRs. We shall assume that the transitory period during which SCRs are transformed into GCRs is much shorter than the preceding confinement time $`T_{SN}`$. The most important factor for CR confinement is the shock strength. Even in the phase where radiative cooling would formally become important, this remains true since CRs prevent cooling compression due to their pressure. Thus the assumed mean life time $`T_{SN}<\mathrm{\hspace{0.17em}10}^5`$ yr is determined by the shock dynamics more than by anything else. Since the $`\gamma `$-ray production due to GCRs is quite well studied (e.g. Hunter et al. 1997b; Mori 1997), it is primarily important to find the relative contribution of the SCR population. ### 2.1 Gamma rays from $`\pi ^0`$-decay The production rate of $`\pi ^0`$-decay $`\gamma `$-rays from inelastic CR - gas collisions, primarily p-p collisions, may be written in the form (Drury et al. 1994) $$Q_\gamma (ϵ)=Z_\gamma \sigma _{pp}cN_gn(ϵ),$$ (1) where $`N_g`$ is the local gas number density, $`\sigma _{pp}`$ is the inelastic p-p cross-section, $`Z_\gamma `$ is the so-called spectrum-weighted moment of the inelastic cross-section, $`n(ϵ)dϵ`$ is the CR spatial number density of CRs in the kinetic energy interval $`dϵ`$, and $`c`$ is the speed of light. Thus we have to primarily calculate $`n(ϵ)`$ for the two CR populations. The quasi-uniform GCR population in the gas disk is assumed to have roughly a power law spectrum in the relativistic range $$n_{GCR}(ϵ)=\frac{n_0^{GCR}(\gamma _{GCR}1)}{mc^2}\left(\frac{ϵ}{mc^2}\right)^{\gamma _{GCR}}.$$ (2) The total number $`n_0^{GCR}`$ of relativistic GCRs with $`ϵ>mc^2`$, per unit volume, can be expressed in terms of the CR energy density $`e_{GCR}`$: $$n_0^{GCR}=\frac{(\gamma _{GCR}2)e_{GCR}}{(\gamma _{GCR}1)mc^2},$$ (3) where $`m`$ is the proton mass. For simplicity we restrict our consideration here to the proton component which is energetically dominant in both the GCR and the SCR populations. In contrast to the GCR population, the SCRs are confined inside a discrete number $`N_{SN}`$ of SNRs. These are assumed to be predominantly located in the Galactic gas disk, of volume $`V_g`$. Spatially averaged over the Galactic disk volume their $`\gamma `$\- ray production rate is determined by an expression analogous to eq. (1), where instead of $`n(ϵ)`$ one should substitute the SCR distribution $$n_{SCR}(ϵ)=N_{SCR}(ϵ)N_{SN}/V_g,$$ (4) with $`N_{SCR}(ϵ)dϵ`$ being the overall (i.e. integrated over the SNR volume) SCR number in the energy interval $`dϵ`$, and should use an appropriate mean local interstellar medium (ISM) number density $`N_g^{SCR}`$ into which the SNe explode. Since the CRs produced inside SNRs have also a power-law spectrum $`N_{SCR}ϵ^{\gamma _{SCR}}`$ in the relativistic range, $`n_{SCR}(ϵ)`$ can be expressed in the same forms (2) and(3), with $$e_{SCR}=N_{SN}\delta E_{SN}/V_g,$$ (5) and putting $`\gamma _{SCR}>2`$. Here $`\delta `$ is the fraction of the SN explosion energy $`E_{SN}`$ converted into SCRs. However, according to the prediction from nonlinear kinetic theory (Berezhko et al. 1996), diffusive shock acceleration produces an extremely hard spectrum of SCRs at the early Sedov phase which is characterized by a power law index $`\gamma _{SCR}=2`$. In this case we have $$n_0^{SCR}=\frac{e_{SCR}}{mc^2\mathrm{ln}(ϵ_{max}/mc^2)},$$ (6) where $`ϵ_{max}`$ is the maximum SCR energy. For the ratio $`R=Q_\gamma ^{SCR}/Q_\gamma ^{GCR}`$ of the $`\gamma `$-ray production rates due to SCRs and GCRs, we have $`R(ϵ_\gamma )`$ $`=`$ $`{\displaystyle \frac{Z_\gamma ^{SCR}N_{SN}\delta E_{SN}}{Z_\gamma ^{GCR}(\gamma _{GCR}2)\mathrm{ln}(ϵ_{max}/mc^2)V_ge_{GCR}}}`$ (7) $`\times `$ $`\zeta \left({\displaystyle \frac{ϵ_\gamma }{mc^2}}\right)^{\gamma _{GCR}2},`$ where $`\zeta `$ is the ratio $`N_g^{SCR}/N_g^{GCR}`$, and $`N_g^{GCR}`$ denotes the average gas density in the disk. In fact we assume that the gas and the CRs are distributed uniformly inside each SNR which is approximately true for the old SNRs which we consider here. The parameter $`\zeta `$ describes a possible spatial correlation between SN occurrence and local ISM density. If on average Supernovae explode in a denser than average medium in the Galactic disk, then $`\zeta >1`$, whereas $`\zeta <1`$ in the opposite case. The main mass of the ISM $`M_g=4\times 10^9M_{}`$ is contained in a Galactic disk region of a thickness of about 240 pc, corresponding to the thickness of the HI gas (Dickey and Lockman 1990) which has the volume $`V_g=2.5\times 10^{66}`$ cm<sup>3</sup> referred to before. Here we take a disk radius of about $`10`$ kpc which implies an average gas density $`N_g^{GCR}=2\mathrm{cm}^3`$. Taking the relativistic part of the GCR spectrum to be characterized by $`e_{GCR}10^{12}`$ erg/cm<sup>3</sup>, and $`\gamma _{GCR}=2.75`$ which results in $`Z_\gamma ^{SCR}/Z_\gamma ^{GCR}=10`$ (Drury et al. 1994), we obtain for the standard set of SN parameters $`E_{SN}=10^{51}`$ erg, $`\nu _{SN}=1/30`$ yr<sup>-1</sup>: $$R(ϵ_\gamma )=0.16\zeta \left(\frac{T_{SN}}{10^5\text{yr}}\right)\left(\frac{ϵ_\gamma }{1\text{GeV}}\right)^{0.75},$$ (8) in addition using the rather moderate parameter values $`\delta =0.1`$ and $`ϵ_{max}=10^5mc^2`$ to characterize CR acceleration inside SNRs (e.g. Berezhko et al. 1996). One can see from this expression that for $`T_{SN}10^5`$ yr the $`\gamma `$ ray production due to SCRs becomes dominant already at energies $`ϵ_\gamma >\mathrm{\hspace{0.17em}10}`$ GeV. Note that the quantity $`\delta E_{SN}N_{SN}/(V_ge_{GCR})`$ represents the ratio of currently existing total SCR energy and GCR energy inside $`V_g`$. For the above set of parameters it is about 0.1. Despite the fact that the SCRs represent only a relatively small fraction of the total CR energy content even in the disk, they may dominate the $`\gamma `$-ray production at sufficiently high energies due to their much harder spectrum. It is clear that the quantity $`R(ϵ_\gamma )`$ determines the average ratio $`(dN_\gamma ^{SCR}/dϵ_\gamma )/(dN_\gamma ^{GCR}/dϵ_\gamma )`$ of $`\gamma `$-ray spectra produced in any region of the disk, by SCRs and GCRs, respectively. Therefore the total $`\gamma `$-ray spectrum measured from an arbitrary Galactic disk volume is expected to be $$\frac{dN_\gamma }{dϵ_\gamma }=\frac{dN_\gamma ^{GCR}}{dϵ_\gamma }[1.4+R(ϵ_\gamma )],$$ (9) where the additional factor 0.4 is introduced to approximately take into account the contribution of GCR electron component to the diffuse $`\gamma `$-ray emission at GeV energies (e.g. Hunter et al. 1997b). In Fig.1 we present the expected differential flux of $`\gamma `$-rays from the inner Galaxy, calculated for $`\zeta =1`$ and $`T_{SN}=10^5`$ yr, with the spectrum $`dN_\gamma ^{GCR}/dϵ_\gamma `$ taken from the paper by Hunter et al. (1997b), and extended into the region $`ϵ_\gamma >30`$ GeV according to the law $`ϵ_\gamma ^{2.75}`$. One can see that after inclusion of the $`\gamma `$-rays produced by SCRs, the calculated flux even exceeds the EGRET flux for $`ϵ_\gamma >\mathrm{\hspace{0.17em}20}`$ GeV. This suggests that expression (7) overestimates the $`\gamma `$-ray production by SCRs. It is possible that the overall source spectral index $`\gamma _{SCR}`$ is somewhat larger than 2 due to very late accumulation of only low-energy particles. It can also not be excluded that the confinement time $`T_{SN}`$ of the SCRs depends on energy. Due to their high mobility, the highest energy particles may leave the vicinity of parent SNR earlier and also more rapidly. This process of SCR escape into the ISM starts for the most energetic particles already at the early Sedov phase of SNR evolution (e.g. Berezhko et al. 1996; Berezhko & Völk 1997). Therefore at time $`T_{SN}`$, when the main part of SCRs are released from the SNR, their spectrum may be somewhat steeper than a $`\gamma _{SCR}=2`$ spectrum. Due to the importance of the problem we derive the relation between $`n_{SCR}`$ and $`n_{GCR}`$ in a different form. It leads to the same results if SNRs are the GCR source. We start from the usual leaky box balance equation $$\frac{n_{GCR}(ϵ)}{\tau _c}=\frac{N_{SCR}}{V_c}\nu _{SN},$$ (10) where $`V_c(ϵ)`$ is the energy-dependent residence volume occupied by GCRs that reach the gas disk during their mean residence time $`\tau _c`$ in $`V_c`$ . In the case of an extended Galactic halo due to a Galactic wind driven by the GCRs themselves, and for energies much larger than a few GeV, $`V_c`$ can be much greater than $`V_g`$. In fact $`V_c(ϵ)ϵ^{0.55}`$ in such a selfconsistent halo model (Ptuskin at al. 1997). Note that the leaky box model deals with a CR distribution $`n_{GCR}(ϵ)`$ averaged over the residence volume $`V_c`$. Therefore it can only be applied to the GCRs which can be assumed to be almost uniformly distributed in the residence volume. It is not valid for the SCRs, because their behavior is determined not only by large-scale transport but also by other physical factors which, for example, lead to their acceleration. Technically the volume $`V_{SCR}`$, occupied by the SCRs, should be excluded from the residence volume $`V_c`$ and the SCRs appear in the balance equation (10) for the GCRs only in the form of a source term $`N_{SCR}\nu _{SN}/V_c`$ as a CR population released from the source region into the ISM after some unspecified evolutionary period $`T_{SN}`$. Therefore eq.(10) does not depend upon $`T_{SN}`$. However, the effects produced by the SCRs confined inside the ensemble of simultaneously existing SNRs, for example the additional $`\gamma `$-ray production, essentially depends on $`T_{SN}`$, since it directly determines the total number $`N_{SN}`$ of simultaneously existing SNRs. Using eq. (4) we can write $$\frac{n_{SCR}}{n_{GCR}}=\frac{V_cT_{SN}}{V_g\tau _c}=\frac{T_{SN}}{\tau _g}.$$ (11) The GCR residence time in the disk volume, $`\tau _g=\tau _cV_g/V_c`$, can be derived from the measured grammage $`x`$, which is the mean mass of Interstellar matter traversed by GCRs of speed $`v`$ in the course of their random walk in the Galaxy: $$\tau _g=\frac{xV_g}{vM_g}.$$ (12) The measured grammage at high energies $`ϵϵ_0=4.4`$ GeV is $`x=14(v/c)(ϵ/4.4\text{GeV})^{0.60}`$ g/cm<sup>2</sup>; for $`ϵ<ϵ_0`$, $`x=14v/c`$ g/cm<sup>2</sup> (Engelman et al. 1990). Therefore the residence time in the gas disk can be written in the form $$\tau _g=\tau _0(ϵ/ϵ_0)^\mu ,$$ (13) where $`\tau _0=4.6\times 10^6`$ yr, $`ϵ_0=4.4`$ GeV, $`\mu =0.6`$, and $`ϵϵ_0`$. This experimentally inferred value for $`\mu `$ closely agrees with the theoretical result of Ptuskin et al. (1997). We note that at relativistic energies $`ϵ>mc^2`$, according to the initial balance eq.(10), the GCR spectrum $`n_{GCR}ϵ^{\gamma _{GCR}}`$ and the overall SCR spectrum $`N_{SCR}ϵ^{\gamma _{SCR}}`$ should be connected by the relation $$\gamma _{SCR}=\gamma _{GCR}\mu .$$ (14) For $`\mu =0.60`$ the source should produce a SCR spectrum with $`\gamma _{SCR}=2.15`$, while for the case $`\gamma _{SCR}=2`$ one would need $`\mu =0.75`$. At the same time the SCR distribution $`n_{SCR}(ϵ)ϵ^{\gamma _{SCR}^{}}`$, averaged over the gas disk, can have a different shape compared to $`N_{SCR}(ϵ)`$ in the case of an energy dependent SNR confinement time $`T_{SN}(ϵ)`$, according to eq. (11). It can be steeper, $`\gamma _{SCR}^{}=\gamma _{GCR}\mu +\beta `$, if $`T_{SN}ϵ^\beta `$ is a decreasing function of energy. It would mean that the highest energy SCRs leave the parent SNR faster than the lower energy SCRs. Eq.(11) leads to a simple expression for the $`\gamma `$-ray production ratio $$R(ϵ_\gamma )=\zeta \frac{Z_\gamma ^{SCR}T_{SN}}{Z_\gamma ^{GCR}\tau _g},$$ (15) independently of $`\nu _{SN}`$. Substituting the residence time in the form (13), and taking $`\gamma _{SCR}=2.15`$, which leads to $`Z_\gamma ^{SCR}/Z_\gamma ^{GCR}=7.5`$ (Drury et al. 1994) in eq.(15), we obtain for $`ϵ_\gamma 4.4`$ GeV: $$R(ϵ_\gamma )=0.07\zeta \left(\frac{T_{SN}}{10^5\text{yr}}\right)\left(\frac{ϵ_\gamma }{1\text{GeV}}\right)^{0.6}.$$ (16) We shall only consider $`\gamma `$-ray energies $`ϵ_\gamma 4.4`$ GeV. One can see here again that the SCR contribution is determined by the value of the confinement time inside SNRs, $`T_{SN}`$. Unfortunately there is no detailed description of when and how CRs, accelerated in SNR, are released into the ISM. Nevertheless one can give some constraints on this process. According to the standard theory, the expanding SNR shock produces a power law CR spectrum up to the maximum energy (Berezhko 1996; Berezhko et al. 1996; Berezhko & Völk 1997) $$ϵ_mR_sV_s,$$ (17) which is determined by the radius $`R_s`$ and speed $`V_s`$ of the shock. The CRs with the highest energy $`ϵ_{max}`$ are produced at the very beginning of the Sedov phase $`tt_0`$ when the product $`R_sV_s`$ has its maximum $`R_0V_0`$, where $$t_0=\frac{R_0}{V_0},R_0=\left(\frac{3M_{ej}}{4\pi \rho _g}\right)^{\frac{1}{3}},V_0=\sqrt{\frac{2E_{sn}}{M_{ej}}}$$ (18) are the sweep-up time, sweep-up radius and initial mean ejecta speed respectively; $`M_{ej}`$ denotes the ejecta mass, and $`\rho _g=N_gm`$ the ISM density. Subsequently, the product $`R_sV_s`$ decreases with time as $`t^{1/5}`$ and the SNR shock produces CRs with progressively lower cutoff energy $`ϵ_m(t)<ϵ_{max}=ϵ_m(t_0)`$. During that phase those CRs that were previously produced with energies $`ϵ_m<ϵ<ϵ_{max}`$ now propagate outward diffusively without significant influence of the SNR shock. If their expansion is still governed by the Bohm diffusion coefficient as during their acceleration, the expansion rate is only slightly higher than the SNR expansion rate and these particles should be considered as confined inside the source (i.e. the SNR). The opposite extreme case corresponds to the assumption that particles with energies $`ϵ>ϵ_m`$ do no more produce a high level of turbulence. Let us consider this pessimistic scenario in terms of confinement here. In this case the propagation of these very high energy particles is governed by the mean Galactic diffusion coefficient which is very much larger than the Bohm diffusion coefficient. In this situation particles with energy $`ϵ`$ should be considered as released from the source at the moment $`tt_0`$ when $`ϵ_m(t)`$ drops below $`ϵ`$. Since the particles with maximum energy are produced at $`tt_0`$, one can write $$T_{SN}(ϵ)=t_0\left(\frac{ϵ}{ϵ_{max}}\right)^5.$$ (19) Due to this strong dependence it is clear that $`T_{SN}(ϵ)`$ will still deviate from the overall gas dynamic life time $`T_{SN}^{tot}`$ only for large energies $`ϵ`$ near $`ϵ_{max}`$. All particles with $$ϵ/ϵ_{max}<(T_{SN}^{tot}/t_0)^{1/5}$$ will remain confined until $`T_{SN}^{tot}`$. Since the majority of Galactic SNe are core collapse SNe from stars with masses exceeding about $`8M_{}`$, we shall use $`M_{ej}=10M_{}`$ for purposes of estimate. Except for progenitor masses exceeding $`15M_{}`$, the progenitors have only a weak stellar wind. Therefore the assumption of a uniform circumstellar medium remains reasonable for the average properties of the SNR population. For the main fraction of CRs the confinement in SNRs terminates when the SNR shock becomes weak and produces CRs with a very steep spectrum that cannot anymore excite a high level of turbulence near the shock front. If we take $`M=4`$ as a critical Mach number, the corresponding SNR age will be $`t=(M_0/M)^{5/3}=1.5\times 10^3t_0`$, where $`M_0=V_0/c_{S0}`$ is the initial shock Mach number and $`c_{S0}`$ is the ISM sound speed. For an ISM with number density $`N_g=1`$ cm<sup>-3</sup>, $`c_{S0}4`$km/sec, and then $`t_01.5\times 10^3`$ yr which gives $`T_{SN}2\times 10^6`$ yr for $`M_{ej}=10M_{}`$ and $`E_{SN}=10^{51}`$ erg. This estimate shows that the SNR shock remains rather strong during a very long period of time. Another physical factor which can restrict the SCR confinement is the radiative cooling of the postshock gas. Approximately it becomes important when the postshock temperature drops below $`10^6`$ K, or when the postshock sound speed drops below $`c_{S2}100`$ km/s. In the case of a strong shock, with Mach number $`M1`$, the postshock sound speed is determined by the shock speed $`c_{S2}\sqrt{5}V_s/3`$. During the Sedov phase the shock speed decreases with time according to the law $`V_s=0.4\times V_0(t/t_0)^{3/5}`$. Therefore the shock speed drops to the value $`V_c=100`$ km/s at the age $`t_c=t_0(0.4V_0/V_c)^{5/3}`$. The above set of SN and ISM parameters gives $`t_c6\times 10^4`$ yr. Since gas clumping as a result of cooling may lead to effective SCR leakage from the SNR, a value of the confinement time $`T_{SN}=10^5`$ yr is reasonable for $`N_g=1\mathrm{cm}^3`$. (Fig. 1). In Fig.1 we present a calculated $`\gamma `$-ray spectrum based on the above expression (16) with $$T_{SN}=min\{10^5,10^3(ϵ/ϵ_{max})^5\}\text{yr},$$ (20) $`ϵ_{max}=10^5`$ GeV, and $`ϵ_\gamma =0.1ϵ`$, which is roughly valid for the hadronic considered $`\gamma `$-ray production process considered. At GeV energies in this case SCRs contribute about 10% of the total $`\gamma `$-ray flux. Due to their hard spectrum this contribution progressively increases with energy and becomes dominant at $`ϵ_\gamma >\mathrm{\hspace{0.17em}100}`$ GeV. It increases the expected TeV $`\gamma `$-ray flux by about a factor of ten. Note that the actual SCR contribution from the inner part of the Galaxy is somewhat higher than the above estimate due to the larger SNR concentration in this region. As one can see from Fig. 1, the SCR contribution in the case $`\gamma _{SCR}=2`$ , using eqs. (8) and (9), is for all energies larger than that for $`\gamma _{SCR}=2.15`$, using eq. (16). This difference is related to the different SCR acceleration efficiencies. In the first case it is characterized by the parameter $`\delta =0.1`$ which, according to eq. (5), directly determines the $`\gamma `$-ray production rate. In the second case the SCR acceleration efficiency $`\delta =[_{mc^2}^{ϵ_{max}}ϵN_{SCR}(ϵ)𝑑ϵ]/E_{SN}`$ is not contained in the final expression (15), but one can derive it easily from the balance equation (10) which gives: $$\delta \nu _{SN}E_{SN}=_{mc^2}^{ϵ_{max}}\frac{n_{GCR}V_g}{\tau _g}ϵ𝑑ϵ.$$ (21) To determine $`\delta `$ we use the (demodulated) GCR distribution $`n_{GCR}(ϵ)`$ $`=`$ $`8.1\times 10^{10}`$ (22) $`\times `$ $`\left({\displaystyle \frac{ϵ}{1\mathrm{GeV}}}+{\displaystyle \frac{mc^2}{1\mathrm{GeV}}}\right)^{2.75}\mathrm{cm}^3(\mathrm{GeV})^1.`$ (Ryan et al. 1972; Perko 1987). Substituting the values of $`\tau _g(ϵ)`$, $`V_g`$, $`E_{SN}`$, and $`\nu _{SN}`$, we obtain $`\delta 0.05`$ for the relativistic part of the GCR spectrum. The required acceleration efficiency is two times lower compared with the case $`\gamma _{SCR}=2`$ due to the essentially steeper spectrum. Note, that in both cases the SCR spectra with $`\gamma _{SCR}=2`$, $`\delta =0.1`$ and $`\gamma _{SCR}=2.15`$, $`\delta =0.05`$ contain about the same number of relativistic CRs. In the first case the required GCR residence time is $`\tau _gϵ^{0.75}`$, whereas the second case with $`\tau _gϵ^{0.6}`$ is close to the experiment (Engelman et al. 1990). Therefore we believe that the dashed line in Fig.1 represents the most reliable estimate for the expected diffuse $`\pi ^0`$-decay $`\gamma `$-ray emission, especially at high energies $`ϵ_\gamma >\mathrm{\hspace{0.17em}100}`$ GeV. Note also that the acceleration efficiency required by eq. (21), using an assumed Galactic SN rate $`\nu _{SN}=1/30`$ yr<sup>-1</sup> and a mean SN explosion energy $`E_{SN}=10^{51}`$ erg, is considerably lower than that predicted by shock acceleration theory, which gives $`\delta =0.2÷0.5`$ (Berezhko et al. 1996). Yet, in contrast to the acceleration models which determine $`\delta `$ from the injection rate and the nonlinear acceleration theory selfconsistently, eq. (21) determines only the product $`\delta E_{SN}\nu _{SN}`$ from observed quantities. The observationally inferred SN explosion energies $`E_{SN}`$ can be at least by a factor 2 smaller than $`10^{51}`$ erg, and from comparisons with galaxies similar to our own $`\nu _{SN}`$ could vary between $`1/30÷1/100`$ yr<sup>-1</sup>. Therefore the empirical value of $`\delta `$ from eq. (21) can vary between 0.05 and 1/3. Nevertheless, the theoretically determined efficiencies appear systematically too high. As a possible solution for this discrepancy one might assume that, just before being released, the SCRs lose an important part of their energy by adiabatic expansion so that the SCRs’ energy content inside a SNR is higher than the energy contained in the released spectrum $`N_{SCR}(ϵ)`$. However, there is little dynamical basis for such an assumption. Much more likely is that the very efficient CR acceleration inside SNRs predicted by the nonlinear kinetic theory, assuming spherical symmetry, takes place in reality only at some fraction of the SN shock surface, because suprathermal positive ion injection into the acceleration process on the highly oblique part of the shock can be significantly suppressed (Bennet & Ellison 1995; Malkov & Völk 1995). In this case the acceleration efficiency, calculated for a spherical SNR shock, should be reduced by a factor of a few. The actual SNR distribution can in fact be rather nonuniform within the disk volume, contrary to what we have assumed implicitly up to now. In this case the estimated value of $`R(ϵ_\gamma )`$, which describes the relative SCR contribution to the diffuse $`\gamma `$-ray emission, should be corrected by a factor $`N_{SN}^a/N_{SN}`$, where $`N_{SN}^a`$ is the expected number of SNRs in the observed region and $`N_{SN}`$ represents this number in the case of uniformly distributed SNRs. It is clear that the expected value of $`R(ϵ_\gamma )`$ is almost independent of the actual SNR distribution if the observed region is an essential part of the whole disk volume $`V_g`$. ### 2.2 Inverse Compton and Bremsstrahlung gamma-rays from SCR electrons Electrons, once being injected into the diffusive shock acceleration process, will be as efficiently accelerated in SNRs as are the protons. Even though there exist theoretical concepts (e.g. Levinson 1994; Galeev et al. 1995; McClements et al. 1997; Bykov & Uvarov 1999), electron injection is not completely understood. However, there is no doubt that electrons undergo continuous acceleration during SNR evolution. The spectral shape of accelerated electrons $`N_{SCR}^e(ϵ)`$ inside SNRs deduced from radio-observations on average agrees with what is expected from shock acceleration. Since relativistic electrons with energies $`ϵ>1`$ GeV are dynamically indistinguishable from protons, their source spectrum $`N_{SCR}^e(ϵ)=K_{ep}N_{SCR}(ϵ)`$ can differ from that of the protons $`N_{SCR}(ϵ)`$ only by some energy independent factor $`K_{ep}`$ that is determined by the injection process. High energy electrons produce $`\gamma `$-ray emission due to IC scattering, especially on the Cosmic Microwave Background (CMB) and by Bremsstrahlung on the interstellar gas. We shall first consider the IC contribution here. #### 2.2.1 Inverse Compton contribution In an approximate form, valid if the generating electron energy distribution is smoothly varying, like in the case of a power law considered here, the IC $`\gamma `$-ray emissivity $`Q_\gamma ^{IC}(ϵ_\gamma )`$ can be written as (e.g. Longair 1981, Berezinsky et al. 1990) $$Q_\gamma ^{IC}(ϵ_\gamma )=\sigma _TcN_{ph}n_{SCR}^e(ϵ_e)\frac{dϵ_e}{dϵ_\gamma },$$ (23) where $$ϵ_e=m_ec^2\sqrt{3ϵ_\gamma /(4ϵ_{ph})}$$ (24) is the energy of electrons which produce an IC photon with mean energy $`ϵ_\gamma `$, $`\sigma _T=6.65\times 10^{25}\mathrm{cm}^2`$ denotes the Thomson cross section, $`ϵ_{ph}=6.7\times 10^4`$ eV and $`N_{ph}=400`$ cm<sup>-3</sup> are the mean energy and number density of the CMB photons, respectively. Finally $`n_{SCR}^e=N_{SCR}^eN_{SN}^e/V_g`$ denotes the average spatial electron SCR number density. Thus the ratio of the IC to the $`\pi ^0`$-decay $`\gamma `$-ray production rate reads as $$\frac{Q_\gamma ^{IC}}{Q_\gamma ^{SCR}}=1028K_{ep}\left(\frac{1\mathrm{cm}^3}{N_g^{SCR}}\right)\left(\frac{ϵ_\gamma }{1\text{TeV}}\right)^{1/2}\frac{N_{SN}^e}{N_{SN}},$$ (25) where we have used $`\gamma _{SCR}=2`$ and $`\sigma _{pp}=4\times 10^{26}`$ cm<sup>-2</sup>; $`N_{SN}^e`$ denotes the number of SNRs which contribute to the IC emission at energy $`ϵ_\gamma `$ from CR electron sources. From this expression it appears as if the IC $`\gamma `$-ray contribution would be dominant at TeV-energies if $`K_{ep}`$ is as large as $`10^2`$ and if the mean gas number density inside SNRs is about $`N_g^{SCR}=1`$ cm<sup>-3</sup>. However, the radiative cooling time $`\tau _e(ϵ_e)`$ of electrons, which produce $`\gamma `$-rays with energy $`ϵ_\gamma =(ϵ_e/17.1\text{TeV})^2`$ TeV, $$\tau _e=7.3\times 10^3\left(\frac{10\mu \text{G}}{B}\right)^2\left(\frac{1\mathrm{TeV}}{ϵ_\gamma }\right)^{1/2}\text{yr},$$ (26) reaches the above assumed overall SNR confinement time $`T_{SN}=10^5`$ yr for $`\gamma `$-ray energies $`ϵ_\gamma <ϵ_\gamma ^{}=5.4(B/10\mu \mathrm{G})^4`$ GeV. Here $`B`$ is the magnetic field strength inside the source, whose typical value inside SNRs is about 10 $`\mu `$G. Therefore, taking into account the obvious relation $`N_{SN}^e/N_{SN}=\tau _e/T_{SN}`$, on average the relative IC contribution of the electron component of SCRs in TeV $`\gamma `$-ray can be written as $$\frac{Q_\gamma ^{IC}}{Q_\gamma ^{SCR}}=75.5K_{ep}\left(\frac{1\mathrm{cm}^3}{N_g^{SCR}}\right)\left(\frac{10^5\mathrm{yr}}{T_{SN}}\right)\left(\frac{10\mu \mathrm{G}}{B}\right)^2.$$ (27) It is independent of the $`\gamma `$-ray energy for $`ϵ_\gamma >ϵ_\gamma ^{}`$, and only given by eq. (25) with $`N_{SN}^e/N_{SN}=1`$ for $`ϵ_\gamma <ϵ_\gamma ^{}`$. This consideration shows that for the parameters assumed, and for $`K_{ep}`$ of the order of $`10^2`$, we have an IC contribution to the average $`\gamma `$-ray background which is comparable to the hadronic background for all energies above a few GeV. #### 2.2.2 Bremsstrahlung contribution At high energies, $`ϵ_e,ϵ_\gamma m_e\mathrm{c}^2`$, we have for the Bremsstrahlung $`\gamma `$-ray emissivity $`Q_\gamma ^{Br}`$ $$Q_\gamma ^{Br}(ϵ_\gamma )=2_{ϵ_{min}}^{\mathrm{}}𝑑ϵ_e\frac{d\sigma _{ep}^{Br}}{dϵ_\gamma }cN_g^{SCR}n_{SCR}^e(ϵ_e),$$ (28) where $`ϵ_{min}`$ is the minimum electron energy necessary to produce a Bremsstrahlung $`\gamma `$-ray of energy $`ϵ_\gamma `$, the factor 2 takes into account the contributions of electron-electron and electron-proton collisions, and where the differential electron-proton Bremsstrahlung cross-section is given by (e.g. Berezinsky et al. 1990) $$\frac{d\sigma ^{Br}(ϵ_e,ϵ_\gamma )}{dϵ_\gamma }=\frac{4\alpha r_0^2}{ϵ_\gamma }\left[\frac{4}{3}\frac{4}{3}\frac{ϵ_\gamma }{ϵ_e}+\left(\frac{ϵ_\gamma }{ϵ_e}\right)^2\right]\left[\mathrm{ln}\left(\frac{2ϵ_e}{mc^2}\frac{ϵ_eϵ_\gamma }{ϵ_\gamma }\right)\frac{1}{2}\right].$$ (29) Here $`\alpha 1/137`$ and $`r_0=2.818\times 10^{13}`$ cm denote the fine structure constant and the classical electron radius, respectively. For our chosen value $`\gamma _{SCR}=2`$, the integral for $`Q_\gamma ^{Br}(ϵ_\gamma )`$ can be calculated in closed form. In the limit $`\mathrm{ln}(ϵ_\gamma /m_ec^2)1`$, of interest here, it reduces to the asymptotic form $$Q_\gamma ^{Br}(ϵ_\gamma )=8\alpha r_0^2cN_g^{SCR}n_{SCR}^e(ϵ_\gamma )\mathrm{ln}\left(\frac{ϵ_\gamma }{m_ec^2}\right).$$ (30) Thus, finally, we obtain $$\frac{Q_\gamma ^{Br}}{Q_\gamma ^{SCR}}=\frac{8\alpha r_0^2K_{ep}}{Z_\gamma ^{SCR}\sigma _{pp}}\frac{N_{SN}^e}{N_{SN}}\mathrm{ln}\left(\frac{ϵ_\gamma }{m_ec^2}\right)=9.3K_{ep}\frac{N_{SN}^e}{N_{SN}}\left[1+0.066\mathrm{ln}\left(\frac{ϵ_\gamma }{1\text{TeV}}\right)\right]$$ (31) This small ratio implies that Bremsstrahlung $`\gamma `$-rays play no role for the average $`\gamma `$-ray background above GeV energies, if $`K_{ep}0.1`$, taking into account, that $`N_{SN}^e/N_{SN}`$ is always smaller than 1. ## 3 Discussion We note that of order ten SNRs of age younger than $`10^5`$ yr can on average lie within a 1 degree field of view of a detector directed towards the Galactic Center. Therefore a moderate fluctuation of the measured $`\gamma `$-ray intensity is expected due to variations of the actual number of SNRs within such a detector’s field of view. At the same time, for directions perpendicular to the Galactic plane, the chance to observe the contribution of SCR $`\gamma `$-ray emission is quite negligible. A question is then how one might best study the nonuniformities of this background experimentally. Clearly this is an investigation of its own. Therefore we would like to restrict ourselves to a few comments here. Due to the spectral form of the background its graininess is most pronounced at high $`\gamma `$-ray energies. This is even more true due to the fact that for individual sources with a very low magnetic field the IC emission could be much stronger than the $`\gamma `$-ray emission due to hadronic interactions; in addition the IC emission has a harder spectrum. This suggests the use of imaging atmospheric Cherenkov telescopes. Their resolution in angle and energy is as good or better than that of other ground-based detectors. However, the study of extended sources is not an easy task with imaging telescopes which have a very limited field of view, even employing the stereoscopic method. For low brightness extended sources a satellite instrument like the future GLAST detector is well suited since it does not have to deal with the charged CR background due to its use of an anticoincidence shield. On the other hand, the statistics achievable with a small area space detector gets very low above a few tens of GeV. Thus one should probably attempt such a study with both types of instruments due to their complementary properties. A different question concerns the limitations of our approach due to its concentration on SNRs as the sources of the GCRs. In fact the considerations in this paper can be applied to any alternative class of dominant GCR sources. The most important aspect, which leads to the dominance of SCRs in high-energy $`\gamma `$-ray production, is that the GCR sources should generate SCRs with a spectrum that is significantly harder than the GCR spectrum. Eq. (15) is valid for an arbitrary class of CR sources if we substitute some other value of the SCR confinement time $`T_S`$ instead of $`T_{SN}`$, since the grammage $`x`$ is an experimentally fixed quantity. Let us then assume that some class of compact CR sources produces an energy $`E_C`$ in the form of CRs with spectrum $`N_{SCR}ϵ^{\gamma _{SCR}}`$ with average frequency $`\nu _S`$, and let us further assume that this spectrum remains unchanged inside the source regions for some period of time $`T_S`$ after which these CRs are released into the ISM as the GCRs. It is obvious that due to the general energy requirement the production rate $`\nu _SE_C`$ should be about the same as $`\nu _{SN}\delta E_{SN}`$. The SCR energy $`E_C`$, deposited in some initial volume $`V`$, will produce a dynamically significant disturbance in the background ISM if we assume that the initial SCR energy density $`E_C/V`$ is much greater than the thermal ISM energy density which in turn is of order $`e_{GCR}`$. This will inevitably lead to the confinement of these SCRs inside an expanding, disturbed volume for some period of time $`T_S`$ before the SCRs will be released. It is difficult to give a general relation between $`E_C`$ and $`T_S`$ and there may exist only lower bounds on $`T_S`$, given $`E_C`$. Therefore, we cannot exclude speculative source classes with many weak but long-lived sources. The opposite case of many weak and short-lived sources is excluded to the extent that the present explanation of the hard $`\gamma `$-ray spectrum by the contribution of the SCRs is unique. ## 4 Summary Our considerations suggest that the SCRs can provide an essential contribution to the high-energy Galactic $`\gamma `$-ray flux. According to our estimates, depending on the parameters, the SCR contribution is less than 10% of the GCR contribution at GeV energies and it dominates at energies greater than 100 GeV due to its essentially harder spectrum. This conclusion is confirmed by calculations performed for the case when SNRs are the main source of GCRs. At TeV energies the SCRs increase the expected $`\gamma `$-ray flux from the Galactic disk by almost an order of magnitude. The single physical parameter which determines the SCR contribution due to hadronic interactions is the SCR confinement time $`T_{SN}`$. As far as the $`\gamma `$-ray emission due to $`\pi _0`$-decay is concerned, the above conclusions are valid for $`T_{SN}10^5`$ yr. Since this SCR contribution is proportional to $`T_{SN}`$, it would be negligible at TeV energies if $`T_{SN}<\mathrm{\hspace{0.17em}10}^4`$ yr. A SNR age of $`10^4`$ yr typically corresponds to the intermediate Sedov phase, when the SNR shock is still quite strong. Therefore it seems to be quite improbable that the GCRs are replenished from SNRs at such an early phase. For the IC contribution even a ten times shorter source life time would be sufficient at TeV energies. In fact, for the TeV IC emission the relevant time scale is the life time $`\tau _e(ϵ_e)`$ of parent SCR electrons due to their synchrotron losses, which is indeed about $`10^4`$ yr. For decreasing $`\gamma `$-ray energies $`\tau _e(ϵ_e)`$ increases beyond $`10^4`$ yrs, and therefore a source life time of this magnitude would become a limiting factor. Our conclusions remain valid for alternative classes of possible GCR sources with comparable overall energy release and comparable individual confinement times. We note that this contribution of the dominant GCR sources necessarily exists. As we argue, it may be sufficient by itself to explain the observed $`\gamma `$-ray excess, at least in the inner Galaxy where it is well documented, without a need to invoke additional particle populations (e.g. Pohl & Schlickeiser 1991). Acknowledgements This work has been supported in part by the Russian Foundation of Basic Research, grants 97-02-16132 and 00-02-17728, and through the Verbundforschung Astronomie/Astrophysik of the German BMBF, grant 05-2HD66A(7). EGB acknowledges the hospitality of the Max-Plank-Institut für Kernphysik where part of this work was carried out. HJV in turn acknowledges the hospitality of the Institute of Cosmophysical Research and Aeronomy where the final parts of this work were done. The authors thank the anonymous referee for directing their attention to the possible role of electron Bremsstrahlung emission.
warning/0004/astro-ph0004333.html	ar5iv	text	# On the supernovae heating of intergalactic medium ## 1 Introduction Hierarchical models of structure formation have been very successful in explaining many observed properties of galaxies and galaxy clusters. Nevertheless, some puzzling problems remain open. Several theoretical studies have demonstrated that some heating of gas, in addition to the heating during the gravitational collapse, is required to explain the observed properties of the intracluster medium (ICM). ? first showed that an early injection of energy results in correlations and evolution of bulk cluster properties (X-ray luminosity, gas temperature, etc.) that match observations. This conclusion was backed by numerical simulations: ? was able to get a better fit to the X-ray luminosity of the Coma cluster in his cosmological simulation by preheating the gas to $`10^7`$ K ($`0.9`$ keV), while ? showed that pre-heated clusters matched the observed slope of the correlation between X-ray luminosity and temperature. Recent semi-analytical studies of cluster evolution have reached similar conclusions (????). The exact amount of required energy injection depends on the epoch, and have been argued to be in the range of $`0.53`$ keV per gas particle (????). Although the problem has been identified, it is not yet clear what processes can provide the required heating. It is clear that identification of these processes is crucial for a complete picture of cluster formation. The candidate process which has been discussed most is supernovae-driven galactic winds. The gas of the galactic interstellar medium (ISM) can be heated by supernovae explosions and acquire energy comparable to or larger than its gravitational binding energy. This heated gas can then flow away and result in additional heating if the winds result in shocks when they encounter intergalactic gas. Although there is observational evidence for such winds in present-day galaxies (e.g., ?), theoretical models of winds are rather ill-constrained due to uncertainties in the efficiency of conversion of supernovae (SNe) explosion energy into thermal energy of the gas and other details. Early estimates of possible energy input from SNe based on observed metal abundances showed that SNe are plausible candidates (e.g., ???). However, in these estimates it was assumed that the distribution of metals in the ICM is uniform and that the efficiency with which energy of SNe explosions can be converted into thermal energy of the gas is close to $`100\%`$. Therefore, there is a need for detailed estimates using new measurements of metal abundances in clusters and current galaxy formation models which have become much more advanced and sophisticated in the last several years. Evaluations of SNe as a heating source have recently been performed by ? and ? using a semi-analytical approach to galaxy modelling. These authors conclude that it is unlikely that supernovae are the only source of heating. ? argue then that radiation from quasar population can provide the required heating much more easily. In this paper we repeat previous estimates of the possible SNe energy input using updated values of observed ICM abundances and relaxing the assumption of abundance uniformity, motivated by recent observations (?; ?). We also make a separate estimate for the cluster AWM7, for which the radial abundance gradient was measured. The details of this estimate are presented in § 2. We complement this analysis with direct estimates of the energy input by counting the total number of supernovae exploded in all cluster galaxies throughout their evolution in self-consistent three-dimensional gasdynamical simulations of galaxy formation. The simulations include cooling, star formation, SNe feedback and a multi-phase model of the interstellar medium in galaxies and have been shown to match many fundamental observed correlations of galactic properties such as the galaxy luminosity function, the Tully-Fisher relation and its scatter, the color-magnitude sequence, and, perhaps most importantly, the evolution of the global star formation rate in the Universe. The details of the simulations are described in § 3. The energy input estimates are presented and compared in § 4 and discussed in § 5. We summarize our main results and conclusions in § 6. ## 2 Supernovae energy input from observed ICM metallicities We will first estimate the energy input from SNe to the ICM using observed metallicities of the cluster gas. We have based the estimate on silicon (Si) and iron (Fe) abundances because these two elements have been most accurately measured for a large sample of galaxy clusters (??). We use average ICM metallicities quoted in Table 1 and photospheric solar abundances of ? ($`n_{Fe}/n_H=4.68\times 10^5`$ and $`n_{Si}/n_H=3.55\times 10^5`$). Given that the mass of an element $`X_i`$, $`M_{X_i}`$, within the cluster virial radius is known, the number of SNe type I and II required to produce this mass is equal to $`f_{X_i}^IM_{X_i}/y_I(X_i)`$ and $`(1f_{X_i}^I)M_{X_i}/y_{II}(X_i)`$, respectively. Here $`f_{X_i}^I`$ is mass the fraction of the element contributed by type I SNe, $`y_I(X_i)`$ and $`y_{II}(X_i)`$ are the mass-weighted yields of the element $`X_i`$ by SNe type I and II respectively. The SNe energy input can then be obtained by multiplying the number of SNe by the energy transferred to the gas during each SN explosion: $$E_{SN}M_{X_i}\left(f_{X_i}^I\frac{ϵ_{SNI}E_{SNI}}{y_I(X_i)}+(1f_{X_i}^I)\frac{ϵ_{SNII}E_{SNII}}{y_{II}(X_i)}\right),$$ (1) where we denote $`E_{SNI}`$ and $`E_{SNII}`$ energies released in explosion of the two types of SNe, and $`ϵ_{SNI}`$ and $`ϵ_{SNII}`$ are fractions of the released energy left after the radiative losses during and after the explosion, which can be actually transferred in the form of thermal and kinetic energy to the ambient gas and lead subsequently to the increase of its entropy. There is a varying degree of uncertainty in our knowledge of the above parameters. First of all, our estimate of the mass of an element depends on the assumption about uniformity of the observed metallicities. If strong radial abundance gradients exist in clusters, the observed metallicity is emission-weighted and therefore corresponds to the metallicity in the cluster core. Numerical simulations of cluster formation (??) that include modelling of galaxy feedback predict the existence of strong radial metallicity gradients in the ICM, as well as patchy spatial distribution of metals. At present, however, it is not clear whether large-scale abundance gradients are universal in clusters. Although abundance gradients have been observed in several clusters (see, e.g., ?? and references therein), these are usually clusters that have a central cD galaxy and exhibit signatures of a central cooling flow (?). The spatial extent of the observed gradients coincides with that of the cooling flow region. It is thus unclear whether such central gradients imply the existence of a larger-scale gradient or they are simply due to the presence of a central cD galaxy and cooling flow. Currently, abundances in the fainter, outer parts of clusters can be measured only for bright nearby systems. In a recent study, ?, found strong large-scale metallicity gradients in the nearby cluster AWM7. The observed iron abundance in this cluster decreases from $`0.5\pm 0.05`$ solar within the central $`60h^1\mathrm{kpc}`$ to $`0.2\pm 0.2`$ at $`300500h^1\mathrm{kpc}`$. The radially averaged gradient can be well fitted by a $`\beta `$-model with a core radius equal to that of the gas and $`\beta =0.8`$. The Ezawa et al. measurement was the first in which the abundance gradient has been found far beyond the cluster core radius. It is not yet clear how common such large-scale gradients are. ? show that large-scale metallicity gradient are indeed observed in many clusters. It is clear, however, that strong large-scale gradients are not universal; for example, no strong gradient was detected in the Coma cluster (?). Clusters may therefore exhibit a variety of metal distributions and span a range in the ICM metallicities. In support of this, ?, present evidence that clusters without strong metallicity gradients have systematically lower metal abundances. For our purposes it suffices to consider two possible extreme assumptions about the metal distribution. In reality the mass of metals will likely lie in between the masses computed under these assumptions. The first assumption is that the metallicity of the ICM is spatially uniform. Observationally, the metallicity derived from a spatially unresolved spectrum is emission-weighted. It is clear then that if a strong metallicity gradient is present in a cluster, the total mass of metals may be significantly overestimated under the assumption of spatial uniformity. Our second assumption is that the metallicity gradient of the form observed by Ezawa et al.: $`Z(r)=Z_0[1+(r/r_c)^2]^{3\beta /2}`$ is a universal property of the cluster ICM. We will assume $`r_c=100h^1\mathrm{kpc}`$ and $`\beta =0.8`$ (?) for all clusters, normalizing $`Z_0`$ to a value such that $`Z(r_c)`$ is equal to the observed value of metallicity. Most of the cluster emission comes from radii $`<2r_c`$, providing an approximate way to account for the emission-weighting of the metallicity. The core radius of $`100h^1\mathrm{kpc}`$ is larger than the best fit value for the cluster AWM7 but is closer to a typical core radius of the gas distribution in rich clusters. In addition to the estimate for the whole range of cluster masses, we will present the SN energy input for the specific case of AWM7 for which the abundance gradient has been observed and its parameters measured. Another source of uncertainty is the relative importance of type Ia SNe (parameter $`f_{X_i}^I`$) in the metal enrichment of the ICM (???). Therefore, in the case of iron, we will treat $`f_{Fe}^I`$ as a free parameter and calculate $`E_{SN}`$ for values ($`f_{X_i}^I=0.0,0.5,1.0`$). Silicon is a special case, because $`y_{II}(Si)y_I(Si)`$. This renders $`E_{SN}`$ almost insensitive to a particular choice of $`f_{Si}^I`$. This insensitivity, together with the fact that silicon abundance was fairly accurately measured by ? and ?, effectively reduces the uncertainties and thus makes Si a very useful element for our estimate. The third major source of uncertainty is the yields predicted by different theoretical models of SN explosions (?). The yields of SNe type II may depend on the initial metallicity of SNe, input physics of a model, and other factors. In our estimate we will use yields calculated by ? for metal-poor ($`Z/Z_{}=10^4`$) SNe with explosion energy of $`1.2\times 10^{51}\mathrm{ergs}`$ (model A) and metal-rich ($`Z/Z_{}=1`$) SNe with explosion energy of $`1.2\times 10^{51}\mathrm{ergs}`$ for SN of mass $`25\mathrm{M}_{}`$ and $`2\times 10^{51}\mathrm{ergs}`$ for SN of mass $`>25\mathrm{M}_{}`$ (Model B). Model B has a higher explosion energy for very massive stars to reduce the effects of reimplosion of explosively synthesized ejecta, thus increasing the yields. The yields for these models approximately give the lower and upper limits of the current theoretical predictions (see ???). Given the wide spread in predictions of theoretical models, the supernovae energy input estimates from the observed metallicities have the uncertainty of up to $`50\%`$, in addition to other possible uncertainties. We calculate the average yield of SNII by averaging the mass-dependent yields with a stellar initial mass function (IMF): $$y_{II}(X_i)=\frac{_{m_l}^{m_u}y_{II}(X_i,m)\varphi (m)𝑑m}{_{m_l}^{m_u}\varphi (m)𝑑m}.$$ (2) For the IMF we assume the (?) function, $`\varphi (m)m^{2.35}`$, with the lower mass limit of $`m_l=12M_{}`$ and $`m_l=11M_{}`$ for the two yield models<sup>2</sup><sup>2</sup>2We choose not to use extrapolation and use mass limits of the yield grid of ?. This does not significantly affect the average yields. ? and ? have used a somewhat larger range of masses and obtained similar average yields. and the upper mass limit of $`m_u=40M_{}`$. The SN with masses $`<11M_{}`$ do not contribute significantly to the metal enrichment, while stars of mass $`>40M_{}`$ are rare. Our choice of the Salpeter IMF does not affect the average yields significantly: averaging with considerably shallower ($`\varphi (m)m^2`$) and steeper ($`\varphi (m)m^{2.7}`$) IMFs results in average yields that differ by less than $`10\%`$ from the values for the Salpeter IMF. For SN type Ia, the yields appear to be independent of the SN mass. We use SNIa yields (see Table 1) predicted by the W7 model of ?, which is a model of simple deflagration. The yields for the supernovae of type Ia and type II used in our analysis are summarized in Table 1. The energy released in a SN explosion ($`E_{SNI}`$ and $`E_{SNII}`$) also depends on a variety of factors (e.g. the mass of supernova). However, it can vary only by a factor of $`2`$ and we will therefore assume for simplicity that $`E_{SNI}=E_{SNII}=1.2\times 10^{51}\mathrm{ergs}`$ (e.g., ?). Not all of this initial kinetic energy of explosion is retained by the ejected gas. Analytical arguments (??) and recent numerical simulations (?) suggest that at most $`10\%`$ of the initial kinetic energy of explosion can be ultimately transferred to the ambient gas. In particular, ? have numerically studied the evolution of the ejected material for a variety of densities and metallicities of the ambient gas and concluded that regardless of the ambient density and metallicity, $`\genfrac{}{}{0pt}{}{_>}{^{}}90\%`$ of the initial energy acquired by ejecta during the explosion is lost to radiation (see, however, discussion in § 5). Therefore, we will assume that only $`10\%`$ of the explosion can actually be transferred to the ambient gas: i.e., $`ϵ_{SNI}=ϵ_{SNII}=0.1`$. Here, we neglect the dependence of $`ϵ`$ on environment, metallicity, and possible systematic differences between $`ϵ_{SNI}`$ and $`ϵ_{SNII}`$. However, according to ?, radiation losses are $`90\%`$ or more for most of the realistic environments and metallicities and the value of $`ϵ=0.1`$ is a reasonable upper limit for both types of SNe. We note that the energy estimates from observed metallicities presented in Figs. 45 can be simply linearly scaled up or down for other values of $`ϵ`$. Clearly, the energy of SNe explosions is actually released into the interstellar medium and needs then to be somehow transferred to the IGM. Even if such transfer is possible, it is likely that it would result in additional energy losses. The estimates we make should therefore be considered as the upper limits on the amount of the energy that could have been available for the IGM heating. All of the estimates are made assuming the low-density flat cold dark matter model with cosmological constant ($`\mathrm{\Lambda }`$CDM). The contributions of baryons, cold dark matter, and vacuum energy are: $`\mathrm{\Omega }_b=0.05`$, $`\mathrm{\Omega }_m=0.25`$, $`\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$, respectively. We assume a Hubble constant of $`H_0=70\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$. The cluster virial radius and mass for this model are defined at the overdensity of $`334`$ and we assume the baryon fraction within the virial radius is $`f_b=\mathrm{\Omega }_b/\mathrm{\Omega }_m0.17`$. ## 3 Energy input from supernovae in numerical simulations of galaxy formation The question we now ask is what energy, $`E_{SN}^g`$, can be expected to have been released in all the SNe in cluster galaxies throughout their evolution? This question can be answered only in the framework of a self-consistent model of galaxy formation. There are currently two independently developing approaches to modelling galaxy formation and evolution: semi-analytic models (SAMs; e.g., ??? and references therein) and numerical models (e.g., ???). In this section we will make an estimate of $`E_{SN}^g`$ using numerical simulations of galaxy formation. The numerical techniques and physical ingredients of the model are described in ?. The model includes a self-consistent treatment of the dark matter and baryonic components and effects of cooling, star formation, and SNe feedback. Simulations include a multi-phase model of interstellar medium. Note that these simulations account only for SN type II so the contribution from SNI is therefore neglected in the estimate of $`E_{SN}^g`$. The ideal simulation for our purpose would be a full modelling of cluster formation that would include formation of the cluster galaxies, their starformation and feedback. However, with the numerical code used here, this would require a significant sacrifice in the spatial dynamic range and mass resolution and would make it impossible to follow reliably the starformation and feedback processes. Such simulation awaits future higher dynamic range simulations using adaptive mesh refinement technique. We choose the following compromise. We use many small-box galaxy formation simulations to determine statistically the number of supernovae, $`N_{\mathrm{SN}}`$, that is expected to explode in a galaxy of a given absolute magnitude, $`M_\mathrm{B}`$. We then use this relation and assume a galaxy luminosity function in clusters to estimate how many supernovae could have exploded in a cluster of a given mass. The energy released in these SN explosions would provide an upper limit on the amount of energy available for IGM heating. While the number of SNe could be estimated by assuming a particular $`N_{\mathrm{SN}}M_\mathrm{B}`$ relation, the use of simulations in this study spares us from making this additional assumption. As we describe below, the simulations reproduce many of the observed galactic properties which provides support to the used $`N_{\mathrm{SN}}M_\mathrm{B}`$ relation. The simulations of the COBE-normalized $`\mathrm{\Lambda }`$CDM model ($`\mathrm{\Omega }_0=1\mathrm{\Omega }_\mathrm{\Lambda }=0.35`$; $`\mathrm{\Omega }_b=0.026`$; $`H_0=70\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$, where $`\mathrm{\Omega }_0`$, $`\mathrm{\Omega }_b`$, and $`H_0`$ are present day values of the matter and baryon densities and the Hubble constant, respectively) used here are described in ?. A total of 11 simulations were run from different realizations of initial conditions. The size of the simulation boxes was fixed to $`L_{box}=3.5h^1\mathrm{Mpc}=5\mathrm{M}\mathrm{p}\mathrm{c}`$ and the simulations were run using $`128^3`$ grid cells and particles which gives mass and spatial resolution of $`2\times 10^6h^1\mathrm{M}_{}`$ and $`27h^1\mathrm{kpc}`$, respectively. A total of $`240`$ galaxies, $`140`$ of which have $`M_B(z=0)<14`$, were formed in all the runs combined. The observed color-magnitude diagram, luminosity function (LF), and Tully-Fisher relation of low-redshift galaxies are reproduced well by the simulated galaxies (??). Simulations used here were done assuming SN feedback parameter of $`A=200`$ (see ?). This value of the parameter means a moderate efficiency of supernovae feedback and, correspondingly, relatively high star formation rate. The redshift dependence of the global star formation rate averaged over all simulations is shown in Fig. 1 together with the current data on observed global star formation history in the Universe (see also ? for comparisons of other cosmological models). The observational data were collected from ????????????, and ?. All data points correspond to measurements of comoving UV or $`H\alpha `$ luminosity densities. In order to transform to star formation densities, we have followed Madau’s prescription (?) to correct the original data for dust extinction and to transform the luminosity densities to star formation densities. All data points were properly rescaled to the $`\mathrm{\Lambda }`$CDM cosmological model used in the numerical simulations. The figure shows that the simulation results are in agreement (within the errors) with the observed evolution of the global star formation rate. The star formation rate in the simulations may actually be a little higher than the average observed rate, implying hence a larger number of exploded SNe. As we derived the SN rate from the galactic star formation rate shown in Fig. 1, this figure may serve as an illustration of how the SN explosion rate evolves with time. We have analyzed two additional simulations to assess the effects of resolution and box size. Particularly, the effects of resolution were checked by re-running one of the $`5`$ Mpc simulations with $`256^3`$ grid cells and particles (i.e., with eight times better mass resolution and twice the dynamic range). We have not found any significant changes in the global star formation rate or in the predicted number of supernovae (see below). To test the effects of the box size, we ran a simulation of $`8.4h^1\mathrm{Mpc}=12\mathrm{Mpc}`$ box using $`300^3`$ grid cells and particles, which gives the same resolution as the $`128^3`$ $`5`$ Mpc runs but in a $`2.4`$ times larger box. The results of this simulations are shown in Figs. 1 and 2 together with the results of other runs. The figures show that results of the large-box simulation are in agreement with results of $`5`$ Mpc runs. As we mentioned above, to estimate $`E_{SN}^g`$ expected from galaxies which end up in a cluster, we make use of the correlation between absolute magnitude of a galaxy at z=0 and number of type II SNe exploded in this galaxy throughout its evolution. The number of type II SNe exploded in a galaxy of absolute magnitude $`M`$, $`N_{SN}(M)`$, is computed as the fraction of gas mass converted into stars of mass $`10M_{}`$ divided by the IMF-weighted mean SNe mass. We use the ? IMF, for which these numbers are $`0.12`$ and $`22M_{}`$ (using lower and upper integration limits of $`0.1M_{}`$ and $`125M_{}`$), respectively. Figure 2 shows $`z=0`$ correlation $`N_{SN}(M_B)`$ for galaxies formed in the eleven $`5`$ Mpc and one $`12`$ Mpc $`\mathrm{\Lambda }`$CDM runs. The correlation at $`z=0`$ can be well fitted by a linear fit (shown by solid line) of the form $`\mathrm{log}(N_{SN})=a+bM_B`$, with $`a=1.411`$ and $`b=0.344`$. Figure 3 shows evolution of this correlation with redshift. The figure shows that by $`z=2`$ the number of exploded SNe is predicted to be $`35`$ smaller than the number exploded by $`z=0`$. Note that galaxies in our simulations are either isolated or are located in poor groups. It can be expected that formation of cluster galaxies occurs somewhat earlier than that of galaxies in poorer environments (by about $`\mathrm{\Delta }z1`$, see, e.g., ?) and the results for $`z>2`$ should probably be interpreted as $`z>33.5`$ instead. To estimate the supernovae energy input in a cluster of a given virial mass, $`M_{vir}`$, we convolve $`N_{SN}(M_B)`$ fit with the ? galaxy luminosity function $$\varphi (M)=0.4\mathrm{ln}10\varphi _{}x^{1+\alpha }e^x;x10^{0.4(M_{}M)};$$ (3) where normalization parameter $`\varphi _{}`$ is assumed to be equal to $`\mathrm{\Delta }_{vir}`$ times its field value. The parameter $`\mathrm{\Delta }_{vir}`$ is the expected virial overdensity in a given cosmological model and is $`334`$ for the $`\mathrm{\Lambda }`$CDM model adopted for our estimate (e.g., ??). The energy input is thus $$E_{SN}^g(M_{vir})=E_s\mathrm{\Delta }_{vir}\left(\frac{4\pi }{3}R_{vir}^3\right)\underset{M_f}{\overset{M_b}{}}N_{SN}(M)\varphi (M)𝑑M;$$ (4) where $`E_s=ϵ_{II}E_{SNII}`$ is the energy input of a single supernova explosion, $`\varphi (M)`$ is the field LF, $`R_{vir}`$ is the virial radius of the cluster, $`M_b`$ and $`M_f`$ are the bright and faint limits of integration, and $`ϵ_{II}`$, $`E_{SNII}`$ have the same meaning as in the previous section. The parameters of the luminosity function of galaxies in clusters appear to be similar to those of the field LF (?) and we will therefore neglect possible small differences between cluster and field LFs and cluster-to-cluster variations. We adopt parameters $`M_B=19.5`$ and $`\alpha =1.2`$ of the Schechter luminosity function consistent with recent measurements of $`B`$-band LF in the field (??) and in clusters (?). The faint-end slope $`\alpha `$ is somewhat steeper than in LFs from most of other field surveys (e.g., ???). However, the steep value $`\alpha =1.2`$ better matches the LF of cluster galaxies and the faint end slope of the LF of the simulated galaxies (see ?). Therefore, we adopt this value in our analysis along with the normalization of the field LF $`\varphi _{}^{field}=0.02h^3\mathrm{Mpc}^3`$ (?). This value may be uncertain by a factor of two (see Table 1 in ?). The $`E_{SN}^g`$ estimate presented below is proportional to $`\varphi _{}`$ and can be simply rescaled for other values. We use the integration limits $`M_B^b=22`$ and $`M_B^f=14`$. The results are insensitive to adopting a brighter $`M_b`$ or a fainter $`M_f`$. For consistency, we use $`ϵ_{II}=0.1`$ and $`E_{SNII}=1.2\times 10^{51}\mathrm{ergs}`$ adopted in the previous section. ## 4 Results Figures 4 and 5 show results of the estimates described in the previous two sections. In Figure 4 we compare estimated energy input from SNe with the thermal energy of the ICM gas. Top row of Fig. 4 shows $`E_{SN}`$ estimate using observed ICM metallicities and model A for SNII yields, while the bottom row shows the same estimate for yield model B (see Table 1). The thermal energy of the gas is computed as $$E_{th}=6\pi \frac{k}{\mu m_p}\underset{0}{\overset{R_{vir}}{}}\rho _g(r)T(r)r^2𝑑r,$$ (5) where $`k`$ is the Boltzmann constant, $`m_p`$ is the mass of proton, $`\mu =0.6`$ is the assumed mean molecular weight of the ICM plasma, $`\rho _g(r)`$ and $`T(r)`$ are its radial density and temperature profiles. A density profile $`\rho _g(r)`$ is assumed, and the temperature profile is calculated from the equation of hydrostatic equilibrium. In Fig. 4, we assume that gas is distributed similarly to dark matter, and is described by the ? (hereafter NFW), functional form, $$\rho (r)=\frac{\rho _s}{(r/r_s)(1+r/r_s)^2}$$ (6) with appropriate scaling of parameter $`c=R_{vir}/r_s`$ with cluster virial mass (NFW). The observed distribution of the ICM gas is more often described by the $`\beta `$-profile: $$\rho (r)=\frac{\rho _0}{[1+(r/r_c)^2]^{3\beta /2}},$$ (7) where $`r_c`$ is the core radius, and the parameter $`\beta `$ controls the outer slope of the distribution. The thermal energy of gas distributed with the above density profile (for values of $`r_c`$ and $`\beta `$ consistent with the observed range) is only $`1020\%`$ higher than $`E_{th}`$ of the NFW-distributed gas, the difference indistinguishable on the scale of Fig. 4. Figure 4 shows that expected energy input from SNe is, depending on the assumed uniformity of the ICM metallicity, $`510\%`$ of the gas thermal energy for poor clusters ($`M_{vir}10^{14}h^1M_{}`$) and $`\genfrac{}{}{0pt}{}{_<}{^{}}5\%`$ for rich clusters ($`M_{vir}10^{15}h^1M_{}`$). In case of the strong metallicity gradients these numbers are $`3\%`$ and $`\genfrac{}{}{0pt}{}{_<}{^{}}1\%`$, respectively. The estimates from both Si and Fe agree very well between each other, for $`f_{Fe}^I\genfrac{}{}{0pt}{}{_<}{^{}}0.5`$. Supernova energy input estimated from the numerical simulations agrees well with $`E_{SN}`$ derived from observed metallicities in the case where a strong metallicity gradient is allowed. This implies that if metallicity gradients exist in clusters, simple galaxy formation models with the Salpeter IMF will have no difficulty in accounting for the observed amount of metals in clusters. Figure 4 also shows estimates for the cluster AWM7 for which a large-scale abundance gradient has been observed and its parameters measured (?). The estimate was made using the observed metallicity gradient: $`Z(r)=Z_0[1+(r/r_c)^2]^{3\beta _Z/2}`$ with $`Z_0=0.59`$ and $`\beta _Z=0.8`$, where the core radius $`r_c=57.5h^1\mathrm{kpc}`$ is equal to that of the gas distribution. The gas distribution is described by the similar $`\beta `$-profile with $`\beta _g=0.58`$ (see ? for details). We have made estimates for different gas fractions, $`f_b`$, within the cluster virial radius but the results are only mildly sensitive to a particular value of $`f_b`$; the estimates shown in figs. 4 and 5 were made assuming $`f_b=0.2`$. The energy input estimate for AWM7 lies even lower than the solid lines at this mass because to calculate the latter we have assumed a core radius of $`r_c=100h^1\mathrm{kpc}`$, which results in a higher mass of metals and consequently larger estimated $`E_{SN}`$. For reference, table 2 gives predictions for the numbers of type II SNe in clusters of different masses based on our model (line 1), as well as the total masses of Fe and Si inferred from observations with assumptions of metallicity gradient and uniform metal distribution (lines 3-6). The table also gives the number of SNe required to deposit 1 keV per gas particle into the ICM (line 2), estimated assuming that average supernova deposits $`1.2\times 10^{50}`$ ergs. The mass of metals predicted in our model can be easily obtained by multiplying the number of SNe in line 1 of table 2 by the corresponding mass-weighted yield given in Table 1. Note, however, that the predicted number refers to the type II SNe only, and the predicted mass of metals is thus only due to SNIIe. The question we would ultimately like to address is whether the energy input from SNe can noticeably affect the thermal state of the ICM. To answer this question, we need to know what energy input is needed to account for the observed properties of the ICM. It is not completely clear what energy is required. However, on theoretical grounds (??) it is known that model predictions are in better agreement with the data when gas is assumed to be preheated (by some non-gravitational process) at an early moment. The preheating results in gas evolution corresponding to a higher adiabat, which affects the evolution of the accreted gas (in particular, some of the accreted gas may avoid being strongly shocked). Numerical simulations (????) and semi-analytic models of cluster evolution (?????) confirm that preheating results in cluster properties that are more in accord with observations. For example, to simulate SNe heating ? preheats the gas in his gasdynamic simulations by injecting 1 keV of energy per nucleon of gas, or, for plasma with primordial composition, $`0.5\mathrm{keV}`$ per gas particle. ? assume in their model that SNe preheat the intergalactic gas to temperatures of $`0.50.7`$ keV, which corresponds to $`(3/2)kT0.751.0\mathrm{keV}`$ per gas particle. ? and ? argue based on their semi-analytic calculations that the energy injection of $`23`$ keV per particle is required to bring model predictions in accord with observations. We can compare our estimate of $`E_{SN}`$ to these numbers calculating the energy per gas particle as $`E_{SN}/N_p`$, where $`N_p=f_bM_{vir}/(\mu m_p)`$ is the number of gas particles within the virial radius of cluster. Figure 5 shows $`E_{SN}/N_p`$ for $`E_{SN}`$ estimated from Si and Fe and from galaxy formation simulations. The figure shows that the maximum energy per particle of $`0.1\mathrm{keV}`$ can be injected by SNe if the ICM metallicity is homogeneous, while in the case of a strong metallicity gradient the typical energy per gas particle is only a few tens eV. In particular, the estimate of $`E_{SN}`$ for the cluster AWM7 is only $`0.0020.01`$ keV per particle. The corresponding estimate from galaxy formation simulations is $`10^2\mathrm{keV}`$ per particle. These numbers are $`520`$ times smaller than the typical energy injection assumed in the cluster formation models quoted above. ## 5 Discussion The results presented in § 4 allow us to assess the conditions required for the SNe energy input to be important in galaxy clusters. The primary conditions that are implied by our estimate of $`E_{SN}`$ from the observed abundances of Si and Fe are (i) large-scale uniformity of the metal abundances throughout the cluster volume and (ii) near $`100\%`$ efficiency in transfer of the energy of SN explosion to the thermal energy of the IGM gas. The latter assumption is rather unlikely and the energies derived from the observed abundances should therefore be considered as the upper limits on the amount of SN energy that could have heated the IGM. There are but a few theoretical predictions and observational data concerning the degree of uniformity of the metal distribution in clusters. Based on the numerical simulations that include galaxy feedback and metal enrichment, ? and ? predict that large-scale metallicity gradients should exist in clusters. On the observational side, ? observed such a gradient in cluster AWM7. More recently, ? reported similar large-scale ($`R\genfrac{}{}{0pt}{}{_<}{^{}}0.51\mathrm{Mpc}`$) metallicity gradients detected using ASCA observations for several other clusters. It is not clear, however, whether such gradients are ubiquitous. Our estimates of energy input from observed metal abundances differ by a factor of $`510`$ if we assume a uniform distribution of metals versus metallicity gradients of the type observed in AWM7. New, deep observations of ICM metallicity profiles are therefore crucial to make this estimate much more reliable. With the launch of the Chandra X-ray satellite, such observations should become available. Our estimate of the SN energy input for AWM7 is two orders of magnitude lower than energy input which seem to be required to sufficiently preheat the ICM gas. Incidentally, the existence of large-scale abundance gradients in clusters would solve the problem of the total iron mass in clusters. ? and ? show that if the contribution of type I SNe to the iron production in clusters is relatively small, the total iron mass in the ICM is too large to be explained by type II SNe produced with Salpeter IMF. ? argue, however, that this solution is unattractive because it makes it difficult to explain the metallicities and radial abundance gradients in massive elliptical galaxies. It is clear from our analysis that the existence of large-scale abundance gradients in the ICM can reduce the estimate of the iron mass by up to an order of magnitude, thereby eliminating the need for a large number of SNII and flatter IMF. The predictions of the SNe energy input, $`E_{SN}^g`$, of the numerical simulations of galaxy formation presented in this paper, although consistent with observed evolution of the global starfomation rate in the Universe, are somewhat lower than the estimate from the metal abundances, $`E_{SN}^m`$. The estimates $`E_{SN}^g`$ and $`E_{SN}^m`$ agree reasonably well if a metallicity gradient is assumed and the contribution of type Ia SNe to the iron enrichment is $`>50\%`$. In particular, the $`E_{SN}^g`$ estimate is actually higher than estimate of $`E_{SN}^m`$ for AWM7. However, the energy input in this case is of the order of $`1050`$ eV per gas particle, which is far short of the energy injection typically assumed to bring theoretical models in accord with observations: $`0.52`$ keV per particle. The estimate will still be short by a factor of $`510`$ even if $`100\%`$ of the energy of every SN explosion goes into heating the ICM gas. The above conclusions are for clusters of virial mass $`M_{vir}\genfrac{}{}{0pt}{}{_>}{^{}}10^{14}h^1\mathrm{M}_{}`$. Figure 4 shows that the ratio of predicted SNe energy input to the thermal energy of the ICM gas increases by about an order of magnitude as the mass is decreased from $`10^{15}h^1\mathrm{M}_{}`$ to $`10^{14}h^1\mathrm{M}_{}`$. This trend means that the SNe energy input may be much more important for clusters of mass $`M_{vir}\genfrac{}{}{0pt}{}{_<}{^{}}5\times 10^{13}h^1\mathrm{M}_{}`$ than for more massive clusters<sup>3</sup><sup>3</sup>3Note that this conclusion depends on our assumption that total luminosity of stars in clusters is proportional to the cluster mass (See § 3). Although this assumption is reasonable, there is evidence that mass-to-light ratio of clusters and groups is a function of system mass. The data indicates that mass-to-light ratio of galaxy groups is somewhat smaller than that of clusters (e.g., ?). In this case our conclusion would not be changed.. The mass $`5\times 10^{13}h^1\mathrm{M}_{}`$ corresponds to the ICM temperature of $`2`$ keV (e.g., ?), while deviations from non-similarity are observed in real clusters for temperatures of $`\genfrac{}{}{0pt}{}{_<}{^{}}2`$ keV (??). Nevertheless, it appears that quantitatively our conclusions will stand for poor clusters. The entropy of the preheated gas required to explain observations is $`100\mathrm{keV}\mathrm{cm}^2`$ (?) which corresponds to an energy of $`1.5(n_e/10^3cm^3)^{2/3}`$ keV per particle, where $`n_e`$ is electron number density. Thus, the energy injection into the gas in cluster cores ($`n_e10^3\mathrm{cm}^3`$) is about $`1.5`$ keV per particle. Semi-analytical calculations of ? and ? show that the energy injection required to explain the data may be even higher: $`23`$ keV per particle<sup>4</sup><sup>4</sup>4? give an estimate of the required entropy of preheated gas as $`S=kT/(\mu m_H\rho ^{2/3})3.7\times 10^{33}\mathrm{ergs}\mathrm{g}^{5/3}\mathrm{cm}^2`$. This corresponds to $`E_p=1.5\mu m_H\rho ^{2/3}S3.5`$ keV per particle if we assume a typical density of gas in cluster cores ($`\rho 10^{27}\mathrm{g}\mathrm{cm}^3`$) and the above value of entropy.. Such energy input is marginally consistent with our $`E_{SN}^m`$ estimate in the case of uniform metallicities and $`ϵ1`$. For the case of metallicity gradients, $`E_{SN}^m`$ is more than an order of magnitude lower. The energy input, $`E_{SN}^g`$, predicted from numerical simulations is even lower and is $`0.1`$ keV per particle even for $`ϵ=1`$. Therefore, the conclusion we draw from this analysis is that it is unlikely that the energy input from SNe is sufficient to preheat the intracluster gas to the required entropy, unless all of the explosion energy goes into heating of the gas and metal abundances are uniform throughout the ICM. Moreover, in light of the $`E_{SN}^g`$ estimates, the SN energy input can only be important if starformation rate in cluster environments is a factor of 10 higher than the average cosmic rate. Similar conclusions were reached by ?, ?, and ?. Recently, ? have also used observed abundance of Si in the ICM to estimate possible SNe heating and found that the implied SNe energies would not be sufficient to heat the entire cluster gas to the required levels (note that this estimate was done assuming uniform distribution of Si and $`ϵ_{SN}=1`$). He pointed out, however, that SNe could still be the source of heating if only the gas in cluster cores was heated. In this case, heating would have to occur after or during formation of a cluster, not at early epochs as was assumed previously, but sufficiently early enough to be consistent with lack of evolution of metal abundances at lower redshifts ($`z\genfrac{}{}{0pt}{}{_<}{^{}}1`$; ?). Details and quantitative predictions of such a scenario are yet to be worked out. It is obvious that there are a number of uncertainties in our estimates of the SNe energy input. The estimates of $`E_{SN}^m`$ made with the assumption of uniform ICM metallicity are by a factor $`310`$ higher than the corresponding estimates in the case when a strong metallicity gradient is assumed. This uncertainty not only makes the $`E_{SN}^m`$ estimate uncertain, but also hinders comparisons of metal abundances predicted by galaxy formation models with observations. This will likely be resolved in the near future with the advent of new, deep X-ray observations of clusters, but it is a major limitation at present. Currently, only one robust measurement of large-scale metallicity gradient has been obtained (?). This cluster, AWM7, confirms the existence of strong metallicity gradients and the estimate of $`E_{SN}^m`$ for this particular cluster supports our conclusions. It is not clear, however, how universal such gradients are in clusters. Note also that our estimates are based on average abundances of Si and Fe from a large sample of clusters. Abundances in individual clusters may vary by a factor of $`\genfrac{}{}{0pt}{}{_>}{^{}}23`$. Thus, for example, abundances of Si and Fe (in solar units) vary in the range $`0.11`$ and $`0.150.45`$, respectively (??). The energy estimates for individual clusters may therefore also vary by a corresponding factor. The theoretical yields of Si and Fe from type Ia and type II SNe used in our analysis depend on specifics of the explosion model. The Si yields from SNIa may be uncertain by a factor of two (??), while all models predict similar (to $`10\%`$) yields of iron. The yields of SNII for Si and Fe vary by $`3040\%`$ between different models (e.g., ?). Yield models A and B used in our analysis approximately represent the range of predictions and should therefore provide a fair estimate of uncertainty. Our conclusions hold for both yield models. The fraction of supernova explosion energy that can be available for gas heating is also rather uncertain. ? and ? give analytical arguments that this fraction should be $`0.1`$. These arguments are supported by recent direct numerical simulations of ? who studied radiative losses of a SN remnant (SNR) for a grid of densities and metallicities of the ambient gas. The arguments and simulations, however, assume spherically symmetric evolution of SNRs in ambient gas of uniform density. The efficiency may be higher if the topology of ambient gas density is very assymetric and the gas has been swept up and preheated by previous, recently exploded SNe (?). This, for example, may be the case during a strong starburst (e.g., ?). The parameter $`ϵ`$ is thus likely to be environment dependent and the average value would be determined by the relative number of SNe exploding during periods of quiescent star formation vs. the number of SNe exploding in starbursts. Regardless of the actual value, considerable radiation losses are expected and therefore it seems very unlikely that the efficiency is close to $`100\%`$ ($`ϵ=1`$). Beside the problem of heating efficiency, it is also not clear how the heated interstellar gas and released SN energy is transferred to the IGM (or ICM). Several transfer mechanisms have been suggested. Gas may be blown away from galaxies by supernova-driven winds (???) which subsequently shock the IGM gas. Evidence for winds is indeed observed in some starburst galaxies (e.g., ?). However, only a small fraction of gas is expected to be blown away by starbursts in massive galaxies (e.g., ?, ?) and therefore ejected gas can only constitute a small fraction of ICM. Clearly, the same questions arise when we consider how energy released by SNe can actually heat the IGM. If only a fraction of this energy is delivered to IGM, this effectively means a smaller value of $`ϵ_{SN}`$ and strengthens our conclusions. The gas can also be transferred to the ICM by ram pressure (e.g., ?) and tidal stripping. The efficiency of ram pressure in clusters is not well known. Recent numerical simulations, however, suggest that it may actually be rather low (?). The tidal stripping is probably the most efficient mechanism of delivering ISM gas to the intracluster medium, especially for low surface brightness galaxies (?). However, in this case the gas is transferred to the ICM relatively late, after the epoch of cluster formation, when a sufficiently deep potential well is formed. This is in conflict with high metal abundances observed in high-redshift clusters (??). Recently, ? suggested that metal-enriched gas can be ejected at early epochs during galactic mergers. This mechanism may transfer metal-rich hot interstellar gas into IGM, where it can be further heated by shocks developed during a merger or after an encounter between ejected material and the ambient IGM gas. Despite the abundance of possible processes, it is not clear which process (or combination thereof) is responsible for the transfer of gas from galaxies into the intergalactic medium. It is clear, however, that this question needs to be clarified if SNe are to be considered a viable source of IGM heating. We have made a number of assumptions to estimate $`E_{SN}^g`$ from the galaxy formation simulations. Changing some of these assumptions can change the energy input estimate. First of all, our assumption of Salpeter IMF directly affects the number of SNe per given mass of formed stars. IMFs Flatter than a Salpeter result in a larger number of supernovae and thus in a larger energy input for the same star formation rate. For instance, a 10% flatter slope with respect to Salpeter’s results in a 50% increase in the number of SNe, given the same low-mass limit of the IMF. Indeed, a flatter IMF has been suggested as an explanation for the observed iron abundances in clusters (e.g., ??). However, note that ? argue that flatter IMF is not consistent with the evolution of elliptical galaxies. The number of SNe depends also on the low-mass limit of the IMF, although in a less sensitive manner. Thus, an increase of the lower-mass limit by a factor of 2 (from 0.1 to 0.2 $`M_{}`$) results in an increase factor of 1.3 in the number of SNe. To calculate the number of SNe exploded during a Hubble time in all cluster galaxies we have assumed that the number density of galaxies in clusters is equal $`(1+\mathrm{\Delta }_{vir})`$ times its field value. This means that clusters represent the same fluctuation in number of galaxies as in their total mass. Although this is a reasonable assumption, we note that in the $`\mathrm{\Lambda }`$CDM model (as well as in other low-matter density CDM cosmologies) studied here, a certain amount of anti-bias ($`b0.5`$) is required for the model to be consistent with observed galaxy clustering (???). This anti-bias arises primarily in the densest regions of galaxy groups and clusters (?). For an anti-bias of $`b0.5`$ the number density of galaxies would be two times lower than assumed in our analysis, which would reduce the estimated energy input by a factor of two. We neglected possible differences between the shape of the field and cluster luminosity functions. These differences appear to be rather small for the $`B`$-magnitude LF used here (?), and we therefore think that the uncertainty associated with this assumption is relatively small. A more important assumption is that the global star formation histories of field and cluster galaxies are similar. At present, there is no convincing evidence otherwise. ?, for example, argue that star formation activity in cluster galaxies is not very different from that in the field. They argue, in fact, that field galaxies may produce more stars (and more type II SNe) than cluster galaxies in which the star formation is being gradually turned off after their infall onto cluster. This is in fact consistent with theoretical predictions of ? who present models for the evolution of the SN rate in clusters and the field. They predict that the rate in clusters is higher than in the field only at $`z\genfrac{}{}{0pt}{}{_>}{^{}}3.5`$, while at lower redshifts it is actually lower due to a decreased contribution from SNe in spiral galaxies. Their predictions for the overall starformation rate in clusters are almost an order of magnitude lower than the starformation rate in the simulations presented here at $`z\genfrac{}{}{0pt}{}{_<}{^{}}3.5`$ and are higher at higher redshifts. It seems unlikely, however, that their prediction can account for the required tenfold increase in number of SNe because only a small fraction of SNe in cluster galaxies explode at $`\genfrac{}{}{0pt}{}{_>}{^{}}4`$. We therefore conclude that possible differences in starformation histories between cluster and field galaxies are too small to change our conclusions. Nevertheless, it is known that rich clusters have properties different than if they would have simply had been constructed from massive ellipticals and small galaxy groups (???). Ellipticals and galaxy groups appear to have smaller gas fractions and lower metal abundances than rich clusters do. In particular, the ratio of iron mass in the ICM to the total blue luminosity of cluster galaxies is consistently higher for clusters than for groups (?). It appears also that parameters of the models of elliptical galaxies that are tuned to produce the observed metal abundances in the ICM are inconsistent with abundance measurements in individual ellipticals, the problem which cannot be solved by adjusting the SNIa contribution to the metal enrichment (?). These problems may indicate that an important component is missing in our understanding of the ICM enrichment history and cluster evolution. However, our conclusions about the importance of SNe energy input can only change if the star formation rate in the volume from which the cluster forms is significantly higher at all epochs than that star formation rate in the field. ## 6 Conclusions We have presented estimates of the possible energy input by supernovae into the intracluster medium. Although these estimates are prone to a number of uncertainties, we have defined conditions which determine whether SNe can be a significant source of ICM heating. The following main conclusions can be drawn from our analysis. The SNe energy input, $`E_{SN}^m`$, estimated from observed ICM abundances of Si and Fe is only significant ($`1`$ keV per particle) when we assumed that the distribution of metals in the ICM is uniform (no significant radial gradients) and that $`100`$% of individual SN explosion energy goes into heating the ambient gas followed by negligible cooling ($`ϵ1`$) (see § 2). If large-scale metallicity gradients are assumed in clusters, the estimated energy input is $`0.10.5`$ keV per particle for $`ϵ=1`$ and, correspondingly, $`0.010.08`$ keV per particle for a more realistic value of $`ϵ=0.1`$. As an example, we present estimates of the energy input for the cluster AWM7 for which the abundance gradient has been measured. We find that the observed abundance of iron in this cluster implies a SNe energy input of $`\genfrac{}{}{0pt}{}{_<}{^{}}0.01`$ and $`\genfrac{}{}{0pt}{}{_<}{^{}}0.1`$ keV per particle for $`ϵ=0.1`$ and $`ϵ=1`$, respectively. The energy input, $`E_{SN}^g`$, estimated using self-consistent three-dimensional numerical simulations of galaxy formation which include effects of shock heating, cooling, SN feedback, and multi-phase model of ISM, are $`0.01`$ and $`0.1`$ keV per gas particle for values of efficiency parameter $`ϵ=0.1`$ and $`ϵ=1`$, respectively. These values are somewhat lower than the values of $`E_{SN}^m`$ (but are in good agreement with estimates for the AWM7). Nevertheless, the two estimates agree reasonably well if the existence of large-scale abundance gradients is assumed in clusters. We therefore emphasize the importance of new measurements of large-scale metallicity gradients for testing the theoretical models. Our estimates of the SN energy input in all cases, except the case of uniform ICM abundances and $`ϵ=1`$, fall short of the energy injection of $`0.53`$ keV per particle required to bring theoretical models of cluster formation in accord with observations. This suggests that supernovae are unlikely to be the only source of the IGM heating and should possibly be supplemented (or substituted) by some other heating mechanism. Similar conclusions have been reached in recent studies of ?, ?, and ?. ? propose radiation from quasars as an alternative heating mechanism. This opens discussion of new possible processes for what appears to be a required high-redshift preheating of the intergalactic medium. ## Acknowledgements We would like to thank Anatoly Klypin for useful discussions and comments. A.V.K. was supported by NASA through Hubble Fellowship grant HF-01121.01-99A from the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS5-26555. GY acknowledges support from S.E.U.I.D under project number PB96-0029. The numerical simulations used in this paper were run at the Centro Europeo de Paralelismo de Barcelona (CEPBA).
warning/0004/astro-ph0004286.html	ar5iv	text	# Constraining the Accretion Rate Onto Sagittarius A∗ Using Linear Polarization ## 1 Introduction In this paper we discuss the effect of Faraday depolarization on synchrotron radiation in spherical accretion flow models of low-luminosity galactic nuclei (see also Bower et al 1999ab; Agol 2000). We focus on the radio source Sgr A at the Galactic Center, but our results can also be applied to other systems (see §4). This paper was motivated by a possible detection of linear polarization from Sgr A (Aitken et al. 2000). The Bondi accretion rate onto the supermassive black hole at the center of our galaxy is estimated to be $`10^410^5M_{}\mathrm{yr}^1`$ (e.g., Coker & Melia 1997; Quataert, Narayan, & Reid 1999), implying a luminosity of $`10^{41}\mathrm{ergs}\mathrm{s}^1`$ if the radiative efficiency is $`10\%`$. This is roughly 5 orders of magnitudes larger than the observed luminosity (see Narayan et al. 1998 for a recent compilation). Comparable “discrepancies” are obtained for massive elliptical galaxies in nearby X-ray clusters (e.g., Fabian & Rees 1995; Di Matteo et al. 1999). One explanation for the low luminosity of nearby supermassive black holes is that they accrete via an advection-dominated accretion flow (ADAF), in which most of the dissipated turbulent energy is stored as thermal energy rather than being radiated (e.g., Rees et al. 1982; Narayan & Yi 1994, 1995; Abramowicz et al. 1995). In such models the accretion rate is of order the Bondi rate while the radiative efficiency is extremely small ($`10^6`$ for Sgr A). Another explanation for very low luminosity accreting systems is that the Bondi accretion rate estimate is inapplicable (e.g., Blandford & Begelman 1999; Gruzinov 1999). In particular, numerical simulations of non-radiating accretion flows with small values of the dimensionless viscosity parameter $`\alpha `$ find that the gas density scales with radius as $`\rho r^{1/2}`$ rather than the canonical Bondi/ADAF scaling of $`\rho r^{3/2}`$ (Stone, Pringle, & Begelman 1999; Igumenshchev & Abramowicz 1999, 2000; Igumenshchev, Abramowicz, & Narayan 2000). Narayan, Igumenshchev, & Abramowicz (2000) and Quataert & Gruzinov (2000) explained these simulations in terms of a “convection-dominated accretion flow” (CDAF). In such a flow angular momentum is efficiently transported inwards by radial convection, nearly canceling the outward transport by magnetic fields. This strongly suppresses the accretion of matter onto the black hole. Broad band spectra have thus far had difficulty distinguishing between these explanations for the low luminosity of nearby supermassive black holes. For example, Quataert & Narayan (1999; hereafter QN) showed that accretion at much less than the Bondi rate could produce spectra quite similar to ADAF models. In this paper we show that linear polarization observations in the radio to sub-mm can provide a sensitive probe of the accretion rate onto the black hole, and help distinguish between degenerate spectral models. In the next section (§2) we present simple estimates of the physical parameters of the accretion flow relevant for our analysis. We then discuss Faraday depolarization in spherical accretion flow models of Sgr A (§3). In §4 we compare these predictions with observational constraints on the linear polarization of Sgr A and summarize our results. We also generalize our analysis to other low-luminosity galactic nuclei. Throughout this paper, we focus on accretion models of Sgr A. An unresolved jet or outflow may, however, dominate the observed emission (e.g., Falcke, Mannheim, & Biermann 1993; Lo et al. 1998; Falcke 1999); this is briefly discussed in §4. ## 2 Plasma Parameters for Sgr A Stellar kinematics show that there are $`2.6\times 10^6M_{}`$ within $`0.015`$ pc of the Galactic Center (Eckart & Genzel 1997, Ghez et al. 1998), centered on the radio source Sgr A (Menten et al. 1997). The most plausible explanation is that Sgr A is a $`2.6\times 10^6M_{}`$ accreting black hole. Sgr A is believed to accrete the winds from nearby ($`0.1`$ pc) massive stars (Krabbe et al. 1991). The Bondi accretion rate of these winds onto the supermassive black hole is estimated to be $`10^410^5M_{}\mathrm{yr}^1`$ (e.g., Coker & Melia 1997; Quataert, Narayan, & Reid 1999). If the accretion rate close to the black hole is of order the Bondi value the gas density near $`r1`$ is $`n10^910^{10}\mathrm{cm}^3`$ (since $`v_rc`$ near the horizon). The corresponding magnetic field strength, assuming rough equipartition with the nearly relativistic protons, is $`B2\times 10^3`$G. At such magnetic field strengths, relativistic electrons cooling by synchrotron radiation would have a cooling time much less than the inflow time of the gas. In order to not overproduce the observed radio to sub-mm luminosity of Sgr A, the bulk of the electrons must therefore be marginally relativistic, with $`T_e10^910^{10}`$ K. These plasma parameters ($`n,B,T_e`$) describe Bondi and ADAF models of Sgr A (e.g., Melia 1992, 1994; Narayan, Yi, & Mahadevan 1995; Narayan et al. 1998). In such models the electrons are assumed to be adiabatically compressed from large radii in the accretion flow, with virtually no additional turbulent heating. QN showed that accretion at much less than the Bondi rate could also produce the observed high frequency emission from Sgr A, provided the electrons were much hotter than in standard ADAF models (see their Table 2 and Fig. 8b). A simple explanation for this result can be obtained by applying the Burbidge (1958) estimate to Sgr A. We consider synchrotron emission from a sphere of radius $`R`$ containing relativistic electrons with a temperature $`kT_e=\gamma m_ec^2`$. We take the electron heating rate to be comparable to the net turbulent (magnetic) heating rate. As can be confirmed a posteriori, the synchrotron cooling time is $``$ the inflow time of the gas. The electron energy density is then similar to the magnetic energy density $$n\gamma m_ec^2\frac{B^2}{8\pi }.$$ (1) The frequency of peak synchrotron emission and the synchrotron luminosity are given by $$\nu 0.1\gamma ^2\frac{eB}{m_ec}$$ (2) and $$L\sigma _TcB^2\gamma ^2R^3n,$$ (3) where $`\sigma _T`$ is the Thomson cross section. We express $`n`$, $`\gamma `$, and $`B`$ in terms of $`R`$, $`\nu `$, and $`L`$ (see also Falcke 1996; Beckert & Duschl 1997) $$\gamma 3.2\left(\frac{m_e}{c}\frac{\nu ^4R^3}{L}\right)^{1/7}100,$$ (4) $$n\frac{4}{\gamma ^5\lambda ^2r_e}10^6\mathrm{cm}^3,$$ (5) and $$B\sqrt{8\pi m_ec^2\gamma n}45\mathrm{G},$$ (6) where $`\lambda =c/\nu `$ is the wavelength and $`r_e=e^2/(m_ec^2)`$ is the classical electron radius. For the numerical estimates in equations (4)-(6), we have used the observed values for Sgr A. The peak synchrotron frequency is at $`\nu 10^3`$ GHz with a luminosity of $`L10^{36}\mathrm{ergs}\mathrm{s}^1`$ (e.g., Serabyn et al. 1997). In spherical accretion models, this high frequency emission arises from very close to the black hole, so we have taken $`RR_g10^{12}`$ cm. Equation (15) gives a density close to the black hole of $`n10^6\mathrm{cm}^3`$; the implied accretion rate is then $`10^8M_{}\mathrm{yr}^1`$, three to four orders of magnitude smaller than the Bondi value. The thermal blackbody emission at frequency $`\nu `$ from a sphere of radius $`R`$ is $$L_t=2\pi \nu ^3\gamma m_e4\pi R^210^{37}\mathrm{ergs}\mathrm{s}^1,$$ (7) where the numerical estimate is for our fiducial parameters. This comparison shows that the synchrotron emission becomes optically thin below the peak frequency, near $`\nu 300`$ GHz. At lower frequencies the emission is self-absorbed. The above considerations show that both low ($`10^8M_{}\mathrm{yr}^1`$) and high ($`10^510^4M_{}\mathrm{yr}^1`$) accretion rate models can explain the observed sub-mm “bump” in Sgr A. Such models can be distinguished by comparing the observed brightness temperature and/or radio image as a function of frequency with the theoretical predictions (see, e.g., Özel, Psaltis, & Narayan 2000). This test has been difficult to implement because interstellar scattering significantly broadens the image of Sgr A.<sup>1</sup><sup>1</sup>1Recent detections of Sgr A’s intrinsic size (Lo et al. 1998; Krichbaum et al. 1998) still have sufficient uncertainties that a range of theoretical models are allowed. In the next section we show that the linear polarization of Sgr A at high frequencies provides an additional discriminant. ## 3 Faraday Depolarization The anisotropic index of refraction of a magnetized plasma leads to a frequency-dependent rotation in the position angle, $`\theta `$, of linearly polarized electromagnetic waves, $$\theta =RM\lambda ^2,$$ (8) where $`RM`$ is the rotation measure. This can lead to significant depolarization of intrinsically linearly polarized synchrotron emission. For a “cold” non-relativistic plasma, RM is given by (e.g., Rybicki & Lightman 1979) $`RM`$ $`=`$ $`{\displaystyle \frac{e^3}{2\pi m_e^2c^4}}{\displaystyle 𝑑𝐥𝐁n}`$ (9) $`=`$ $`2.63\times 10^{13}\times {\displaystyle 𝑑𝐥𝐁n\frac{\mathrm{rad}}{\mathrm{m}^2}},`$ where $`d𝐥`$ is the differential path length from the observer to the source. In the Appendix we show that the rotation measure for an ultrarelativistic thermal plasma is given by $`RM_\gamma `$ $`=`$ $`{\displaystyle \frac{e^3}{2\pi m_e^2c^4}}{\displaystyle 𝑑𝐥𝐁n\frac{\mathrm{log}\gamma }{2\gamma ^2}}`$ (10) $`=`$ $`2.63\times 10^{13}\times {\displaystyle 𝑑𝐥𝐁n\frac{\mathrm{log}\gamma }{2\gamma ^2}\frac{\mathrm{rad}}{\mathrm{m}^2}},`$ where $`\gamma =kT_e/m_ec^2`$. A comparable expression is obtained for a power law distribution of relativistic electrons, with $`\gamma `$ replaced by $`\gamma _{min}`$, the minimum Lorentz factor of the electrons (Jones & O’Dell 1977). In what follows, we define $`RM(r)`$ to be the contribution to the net rotation measure from radii within $`drr`$ of radius $`r`$ in the accretion flow. ### 3.1 ADAF/Bondi Models In spherical accretion flow models, higher frequency radio emission arises from closer to the black hole, where the electron temperature and magnetic field strengths are the largest; this is also true for jet models (e.g., Falcke 1999). Özel et al. (2000) show that in ADAF models of Sgr A the synchrotron emission at frequency $`\nu =100\nu _{100}`$ GHz arises from a radius $`r_\nu 20\nu _{100}^{0.9}`$ (see their Fig. 5).<sup>2</sup><sup>2</sup>2Falcke (1999) finds a similar expression for $`r_\nu `$ in the jet model. This radius defines the $`\tau =1`$ surface of the synchrotron emission. For smaller radii the emission is self-absorbed while for larger radii it is optically thin. Faraday rotation is only important for r >rν >𝑟subscript𝑟𝜈r\mathrel{\mathchoice{\vbox{\offinterlineskip\halign{$\m@th\displaystyle\hfil#\hfil$\cr>\crcr\sim\crcr} }}{\vbox{\offinterlineskip\halign{$\m@th\textstyle\hfil#\hfil$\cr>\crcr\sim\crcr} }}{\vbox{\offinterlineskip\halign{$\m@th\scriptstyle\hfil#\hfil$\cr>\crcr\sim\crcr} }}{\vbox{\offinterlineskip\halign{$\m@th\scriptscriptstyle\hfil#\hfil$\cr>\crcr\sim\crcr} }}}r_{\nu}, where the photons “free stream” out of the accretion flow. In Bondi/ADAF models Faraday rotation is so strong that the synchrotron emission is completely depolarized by the plasma within $`drr`$ of $`r_\nu `$, i.e., in the vicinity of the $`\tau =1`$ surface where it is emitted (depolarization and emission are thus virtually co-spatial). Taking $`nr^{3/2}`$ and $`Br^{5/4}`$; the rotation measure scales roughly as $`RMr^{7/4}\gamma ^2`$. The relativistic suppression of the rotation measure is small in all models which have an accretion rate comparable to the Bondi rate, because the electrons must then be at most marginally relativistic (§2). The rotation measure as a function of radius is thus given by (see also Bower et al. 1999ab) $$RM10^{13}r^{7/4}\mathrm{rad}\mathrm{m}^2.$$ (11) In equation (11) we have assumed that the magnetic field is in rough equipartition with the gas pressure, has a significant component along the line of sight, and has a coherence length $`\mathrm{}r`$; for $`\mathrm{}r`$, $`RM`$ is reduced by $`(\mathrm{}/r)^{1/2}`$. The normalization in equation (11) is set by the Bondi accretion rate. The large $`RM`$ in ADAF/Bondi models leads to a significant rotation in the position angle of linearly polarized waves. Photons of frequency $`\nu `$ emitted at radius $`r_\nu `$ undergo Faraday rotation through an angle $$\theta _\nu \lambda ^2RM(r_\nu )10^8\nu _{100}^2r_\nu ^{7/4}10^6\nu _{100}^{0.43}\mathrm{rad},$$ (12) where the last estimate uses Özel et al.’s (2000) fit to $`r_\nu (\nu )`$. These rotation angles are so large that the synchrotron emission in ADAF/Bondi models of Sgr A is completely depolarized by Faraday rotation. For example, in a simple uniform source model, the observed polarization is $`\theta _\nu ^1`$ (Pacholczyk 1970). In general, the observed polarization depends on the rotation measure power spectrum, but is $`1`$ for $`\theta _\nu 1`$ (e.g., Tribble 1991).<sup>3</sup><sup>3</sup>3One way of evading this conclusion is to posit that the magnetic field is sufficiently tangled ($`\mathrm{}r`$) to decrease $`RM`$ to <106radm2 <absentsuperscript106radsuperscriptm2\mathrel{\mathchoice{\vbox{\offinterlineskip\halign{$\m@th\displaystyle\hfil#\hfil$\cr<\crcr\sim\crcr} }}{\vbox{\offinterlineskip\halign{$\m@th\textstyle\hfil#\hfil$\cr<\crcr\sim\crcr} }}{\vbox{\offinterlineskip\halign{$\m@th\scriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr} }}{\vbox{\offinterlineskip\halign{$\m@th\scriptscriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr} }}}10^{6}\rm\ rad\ m^{-2}. This tangling would, however, also eliminate linear polarization. ### 3.2 $`\dot{M}\dot{M}_{\mathrm{Bondi}}`$ If the accretion rate onto Sgr A is much less than the Bondi rate, significant polarization may be observable at high frequencies; we show this using an order of magnitude estimate. For the $`\dot{M}10^8M_{}\mathrm{yr}^1`$ model of §2 the rotation measure calculated using equation (10) is $`RM10^3\mathrm{rad}\mathrm{m}^2`$ near $`r1`$. Moreover, if $`nr^{1/2}`$, as in CDAF models, the magnetic field scales as $`Br^{3/4}`$ and $$RM10^3r^{1/4}\left(\frac{\gamma }{100}\right)^2\mathrm{rad}\mathrm{m}^2.$$ (13) The variation of the electron Lorentz factor with radius is somewhat uncertain, but we expect roughly $`\gamma r^1`$, so that the electrons become non-relativistic by $`r10^2`$. Equation (13) then shows that RM has its maximal value at large radii, $`r10^2`$, where $`RM3\times 10^6\mathrm{rad}\mathrm{m}^2`$. Equation (13) demonstrates that, in contrast to ADAF models, there is no depolarization of synchrotron emission at small radii in models with accretion rates much less than the Bondi rate; $`RM`$ is negligible in the region where the synchrotron emission is produced. Depolarization can still be important, however, because observed photons experience different Faraday rotation at large radii, r >102 >𝑟superscript102r\mathrel{\mathchoice{\vbox{\offinterlineskip\halign{$\m@th\displaystyle\hfil#\hfil$\cr>\crcr\sim\crcr} }}{\vbox{\offinterlineskip\halign{$\m@th\textstyle\hfil#\hfil$\cr>\crcr\sim\crcr} }}{\vbox{\offinterlineskip\halign{$\m@th\scriptstyle\hfil#\hfil$\cr>\crcr\sim\crcr} }}{\vbox{\offinterlineskip\halign{$\m@th\scriptscriptstyle\hfil#\hfil$\cr>\crcr\sim\crcr} }}}10^{2}, on their way out of the accretion flow (e.g., Bower et al. 1999ab). Spatial variation in the rotation measure will depolarize Sgr A at frequencies for which δθ=λ2δRM >π𝛿𝜃superscript𝜆2𝛿𝑅𝑀 >𝜋\delta\theta=\lambda^{2}\delta RM\mathrel{\mathchoice{\vbox{\offinterlineskip\halign{$\m@th\displaystyle\hfil#\hfil$\cr>\crcr\sim\crcr} }}{\vbox{\offinterlineskip\halign{$\m@th\textstyle\hfil#\hfil$\cr>\crcr\sim\crcr} }}{\vbox{\offinterlineskip\halign{$\m@th\scriptstyle\hfil#\hfil$\cr>\crcr\sim\crcr} }}{\vbox{\offinterlineskip\halign{$\m@th\scriptscriptstyle\hfil#\hfil$\cr>\crcr\sim\crcr} }}}\pi, i.e., for ν <100(δRM106radm2)1/2GHz, <𝜈100superscript𝛿𝑅𝑀superscript106radsuperscriptm212GHz\nu\mathrel{\mathchoice{\vbox{\offinterlineskip\halign{$\m@th\displaystyle\hfil#\hfil$\cr<\crcr\sim\crcr} }}{\vbox{\offinterlineskip\halign{$\m@th\textstyle\hfil#\hfil$\cr<\crcr\sim\crcr} }}{\vbox{\offinterlineskip\halign{$\m@th\scriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr} }}{\vbox{\offinterlineskip\halign{$\m@th\scriptscriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr} }}}100\left({\delta RM\over 10^{6}\rm\ rad\ m^{-2}}\right)^{1/2}{\rm GHz}, (14) where $`\delta RM`$ is the difference in the rotation measure for photons of a given frequency which travel through different parts of the accretion flow.<sup>4</sup><sup>4</sup>4To be precise, $`\delta RM`$ is the difference in the rotation measure at $`r10010^4`$ on scales of $`r_\nu `$, the source size at frequency $`\nu `$. This is difficult to calculate analytically, but could be determined from future MHD simulations of non-radiating accretion flows. Quantitative calculations of depolarization by differential Faraday rotation are uncertain; two points are, however, clear: (1) At low frequencies, $`100`$ GHz, Sgr A is easily depolarized at r >102 >𝑟superscript102r\mathrel{\mathchoice{\vbox{\offinterlineskip\halign{$\m@th\displaystyle\hfil#\hfil$\cr>\crcr\sim\crcr} }}{\vbox{\offinterlineskip\halign{$\m@th\textstyle\hfil#\hfil$\cr>\crcr\sim\crcr} }}{\vbox{\offinterlineskip\halign{$\m@th\scriptstyle\hfil#\hfil$\cr>\crcr\sim\crcr} }}{\vbox{\offinterlineskip\halign{$\m@th\scriptscriptstyle\hfil#\hfil$\cr>\crcr\sim\crcr} }}}10^{2}. The required $`\delta RM`$ is $`10^6\mathrm{rad}\mathrm{m}^2`$, orders of magnitudes smaller than the values of $`RM`$ obtained at $`r10^210^4`$. (2) Emission above $`100`$ GHz can plausibly be linearly polarized if the accretion rate onto Sgr A is much less than the Bondi rate. In particular, equations (13) and (14) show that for $`\dot{M}\dot{M}_{\mathrm{Bondi}}`$, emission above $`100`$ GHz is not depolarized propagating out of the accretion flow. ## 4 Discussion ADAF/Bondi models assume that the accretion rate onto Sgr A is of order the Bondi rate ($`10^410^5M_{}\mathrm{yr}^1`$) and that the radio to infrared emission is produced by synchrotron emission from marginally relativistic electrons ($`T_e10^910^{10}`$ K). In such models the rotation measure is >1010radm2 >absentsuperscript1010radsuperscriptm2\mathrel{\mathchoice{\vbox{\offinterlineskip\halign{$\m@th\displaystyle\hfil#\hfil$\cr>\crcr\sim\crcr} }}{\vbox{\offinterlineskip\halign{$\m@th\textstyle\hfil#\hfil$\cr>\crcr\sim\crcr} }}{\vbox{\offinterlineskip\halign{$\m@th\scriptstyle\hfil#\hfil$\cr>\crcr\sim\crcr} }}{\vbox{\offinterlineskip\halign{$\m@th\scriptscriptstyle\hfil#\hfil$\cr>\crcr\sim\crcr} }}}10^{10}\rm\ rad\ m^{-2} inside $`100`$ Schwarzschild radii where the synchrotron emission is produced. ADAF/Bondi models thus predict that Sgr A should be depolarized by Faraday rotation over the entire radio to infrared spectrum, and should have nearly zero linear polarization. The theoretical arguments summarized in §1 propose that the accretion rate onto Sgr A is much less than the Bondi rate. We have described one such model, in which the electron heating rate is of order the rate of change of the magnetic energy density. For an accretion rate $`10^3`$ times smaller than the Bondi rate, i.e., $`10^8M_{}\mathrm{yr}^1`$, and with relativistic electrons with $`\gamma 100`$, this model can explain the observed high frequency emission from Sgr A. Moreover, it predicts that the rotation measure in the accretion flow is much smaller than in ADAF/Bondi models. This is because the gas density and magnetic field strength close to the black hole are much smaller, and because the electrons are relativistic ($`RM\gamma ^2\mathrm{log}\gamma `$ for $`\gamma 1`$; see §3 and the Appendix). The maximal contribution to the rotation measure comes from $`10^210^3`$ Schwarzschild radii, where $`RM10^6\mathrm{rad}\mathrm{m}^2`$. Rotation measures of $`10^6\mathrm{rad}\mathrm{m}^2`$ can depolarize Sgr A at $`\nu 100`$ GHz by differential Faraday rotation; photons of a given frequency travel through different rotation measures on their way out of the accretion flow. Following Bower et al. (1999ab), we believe that this accounts for the <0.2% <absentpercent0.2\mathrel{\mathchoice{\vbox{\offinterlineskip\halign{$\m@th\displaystyle\hfil#\hfil$\cr<\crcr\sim\crcr} }}{\vbox{\offinterlineskip\halign{$\m@th\textstyle\hfil#\hfil$\cr<\crcr\sim\crcr} }}{\vbox{\offinterlineskip\halign{$\m@th\scriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr} }}{\vbox{\offinterlineskip\halign{$\m@th\scriptscriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr} }}}0.2\% linear polarization of Sgr A at low frequencies (from $`4`$ to $`23`$ GHz; see Bower et al. 1999ab);<sup>5</sup><sup>5</sup>5Bower et al. (1999ab) showed that “bandwidth” depolarization at $`\nu 8`$ GHz requires $`RM>10^7\mathrm{rad}\mathrm{m}^2`$; this constraint does not, however, apply to the depolarization discussed here, namely that due to a spatially varying $`RM`$. it is less clear, however, that it can account for Bower et al.’s (1999b) limit of <1% <absentpercent1\mathrel{\mathchoice{\vbox{\offinterlineskip\halign{$\m@th\displaystyle\hfil#\hfil$\cr<\crcr\sim\crcr} }}{\vbox{\offinterlineskip\halign{$\m@th\textstyle\hfil#\hfil$\cr<\crcr\sim\crcr} }}{\vbox{\offinterlineskip\halign{$\m@th\scriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr} }}{\vbox{\offinterlineskip\halign{$\m@th\scriptscriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr} }}}1\% linear polarization at $`86`$ GHz (see below). In fact, rotation measures of $`10^6\mathrm{rad}\mathrm{m}^2`$ are insufficient to depolarize emission above $`100`$ GHz. As a result, in models with accretion rates much less than the Bondi rate, >100 >absent100\mathrel{\mathchoice{\vbox{\offinterlineskip\halign{$\m@th\displaystyle\hfil#\hfil$\cr>\crcr\sim\crcr} }}{\vbox{\offinterlineskip\halign{$\m@th\textstyle\hfil#\hfil$\cr>\crcr\sim\crcr} }}{\vbox{\offinterlineskip\halign{$\m@th\scriptstyle\hfil#\hfil$\cr>\crcr\sim\crcr} }}{\vbox{\offinterlineskip\halign{$\m@th\scriptscriptstyle\hfil#\hfil$\cr>\crcr\sim\crcr} }}}100 GHz emission is not depolarized propagating out of the accretion flow; intrinsically polarized synchrotron emission may therefore be observable at high frequencies. The above considerations show that the linear polarization of Sgr A at high frequencies provides a means of distinguishing between accretion at the Bondi rate, and accretion at a much smaller rate. In fact, Aitken et al. (2000) report a possible detection of $`10\%`$ linear polarization from Sgr A between $`150`$ and $`400`$ GHz. If confirmed, these observations require an accretion rate onto Sgr A\* much less than the Bondi rate, roughly $`\dot{M}10^8M_{}\mathrm{yr}^1`$. One difficulty in interpreting Aitken et al’s results is the large beam ($`20^{\prime \prime }`$) of the SCUBA camera on the JCMT. This large beam forced Aitken et al. to subtract out free-free and (polarized!) dust emission in order to isolate the flux and polarization of Sgr A. Future high resolution polarimetry at mm wavelengths is clearly necessary to further address this important issue. Aitken et al. find that the position angle of Sgr A changes by <10o <absentsuperscript10𝑜\mathrel{\mathchoice{\vbox{\offinterlineskip\halign{$\m@th\displaystyle\hfil#\hfil$\cr<\crcr\sim\crcr} }}{\vbox{\offinterlineskip\halign{$\m@th\textstyle\hfil#\hfil$\cr<\crcr\sim\crcr} }}{\vbox{\offinterlineskip\halign{$\m@th\scriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr} }}{\vbox{\offinterlineskip\halign{$\m@th\scriptscriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr} }}}10^{o} between $`\lambda =0.135`$ cm and $`\lambda =0.2`$ cm; at face value this implies RM <105radm2 <𝑅𝑀superscript105radsuperscriptm2RM\mathrel{\mathchoice{\vbox{\offinterlineskip\halign{$\m@th\displaystyle\hfil#\hfil$\cr<\crcr\sim\crcr} }}{\vbox{\offinterlineskip\halign{$\m@th\textstyle\hfil#\hfil$\cr<\crcr\sim\crcr} }}{\vbox{\offinterlineskip\halign{$\m@th\scriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr} }}{\vbox{\offinterlineskip\halign{$\m@th\scriptscriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr} }}}10^{5}\rm\ rad\ m^{-2}, somewhat smaller than the values of $`10^6\mathrm{rad}\mathrm{m}^2`$ in our model. This assumes, however, that the intrinsic position angle of Sgr A is the same at $`\lambda =0.135`$ cm and $`\lambda =0.2`$ cm, which need not be the case. Moreover, our estimates of $`RM`$ are actually upper limits, since they assume (1) equipartition magnetic fields aligned along the line of sight and (2) that our line of sight passes through the equatorial plane of the accretion flow. Our analysis of depolarization is applicable even if the radio emission from Sgr A is dominated by a jet/outflow, rather than the accretion flow as we have assumed. In jet models, it is still natural for the highest frequency emission to originate very close to the black hole; in Falcke’s model, e.g., the >100 >absent100\mathrel{\mathchoice{\vbox{\offinterlineskip\halign{$\m@th\displaystyle\hfil#\hfil$\cr>\crcr\sim\crcr} }}{\vbox{\offinterlineskip\halign{$\m@th\textstyle\hfil#\hfil$\cr>\crcr\sim\crcr} }}{\vbox{\offinterlineskip\halign{$\m@th\scriptstyle\hfil#\hfil$\cr>\crcr\sim\crcr} }}{\vbox{\offinterlineskip\halign{$\m@th\scriptscriptstyle\hfil#\hfil$\cr>\crcr\sim\crcr} }}}100 GHz emission arises from <10Rg <absent10subscript𝑅𝑔\mathrel{\mathchoice{\vbox{\offinterlineskip\halign{$\m@th\displaystyle\hfil#\hfil$\cr<\crcr\sim\crcr} }}{\vbox{\offinterlineskip\halign{$\m@th\textstyle\hfil#\hfil$\cr<\crcr\sim\crcr} }}{\vbox{\offinterlineskip\halign{$\m@th\scriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr} }}{\vbox{\offinterlineskip\halign{$\m@th\scriptscriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr} }}}10\ R_{g}, in what is really a “transition region” between the accretion flow and the jet (Falcke 1999). In order for this emission to not be depolarized (either in situ or propagating through the accretion flow), our constraints on the rotation measure and the plasma conditions close to the black hole still apply. Two scenarios in which accretion at the Bondi rate could be consistent with observed linear polarization at high frequencies are (1) if the high frequency emission arises close to the black hole, but in a nearly empty funnel pointed directly towards us (e.g., along the rotation axis of an ADAF) or (2) if the high frequency emission from Sgr A is produced at very large distances from the black hole, r >103 >𝑟superscript103r\mathrel{\mathchoice{\vbox{\offinterlineskip\halign{$\m@th\displaystyle\hfil#\hfil$\cr>\crcr\sim\crcr} }}{\vbox{\offinterlineskip\halign{$\m@th\textstyle\hfil#\hfil$\cr>\crcr\sim\crcr} }}{\vbox{\offinterlineskip\halign{$\m@th\scriptstyle\hfil#\hfil$\cr>\crcr\sim\crcr} }}{\vbox{\offinterlineskip\halign{$\m@th\scriptscriptstyle\hfil#\hfil$\cr>\crcr\sim\crcr} }}}10^{3}. The former possibility requires a rather special geometry<sup>6</sup><sup>6</sup>6For example, in Stone et al.’s (1999) simulations of non-radiating accretion flows, the density varies with polar angle roughly as $`\rho \mathrm{sin}^2\theta `$ (see also Quataert & Gruzinov 2000); thus for $`\dot{M}\dot{M}_{\mathrm{Bondi}}`$, the emission must be confined to θ <3o <𝜃superscript3𝑜\theta\mathrel{\mathchoice{\vbox{\offinterlineskip\halign{$\m@th\displaystyle\hfil#\hfil$\cr<\crcr\sim\crcr} }}{\vbox{\offinterlineskip\halign{$\m@th\textstyle\hfil#\hfil$\cr<\crcr\sim\crcr} }}{\vbox{\offinterlineskip\halign{$\m@th\scriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr} }}{\vbox{\offinterlineskip\halign{$\m@th\scriptscriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr} }}}{\rm 3}^{o} in order for the density to be sufficiently small that high frequency emission is not depolarized. In addition our line of sight must lie within <3o <absentsuperscript3𝑜\mathrel{\mathchoice{\vbox{\offinterlineskip\halign{$\m@th\displaystyle\hfil#\hfil$\cr<\crcr\sim\crcr} }}{\vbox{\offinterlineskip\halign{$\m@th\textstyle\hfil#\hfil$\cr<\crcr\sim\crcr} }}{\vbox{\offinterlineskip\halign{$\m@th\scriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr} }}{\vbox{\offinterlineskip\halign{$\m@th\scriptscriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr} }}}3^{o} of the rotation axis of the flow. and the latter is ruled out by the VLBI source size of $`10R_g`$ (Krichbaum et al. 1998) and the variability of Sgr A at $`100`$ GHz (Tsuboi, Miyazaki, & Tsutsumi 1999). ### 4.1 Application to Other Systems Although we have have focused our analysis on Sgr A at the Galactic Center, linear polarization of high frequency radio emission can be used as a probe of the accretion physics in other low-luminosity galactic nuclei (see, e.g., Nagar et al. 2000 for recent VLA observations of LLAGN). For a black hole of mass $`M=m_910^9M_{}`$ accreting (spherically) at a rate $`\dot{M}=10^4\dot{m}_4\dot{M}_{\mathrm{edd}}10^{23}\dot{m}_4m_9`$ g s<sup>-1</sup>, the density, magnetic field strength, and rotation measure in ADAF models are $$n3\times 10^6\dot{m}_4m_9^1r^{3/2}\mathrm{cm}^3,$$ (15) $$B100\dot{m}_4^{1/2}m_9^{1/2}r^{5/4}\mathrm{G},$$ (16) and $$RM3\times 10^{10}\dot{m}_4^{3/2}m_9^{1/2}r^{7/4}\mathrm{rad}\mathrm{m}^2.$$ (17) Equation (17) shows that large rotation measures and the associated depolarization of synchrotron emission by Faraday rotation are generic features of ADAF models (unless $`\dot{m}_41`$). The absence of observed linear polarization in the radio spectrum of a low-luminosity galactic nucleus would be consistent with ADAF models. By contrast, detected linear polarization would argue against an ADAF as the source of the observed radio emission. A particularly interesting class of systems for future polarimetry are elliptical galaxies in nearby X-ray clusters (e.g., NGC 4649, 4472, and 4636 in the Virgo cluster). As discussed by, e.g., Fabian & Canizares (1988), Fabian & Rees (1995), and Di Matteo et al. (1999, 2000), many of these galaxies have extremely dim nuclei given the inferred black hole masses ($`10^9M_{}`$) and Bondi accretion rates. Linear polarization may shed important light on the physics of these systems. ###### Acknowledgements. We thank Don Backer, Roger Blandford, and Mark Reid for useful correspondence, Bruce Draine for useful conversations, and John Bahcall, Heino Falcke, and Feryal Özel for helpful comments on the paper. EQ is supported by NASA through Chandra Fellowship PF9-10008, awarded by the Chandra X–ray Center, which is operated by the Smithsonian Astrophysical Observatory for NASA under contract NAS 8-39073. AG was supported by the W. M. Keck Foundation and NSF PHY-9513835. ## Appendix A Faraday rotation in an ultra-relativistic Maxwellian plasma Faraday rotation in a cold plasma is described by a change in position angle given by $$\frac{d\theta }{dl}=\frac{k_{}}{2}\frac{\omega _p^2\omega _B}{\omega ^3},$$ (A1) where $`\omega =ck`$ is the frequency of the radio wave, $`k_{}`$ is the projection of the wavenumber along the magnetic field, $`\omega _p^2=4\pi ne^2/m_e`$ is the plasma frequency, and $`\omega _B=eB/(m_ec)`$ is the cyclotron frequency. This corresponds to the usual rotation measure $$RM\frac{\theta }{\lambda ^2}=\frac{e^3}{2\pi m_e^2c^4}𝑑𝐥𝐁n=2.63\times 10^{13}\times 𝑑𝐥𝐁n\frac{\mathrm{rad}}{\mathrm{m}^2}.$$ (A2) Here we derive the rotation measure for an ultrarelativistic Maxwellian plasma: $$RM_\gamma =\frac{e^3}{2\pi m_e^2c^4}𝑑𝐥𝐁n\frac{\mathrm{log}\gamma }{2\gamma ^2}=2.63\times 10^{13}\times 𝑑𝐥𝐁n\frac{\mathrm{log}\gamma }{2\gamma ^2}\frac{\mathrm{rad}}{\mathrm{m}^2},$$ (A3) where we have defined $`\gamma kT_e/(m_ec^2)`$. The dominant correction to the non-relativistic expression is the relativistic mass: $`m_e\gamma m_e`$. We use the Vlasov equations to calculate the plasma permittivity and hence the dispersion relation for electromagnetic waves. For a magnetic field and wavenumber along the z axis, the first-order (in the unperturbed magnetic field) permittivity is given by $$ϵ_{xy}^{(1)}=\frac{i}{2\omega }\frac{4\pi e^2}{m_e}\frac{eB}{m_ec}d^3p\frac{1}{(\omega kv_z)^2}\frac{p_{}^2}{p}\frac{dF}{dp}\frac{m_e^2c^2}{p^2+m_e^2c^2},$$ (A4) where the unperturbed distribution function is normalized by $`d^3pF(p)=n`$, and $`p_{}^2p_x^2+p_y^2`$. For a cold plasma, equation (A4) gives $$ϵ_{xy}=\frac{i\omega _p^2\omega _B}{\omega ^3}.$$ (A5) Using standard arguments (e.g., Rybicki & Lightman 1979) this leads to the RM for a cold plasma given by equation (A2). For an ultra-relativistic plasma, equation (A4) gives $$ϵ_{xy}=\frac{i\omega _p^2\omega _B}{\omega ^3}\frac{\mathrm{log}\gamma }{2\gamma ^2},$$ (A6) where we have not changed the definition of $`\omega _p`$ and $`\omega _B`$ in the ultra-relativistic regime. Equation (A6) for the permittivity gives the ultra-relativistic RM in equation (A3).
warning/0004/math0004174.html	ar5iv	text	# Weyl modules for classical and quantum affine algebras ## 0. Introduction The study of the irreducible finite–dimensional representations of quantum affine algebras has been the subject of a number of papers, \[AK\], \[CP3\], \[CP5\], \[FR\], \[FM\], \[GV\], \[KS\] to name a few. However, the structure of these representations is still unknown except in certain special cases. In this paper, we approach the problem by studying the classical ($`q1`$) limits of these representations. Standard results imply, for example, that if $`V`$ is a finite–dimensional representation of $`𝐔_q(\widehat{𝔤})`$, its $`q1`$ limit $`\overline{V}`$ has the same structure as a $`𝔤`$–module as $`V`$ has as a $`U_q(𝔤)`$–module. We begin by studying an appropriate class of representations of the affine Lie algebra $`\widehat{𝔤}`$. The finite–dimensional irreducible representations of $`\widehat{𝔤}`$ were classified in \[C\], \[CP1\], where it was shown that such representations are highest weight in a suitable sense, the highest weight being an $`n`$–tuple of polynomials $`𝝅`$, where $`n`$ is the rank of $`𝔤`$. We therefore study the class of all highest weight finite–dimensional representations of $`\widehat{𝔤}`$. In fact, we prove that corresponding to each irreducible finite–dimensional representation $`V(𝝅)`$ there exists a unique (up to isomorphism) finite–dimensional highest weight module $`W(𝝅)`$, such that any finite dimensional highest weight module $`V`$ with highest weight $`𝝅`$ is a quotient of $`W(𝝅)`$. We call these modules the Weyl modules because of an analogy with the modular representation theory of $`𝔤`$, which we now explain. In \[CP5\], we showed that the irreducible representations of $`U_q(\widehat{𝔤})`$ are also highest weight and that their isomorphism classes are parametrized by a $`n`$–tuples of polynomials $`𝝅_q`$ with coefficients in $`𝐂(q)`$. Under a natural condition on $`𝝅_q`$, the corresponding representation $`V_q(𝝅_q)`$ of $`𝐔_q(\widehat{𝔤})`$ specializes as $`q1`$ to a representation $`\overline{V_q(𝝅_q)}`$ of $`\widehat{𝔤}`$ and is a quotient of $`W(𝝅)`$, where $`𝝅`$ is obtained from $`𝝅_q`$ by setting $`q=1`$. We conjecture that every $`W(𝝅)`$ is the classical limit of an irreducible $`U_q(\widehat{𝔤})`$-module. This is analogous to the fact that the Weyl modules for $`𝔤`$ in characteristic $`p`$ are the mod $`p`$ reductions of the irreducible modules in characteristic zero. We prove the conjecture in the case of $`𝔤=sl_2`$ in this paper. The conjecture is also true for the fundamental representations of $`𝐔_q(\widehat{𝔤})`$, but the details of that will appear elsewhere. In Section 3, we prove a factorization property of Weyl modules analogous the one for the irreducible modules proved in \[CP1\]. In Section 5We obtain a necessary and sufficient condition for the Weyl modules to be irreducible: the interesting feature of this proof is that it uses the fact that the specialized irreducible modules for the quantum algebra are quotients of the Weyl module. Further, the condition for irreducibility of the Weyl modules is the same as a condition that appears first in the work of Drinfeld on the closely related Yangians, \[Dr1\]. The Weyl modules we define are quotients of a family of level zero integrable modules for the extended affine Lie algebra, one corresponding to each dominant integral weight of $`𝔤`$. We call these modules $`W(\lambda )`$. According to unpublished work of Kashiwara, these modules are the classical analogues of the modules $`V^{max}(\lambda )`$ defined in \[K\]. Further, Kashiwara has a number of conjectures on the crystal basis of $`V^{max}(\lambda )`$. In Section 6, we identify the modules $`W(\lambda )`$ explicitly in the case of $`sl_2`$. A similar identification can then be proved for the modules $`V^{max}(\lambda )`$ which settles one of Kashiwara’s conjectures, but the details of this will appear elsewhere. Acknowledgements. We thank F.Knop, D.Rush and R.Sujatha for discussions. We also thank M.Kashiwara for explaining his conjectures to us and for drawing our attention to the modules $`V^{max}(\lambda )`$. ## 1. Preliminaries and Some Identities Let $`𝔤`$ be a finite-dimensional complex simple Lie algebra of rank $`n`$, $`𝔥`$ a Cartan subalgebra of $`𝔤`$ and $`R`$ the set of roots of $`𝔤`$ with respect to $`𝔥`$. Let $`I=\{1,2,\mathrm{},n\}`$, fix a set of simple roots $`\alpha _i`$ ($`iI`$), and let $`R^+`$ be the corresponding set of positive roots. Let $`\theta R^+`$ be the highest root in $`R^+`$. For $`\alpha R^+`$, fix non-zero elements $`x_\alpha ^\pm 𝔤`$, $`h_\alpha 𝔥`$ such that $$[x_\alpha ^+,x_\alpha ^{}]=h_\alpha ,[h_\alpha ,x_\alpha ^\pm ]=\pm 2x_\alpha ^\pm .$$ Let $`Q=_{i=1}^n𝐙\alpha _i`$ (resp. $`Q_+=_{i=1}^n𝐍\alpha _i`$) denote the root (resp. positive root) lattice of $`𝔤`$. For $`\eta Q^+`$, $`\eta =_ir_i\alpha _i`$, we set $`\text{ht}\eta =_ir_i`$. The lattice $`P`$ (resp. $`P_+`$) of integral (resp. dominant integral) weights is the set of elements $`\lambda 𝔥^{}`$ such that $`\lambda (h_\alpha )𝐙`$ for all $`\alpha R`$ (resp. $`\lambda (h_\alpha )0`$ for all $`\alpha R^+`$). For $`iI`$, the fundamental weight $`\omega _i`$ of $`𝔤`$ is given by $`\omega _i(\alpha _j)=\delta _{ij}`$. Let $`<,>`$ be the bilinear pairing on $`P`$ such that $`<\alpha _i,\omega _j>=\delta _{ij}`$. Set $`a_{ij}=<\alpha _i,\alpha _j>`$ for $`i,jI`$. The bilinear form induces an isomorphism $`𝔥𝔥^{}`$ such that, if $`\beta =_ir_i\alpha _iR^+`$, then $$h_\beta =\underset{j}{}\frac{d_jr_j}{d_\beta }h_j,$$ where for a root $`\alpha R`$, we set $`d_\alpha =\frac{1}{2}<\alpha ,\alpha >`$. Let $`W\text{Aut}(𝔥^{})`$ be the Weyl group of $`𝔤`$; it is well known that $`W`$ is generated by simple reflections $`s_i`$ ($`iI`$). The extended loop algebra of $`𝔤`$ is the Lie algebra $$L^e(𝔤)=𝔤𝐂[t,t^1]𝐂d,$$ with commutator given by $$[d,xt^r]=rxt^r,[xt^r,yt^s]=[x,y]t^{r+s}$$ for $`x,y𝔤`$, $`r,s𝐙`$. The loop algebra $`L(𝔤)`$ is the subalgebra $`𝔤𝐂[t,t^1]`$ of $`L^e(𝔤)`$. Let $`𝔥^e=𝔥𝐂d`$. Define $`\delta (𝔥^e)^{}`$ by $$\delta (𝔥)=0,\delta (d)=1.$$ Extend $`\lambda 𝔥^{}`$ to an element of $`(𝔥^e)^{}`$ by setting $`\lambda (d)=0`$. Set $`P^e=_{i=1}^n𝐙\omega _i𝐙\delta `$, and define $`P_+^e`$ in the obvious way. We regard $`W`$ as acting on $`(𝔥^e)^{}`$ by setting $`w(\delta )=\delta `$ for all $`wW`$. For any $`x𝔤`$, $`m𝐙`$, we denote by $`x_m`$ the element $`xt^mL^e(𝔤)`$. Set $`e_i^\pm =x_{\alpha _i}^\pm 1`$ and $`e_0^\pm =x_\theta ^{}t^{\pm 1}`$. Then, the elements $`e_i^\pm `$ ($`i=0,\mathrm{},n`$) and $`d`$ generate $`L^e(𝔤)`$. For any Lie algebra $`𝔞`$, the universal enveloping algebra of $`𝔞`$ is denoted by $`𝐔(𝔞)`$. We set $$𝐔(L^e(𝔤))=𝐔^e,𝐔(L(𝔤))=𝐔,𝐔(𝔤)=𝐔^{fin}.$$ Let $`𝐔(<)`$ (resp. $`𝐔(>)`$) be the subalgebra of $`𝐔`$ generated by the $`x_{\alpha _i,m}^{}`$, (resp. $`x_{\alpha _i,m}^+`$) for $`iI`$, $`m𝐙`$. Clearly, $`x_{\alpha ,m}^{}𝐔(<)`$ (resp. $`x_{\alpha ,m}^+𝐔(>)`$) for all $`\alpha R^+`$, $`m𝐙`$. Set $`𝐔^{fin}(<)=𝐔(<)𝐔^{fin}`$ and define $`𝐔^{fin}(>)`$ similarly. Finally, let $`𝐔(0)`$ be the subalgebra of $`𝐔`$ generated by $`h_{\alpha ,m}`$ for $`\alpha R`$, $`m𝐙`$, $`m0`$. We have $`𝐔^{fin}`$ $`=𝐔^{fin}(<)𝐔(𝔥)𝐔^{fin}(>),`$ $`𝐔^e`$ $`=𝐔(<)𝐔(0)𝐔(𝔥^e)𝐔(>).`$ ###### Lemma 1.1. The assignment $`T(x_{\alpha _i,m}^\pm )=x_{\alpha _i,m\pm 1}^\pm `$, for $`iI`$, $`m𝐙`$, defines an algebra automorphism of $`𝐔`$. ∎ We next recall some identities in $`𝐔^e`$, which are most conveniently stated using generating series. Thus, for any $`\beta R^+`$, we introduce the following power series in an indeterminate $`u`$: $`\stackrel{~}{X}_\beta ^{}(u)={\displaystyle \underset{m=\mathrm{}}{\overset{\mathrm{}}{}}}x_{\beta ,m}^{}u^{m+1},`$ $`X_\beta ^{}(u)={\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}x_{\beta ,m}^{}u^m,`$ $`X_\beta ^+(u)={\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}x_{\beta ,m}^+u^m,`$ $`X_{\beta ,0}^{}(u)={\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}x_{\beta ,m}^{}u^{m+1},`$ $`\stackrel{~}{H}_\beta (u)={\displaystyle \underset{m=\mathrm{}}{\overset{\mathrm{}}{}}}h_{\beta ,m}u^{m+1},`$ $`\mathrm{\Lambda }_\beta ^\pm (u)={\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}\mathrm{\Lambda }_{\beta ,\pm m}u^m=\text{exp}\left({\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{h_{\beta ,\pm k}}{k}}u^k\right).`$ Set $`x_{\alpha _i}^\pm =x_i^\pm `$, $`x_{\alpha _i,m}^\pm =x_{i,m}^\pm `$, and define $`h_i`$, $`\mathrm{\Lambda }_i^\pm `$, etc., similarly. The next lemma follows easily from the definition of the $`\mathrm{\Lambda }_{i,m}`$ ($`=\mathrm{\Lambda }_{\alpha _i,m}`$). ###### Lemma 1.2. The subalgebra $`𝐔(0)`$ of $`𝐔`$ is generated by the elements $`\mathrm{\Lambda }_{i,m}`$, for $`iI`$, $`m𝐙`$. ∎ For any power series $`f`$ in $`u`$ with coefficients in an algebra $`A`$, let $`f_m`$ be the coefficient of $`u^m`$ ($`m𝐙`$) and let $`f^{}`$ denote the derivative of $`f`$ with respect to $`u`$. For $`x𝐔`$, $`r𝐙^+`$, set $$x^{(r)}=\frac{x^r}{r!}.$$ For an algebra $`A`$, let $`A_+`$ denote the augmentation ideal. The next result is a reformulation of a result of Garland, \[G\]. ###### Lemma 1.3. Let $`sr1`$, $`\beta R^+`$. 1. $$(x_{\beta ,0}^+)^{(r)}(x_{\beta ,1}^{})^{(s)}=(1)^r(X_\beta ^{}(u)^{(sr)}\mathrm{\Lambda }_\beta ^+(u))_smod\mathrm{𝐔𝐔}(>)_+.$$ 2. $$(x_{\beta ,1}^+)^{(r)}(x_{\beta ,0}^{})^{(s)}=(1)^r(X_{\beta ,0}^{}(u)^{(sr)}\mathrm{\Lambda }_\beta ^+(u))_smod\mathrm{𝐔𝐔}(>)_+.$$ ###### Proof. We use the identity in \[G, Lemma 7.1\], in the form given in \[CP6, Lemma 5.1\] for its quantum version, namely $`(x_{\beta ,0}^+)^{(r)}(x_{\beta ,1}^{})^{(s)}`$ $`={\displaystyle \underset{t=0}{\overset{r}{}}}{\displaystyle \underset{m=0}{\overset{t}{}}}{\displaystyle \underset{k=0}{\overset{m}{}}}(1)^t\left(u^{s+t}X_\beta ^{}(u)^{(st)}\right)_{tm}\mathrm{\Lambda }_{\beta ,k}\left(X_\beta ^+(u)^{(rt)}\right)_{mk}.`$ In the sum on the right-hand side of the equality, we get an element of $`\mathrm{𝐔𝐔}(>)_+`$ unless $`t=r`$ and $`m=k`$, and so we have $`(x_{\beta ,0}^+)^{(r)}(x_{\beta ,1}^{})^{(s)}`$ $`=(1)^r{\displaystyle \underset{m=0}{\overset{r}{}}}\left(u^{s+r}X_\beta ^{}(u)^{(sr)}\right)_{rm}\mathrm{\Lambda }_{\beta ,m}mod\mathrm{𝐔𝐔}(>)_+`$ $`={\displaystyle \underset{r=0}{\overset{m}{}}}(1)^r\left(X_\beta ^{}(u)^{(sr)}\right)_{sm}\mathrm{\Lambda }_\beta ^+(u)_mmod\mathrm{𝐔𝐔}(>)_+`$ $`=(1)^r\left(X_\beta ^{}(u)^{(sr)}\mathrm{\Lambda }_\beta ^+(u)\right)_smod\mathrm{𝐔𝐔}(>)_+.`$ The identity in (ii) follows from (i) by applying the automorphism $`T:𝐔𝐔`$. ∎ We conclude this section with some elementary properties of integrable $`L^e(𝔤)`$-modules. A representation $`V`$ of $`𝐔^e`$ is called integrable if the Chevalley generators $`e_i^\pm `$, for $`i=0,1,\mathrm{},n`$, act locally nilpotently on $`V`$ and $$V=\underset{\lambda (𝔥^e)^{}}{}V_\lambda ,$$ where $$V_\lambda =\{vV:h.v=\lambda (h)vh𝔥^e\}.$$ It is well known that this implies that the elements $`x_{\beta ,m}^\pm `$ act locally nilpotently on $`V`$ for all $`\beta R^+`$ and $`m𝐙`$. Set $$V_\lambda ^+=\{vV_\lambda :x_{\beta ,m}^+.v=0\beta R^+,m𝐙\}.$$ It is easy to see that, if $`V`$ is integrable, then $`V_\lambda 0`$ (resp. $`V_\lambda ^+0`$) only if $`\lambda P^e`$ (resp. $`\lambda P_+^e`$). Further, if $`vV_\lambda ^+`$, then (1.1) $`(x_{\beta ,m}^{})^{\lambda (h_\beta )+1}.v`$ $`=0,`$ (1.2) $`V_\lambda 0V_{w\lambda }0`$ $`wW.`$ If $`\lambda P_+`$, let $`V^{fin}(\lambda )`$ be the finite-dimensional irreducible $`𝐔^{fin}`$-module with highest weight $`\lambda `$. If $`\lambda P_+^e`$, the restriction of $`\lambda `$ to $`𝔥^{}`$, also denoted by $`\lambda `$, is in $`P_+`$. If $`V`$ is an integrable $`𝐔^e`$-module and $`0vV_\lambda ^+`$, then $`𝐔^{fin}.v`$ is a $`𝐔^{fin}`$-submodule of $`V`$ isomorphic to $`V^{fin}(\lambda )`$. ###### Proposition 1.1. Let $`V`$ be an integrable $`𝐔^e`$-module. Let $`\lambda P_+^e`$, let $`0vV_\lambda ^+`$ and let $`\beta R^+`$. Then: 1. $`\mathrm{\Lambda }_{\beta ,m}.v=0`$ for $`m>\lambda (h_\beta )`$; 2. for $`r1`$, $`s>\lambda (h_\beta )`$, $$\left(X_\beta ^{}(u)^r\mathrm{\Lambda }_\beta ^+(u)\right)_s.v=0,\left(X_{\beta ,0}^{}(u)^r\mathrm{\Lambda }_\beta ^+(u)\right)_s.v=0;$$ 3. for all $`s𝐙`$, $$\left(\stackrel{~}{X}_\beta ^{}(u)\mathrm{\Lambda }_\beta ^+(u)\right)_s.v=0,\left(\stackrel{~}{H}_\beta (u)\mathrm{\Lambda }_\beta ^+(u)\right)_s.v=0;$$ 4. $`\mathrm{\Lambda }_{\beta ,m}.v=0`$ for all $`m>\lambda (h_\beta )`$; 5. for $`0m\lambda (h_\beta )`$, $$\mathrm{\Lambda }_{\beta ,\lambda (h_\beta )}\mathrm{\Lambda }_{\beta ,m}.v=\mathrm{\Lambda }_{\beta ,\lambda (h_\beta )m}.v.$$ ###### Proof. (i) This follows by taking $`r=s>\lambda (h_\beta )`$ in Lemma 1.3(i) and using equation (1.1). (ii) This follows from Lemma 1.3 by replacing $`r`$ by $`sr`$ and using equation (1.1). (iii) Taking $`r=\lambda (h_\beta )`$, $`s=\lambda (h_\beta )+1`$ in Lemma 1.3(i) gives (1.3) $$\underset{m=0}{\overset{\lambda (h_\beta )}{}}x_{\beta ,m+1}^{}\mathrm{\Lambda }_{\beta ,\lambda (h_\beta )m}.v=0.$$ Applying $`h_{\beta ,k}`$, for any $`k𝐙`$, to the above equation and noting that $`h_{\beta ,k}.vV_\lambda ^+`$, we get $$\underset{m=0}{\overset{\lambda (h_\beta )}{}}x_{\beta ,k+m+1}^{}\mathrm{\Lambda }_{\beta ,\lambda (h_\beta )m}.v=0,$$ which can be written as $$\left(\stackrel{~}{X}_\beta ^{}(u)\mathrm{\Lambda }_\beta ^+(u)\right)_{k+1}.v=0.$$ Applying $`x_{\beta ,s1}^+`$, for $`s𝐙`$, to both sides of equation (1.3) gives $$\underset{m=0}{\overset{\lambda (\beta )}{}}h_{\beta ,ms}\mathrm{\Lambda }_{\beta ,\lambda (h_\beta )m}.v=0,$$ i.e., $$\underset{m=s}{\overset{\lambda (h_\beta )s}{}}h_{\beta ,m}\mathrm{\Lambda }_{\beta ,\lambda (h_\beta )sm}.v=0.$$ Replacing $`s`$ by $`\lambda (h_\beta )s+1`$ and using part (i) of the lemma, one sees that this identity is equivalent to the second identity in (iii). (iv) and (v). During the remainder of this proof, write $`\mathrm{\Lambda }_\beta =\mathrm{\Lambda }_\beta ^+`$, $`\stackrel{~}{\mathrm{\Lambda }}_\beta (u)=\mathrm{\Lambda }_\beta ^{}(u^1)`$, so that $$\stackrel{~}{\mathrm{\Lambda }}_\beta (u)=\text{exp}\left(\underset{k=1}{\overset{\mathrm{}}{}}\frac{h_{\beta ,k}}{k}u^k\right).$$ Note that, as operators on $`V_\lambda ^+`$, $$\left(\lambda (h_\beta )u\frac{\mathrm{\Lambda }_\beta ^{}}{\mathrm{\Lambda }_\beta }+u\frac{\stackrel{~}{\mathrm{\Lambda }}_\beta ^{}}{\stackrel{~}{\mathrm{\Lambda }}_\beta }\right)=\stackrel{~}{H}_\beta .$$ By (iii), we have $$\mathrm{\Lambda }_\beta (u)\left(\lambda (h_\beta )u\frac{\mathrm{\Lambda }_\beta ^{}}{\mathrm{\Lambda }_\beta }+u\frac{\stackrel{~}{\mathrm{\Lambda }}_\beta ^{}}{\stackrel{~}{\mathrm{\Lambda }}_\beta }\right)=\stackrel{~}{H}_\beta \mathrm{\Lambda }_\beta =0,$$ as operators on $`V_\lambda ^+`$, so, since $`\mathrm{\Lambda }_\beta (u)`$ is invertible, $$\lambda (h_\beta )\mathrm{\Lambda }_\beta (u)\stackrel{~}{\mathrm{\Lambda }}_\beta (u)=u(\mathrm{\Lambda }_\beta ^{}\stackrel{~}{\mathrm{\Lambda }}_\beta \mathrm{\Lambda }_\beta \stackrel{~}{\mathrm{\Lambda }}_\beta ^{}).$$ Note that both sides of this equation make sense as power series in $`u`$ since $`\mathrm{\Lambda }_\beta (u)`$ is already known by (i) to involve only finitely many positive powers of $`u`$. Hence, as series with only finitely many positive powers (but possibly infinitely many negative powers), we have $$\left(\frac{\mathrm{\Lambda }_\beta }{\stackrel{~}{\mathrm{\Lambda }}_\beta }\right)^{}=\lambda (h_\beta )u^1\left(\frac{\mathrm{\Lambda }_\beta }{\stackrel{~}{\mathrm{\Lambda }}_\beta }\right),$$ and so $$\frac{\mathrm{\Lambda }_\beta (u)}{\stackrel{~}{\mathrm{\Lambda }}_\beta (u)}=A_\beta u^{\lambda (h_\beta )},$$ where $`A_\beta `$ is an operator on $`V_\lambda ^+`$ independent of $`u`$. Equating coefficients of $`u^{\lambda (h_\beta )}`$ shows that $`A_\beta =\mathrm{\Lambda }_{\beta ,\lambda (h_\beta )}`$ and then the equation (of operators on $`V_\lambda ^+`$) $$\mathrm{\Lambda }_{\beta ,\lambda (h_\beta )}(u)\stackrel{~}{\mathrm{\Lambda }}_\beta (u)=u^{\lambda (h_\beta )}\mathrm{\Lambda }_\beta (u),$$ proves both (iv) and (v). ∎ Fix a total order $``$ on $`R^+`$. ###### Proposition 1.2. Let $`V`$ be an integrable $`𝐔^e`$-module, let $`\lambda P_+^e`$ and let $`0vV_\lambda ^+`$ be such that $`V=𝐔^e.v`$. 1. If $`V_\mu 0`$, then $`\mu =\lambda \eta +r\delta `$ for some $`\eta Q_+`$, $`r𝐙`$ such that $`V^{fin}(\lambda )_{\lambda \eta }0`$. 2. $`V`$ is spanned by the elements $$x_{\beta _1,r_1}^{}x_{\beta _2,r_2}^{}\mathrm{}x_{\beta _s,r_s}^{}𝐔(0).v,$$ for $`0r_t<\lambda (h_{\beta _t})`$, $`0ts`$, $`\beta _1\beta _2\mathrm{}\beta _sR^+`$. ###### Proof. Since $`𝐔(>)_+.v=0`$, it follows that $`V=𝐔(<)𝐔(0).v`$ and hence $$V_\mu 0\mu =\lambda \eta +r\delta ,$$ for some $`\eta Q^+`$ and $`r𝐙`$. Choose $`\sigma W`$ such that $`\sigma (\lambda \eta )P_+`$. Since $`V`$ is integrable, $`V_{\sigma (\mu )}0`$, hence $`\sigma (\lambda \eta )=\lambda \eta ^{}`$ for some $`\eta ^{}Q_+`$. This implies that $`V^{fin}(\lambda )_{\sigma (\lambda \eta )}0`$, and hence that $`V^{fin}(\lambda )_{\lambda \eta }0`$. To prove (ii), note first that it is clear that elements of the form $$x_{\beta _1,r_1}^{}x_{\beta _2,r_2}^{}\mathrm{}x_{\beta _s,r_s}^{}𝐔(0).v,(0ts,\beta _tR^+,r_t𝐙)$$ span $`V`$. We prove by induction on $`s`$ that any such element is in the span of the elements $$x_{\gamma _1,k_1}^{}x_{\gamma _2,k_2}^{}\mathrm{}x_{\gamma _m,k_m}^{}𝐔(0).v,(0k_t<\lambda (h_{\gamma _t}),\gamma _1\gamma _2\mathrm{}\gamma _m,0tm).$$ For $`s=1`$, $`r_1\lambda (h_{\beta _1})`$, we have, by Proposition 1.1(iii), that $$\underset{r=0}{\overset{r_1}{}}x_{\beta _1,r}^{}\mathrm{\Lambda }_{\beta _1,r_1r}𝐔(0).v=0.$$ Since $`\mathrm{\Lambda }_{\beta _1,0}=1`$, this implies that $`x_{\beta _1,r_1}^{}.v`$ is in the span of the elements $`x_{\beta _1,r}^{}𝐔(0).v`$ with $`0r<r_1`$, from which the assertion follows. If $`r_1<0`$, we use $$\underset{r=0}{\overset{\lambda (h_{\beta _1})}{}}x_{\beta _1,r+r_1}^{}\mathrm{\Lambda }_{\beta _1,\lambda (h_{\beta _1})r}.v=0,$$ which follows from parts (i) and (ii) of Proposition 1.1. Since, by Proposition 1.1(v), $`\mathrm{\Lambda }_{\beta ,\lambda (h_\beta )}`$ is invertible, it follows that $`x_{\beta _1,r_1}^{}.v`$ is in the span of the elements $`x_{\beta _1,r}^{}.v`$ ($`0r>r_1`$). An obvious induction now gives the result. Suppose that we know the result for some $`s1`$. Let $`\beta _tR^+`$ and $`r_t𝐙`$ for $`0ts`$ be such that $`0r_t<\lambda (h_{\beta _t})`$ for all $`t`$. We have $`x_{\beta _0,r_0}^{}x_{\beta _1,r_1}^{}x_{\beta _2,r_2}^{}\mathrm{}x_{\beta _s,r_s}^{}.v`$ $`=x_{\beta _1,r_1}^{}x_{\beta _0,r_0}^{}x_{\beta _2,r_2}^{}\mathrm{}x_{\beta _s,r_s}^{}.v`$ $`+[x_{\beta _0,r_0}^{},x_{\beta _1,r_1}^{}]x_{\beta _2,r_2}^{}\mathrm{}x_{\beta _s,r_s}^{}.v.`$ Since $`[x_{\beta _0,r_0}^{},x_{\beta _1,r_1}^{}]𝔤t^{r_0+r_1}`$, the induction hypothesis applies to the second term on the right-hand side of the equality. As for the first term, notice that the induction hypothesis applied to $`x_{\beta _0,r_0}^{}x_{\beta _2,r_2}^{}\mathrm{}x_{\beta _s,r_s}^{}.v`$ implies that this is in the span of the elements obtained by applying ordered monomials in the $`x_{\gamma ,k}^{}`$ to $`v`$, for $`\gamma R^+`$ and $`0k<\lambda (h_\gamma )`$. Thus, we must show that every element of the form $$x_{\beta _1,r_1}^{}x_{\gamma _1,k_1}^{}x_{\gamma _2,k_2}^{}\mathrm{}x_{\gamma _m,k_m}^{}𝐔(0).v,$$ where $`\gamma _1\gamma _2\mathrm{}\gamma _m`$ and $`0k_t<\lambda (h_{\gamma _t})`$ for $`1tm`$, can be rewritten in the desired form. If $`\beta _1\gamma _1`$, there is nothing to prove. Otherwise, we have $$[x_{\beta _1,r_1}^{},x_{\gamma _1,k_1}^{}]𝔤_{\beta _1+\gamma _1}t^{r_1+k_1}$$ and $`\lambda (h_{\beta _1+\gamma _1})r_1+k_1`$. Using the induction hypothesis gives the result. ∎ ## 2. Maximal Integrable and Maximal Finite-Dimensional Modules In this section we define, for every $`\lambda P_+^e`$ an integrable $`𝐔^e`$-module $`W(\lambda )`$. Further, for any $`n`$-tuple $`𝝅=(\pi _1(u),\pi _2(u),\mathrm{},\pi _n(u))`$ of polynomials $`\pi _i(u)`$ in an indeterminate $`u`$ with constant term 1 and degree $`\lambda (h_i)`$, we define a finite-dimensional quotient $`𝐔`$-module $`W(𝝅)`$ of $`W(\lambda )`$. For $`\lambda P_+^e`$, let $`I_\lambda `$ be the left ideal in the subalgebra $`𝐔(<)𝐔(0)𝐔(𝔥)`$ of $`𝐔^e`$ generated by the following elements: $`h\lambda (h)(h𝔥),\mathrm{\Lambda }_{i,m}(iI,\|m\|>\lambda (h_i)),`$ $`\mathrm{\Lambda }_{i,\lambda (h_i)}\mathrm{\Lambda }_{i,m}\mathrm{\Lambda }_{i,\lambda (h_i)m}(iI,1m\lambda (h_i)),`$ $`\left(\stackrel{~}{X}_i^{}(u)\mathrm{\Lambda }_i^+(u)\right)_m𝐔(0)(iI,m𝐙),`$ $`\left(X_{i,0}^{}(u)^r\mathrm{\Lambda }_i^+(u)\right)_m𝐔(0)(iI,r1,\|m\|>\lambda (h_i)).`$ Let $`\stackrel{~}{I}_\lambda `$ be the left ideal in $`𝐔`$ generated by $`I_\lambda `$ and the $`x_{i,m}^+`$ for all $`iI`$, $`m𝐙`$, and let $`\stackrel{~}{I}_\lambda ^e`$ be the left ideal in $`𝐔^e`$ generated by $`\stackrel{~}{I}_\lambda `$ and $`d\lambda (d)`$. Set $$W(\lambda )=𝐔^e/\stackrel{~}{I}_\lambda ^e=𝐔/\stackrel{~}{I}_\lambda .$$ Clearly, $`W(\lambda )`$ is a left $`𝐔^e`$-module (and a left $`𝐔`$–module) through left multiplication. Let $`w_\lambda `$ be the image of $`1`$ in $`W(\lambda )`$. Then, $$𝐔(>)_+.w_\lambda =0,W(\lambda )=𝐔^e.w_\lambda =𝐔.w_\lambda .$$ Since $`\stackrel{~}{I}_\lambda 𝐔(0)\stackrel{~}{I}_\lambda `$, we can and do regard $`W(\lambda )`$ as a right $`𝐔(0)`$–module as well. For $`\eta Q^+`$, we set $$W(\lambda )[\eta ]=\underset{r𝐙}{}W(\lambda )_{\lambda \eta +r\delta }.$$ Clearly, $`W(\lambda )[\eta ]`$ is a right $`𝐔(0)`$–module for all $`\eta Q^+`$ and we have $$W(\lambda )=\underset{\eta Q^+}{}W(\lambda )[\eta ]$$ as right $`𝐔(0)`$–modules. Let $`I_\lambda (0)=I_\lambda 𝐔(0)`$. By the PBW theorem, it is easy to see that $$𝐔(0)/I_\lambda (0)𝐂[\mathrm{\Lambda }_{i,m},\mathrm{\Lambda }_{i,\lambda (h_i)}^1:iI,1m\lambda (h_i)].$$ In particular, $$W(\lambda )[0]𝐂[\mathrm{\Lambda }_{i,m},\mathrm{\Lambda }_{i,\lambda (h_i)}^1:iI,1m\lambda (h_i)]$$ as right $`𝐔(0)`$–modules. It follows immediately from Proposition 1.2(ii) that, for all $`\eta Q^+`$, $`W(\lambda )[\eta ]`$ is a finitely-generated $`𝐔(0)`$–module. Next, let $`𝝅=(\pi _1,\pi _2,\mathrm{},\pi _n)`$ be an $`n`$-tuple of polynomials in an indeterminate $`u`$ with constant term $`1`$, and define, for $`a𝐂`$, an element $`\lambda _{𝝅,a}(𝔥^e)^{}`$ by setting $`\lambda _{𝝅,a}(h_i)=\text{deg}\pi _i`$ ($`iI`$) and $`\lambda _{𝝅,a}(d)=a`$. Set (2.1) $$\pi _i^+=\pi _i,\pi _i^{}(u)=\frac{u^{\text{deg}\pi _i}\pi _i(u^1)}{u^{\text{deg}\pi _i}\pi _i(u^1)\|_{u=0}}.$$ Let $`I_𝝅(0)`$ be the maximal ideal in $`𝐔(0)`$ generated by $$(\mathrm{\Lambda }_i^\pm (u)\pi _i^\pm (u))_s(iI,s0),$$ and let $`𝐂_𝝅=𝐔(0)/I_𝝅(0)`$ be the one–dimensional $`𝐔(0)`$-module. Set $$W(𝝅)=W(\lambda _{𝝅,a})_{𝐔(0)}𝐂_𝝅.$$ Then, $`W(𝝅)`$ is a left $`𝐔`$-module (but not a $`𝐔^e`$-module) with $`x𝐔`$ acting as $`x1`$. Let $`w_𝝅`$ be the image of $`1`$ in $`W(𝝅)`$. Note that $`\mathrm{\Lambda }_i^\pm (u).w_𝝅=\pi _i^\pm (u)w_𝝅`$ ($`iI`$). The assignment $`w_{\lambda _{𝝅,a}}w_𝝅`$ extends to a surjective $`𝐔`$-module homomorphism $`W(\lambda _{𝝅,a})W(𝝅)`$. Recall that the affine Lie algebra $`\widehat{𝔤}`$ is an extension of $`L^e(𝔤)`$ by a 1-dimensional central subalgebra $`𝐂c`$. Any representation of $`L^e(𝔤)`$ is a representation of $`\widehat{𝔤}`$ by making $`c`$ act as zero. Set $`h_0=[e_0^+,e_0^{}]`$, the bracket being evaluated in $`\widehat{𝔤}`$. Then, $`h_0=ch_\theta `$. The following result is proved in \[CP2\]. ###### Lemma 2.1. Let $`V`$ be a $`\widehat{𝔤}`$-module generated by an element $`vV_\lambda `$ ($`\lambda P_+^e`$) such that $`(e_i^+)^{\lambda (h_i)+1}.v=0`$ $`\text{if}\lambda (h_i)0,`$ $`(e_i^{})^{\lambda (h_i)+1}.v=0`$ $`\text{if}\lambda (h_i)0,`$ for $`iI`$. Then, $`V`$ is an integrable $`\widehat{𝔤}`$-module. ∎ ###### Theorem 1. 1. Let $`\lambda P_+^e`$. Then, $`W(\lambda )`$ is an integrable $`𝐔^e`$-module. 2. Let $`𝝅`$ be an $`n`$-tuple of polynomials with constant term one. Then, $`W(𝝅)`$ is a finite-dimensional $`𝐔`$-module. ###### Proof. To prove (i), note that by Lemma 2.1 it suffices to show that $`e_i^+.w_\lambda =0,`$ $`(e_i^{})^{\lambda (h_i)+1}.w_\lambda =0(iI),`$ $`e_0^{}.w_\lambda =0,`$ $`(e_0^+)^{\lambda (h_\theta )+1}.w_\lambda =0.`$ Suppose that $`iI`$. Then, $`e_i^\pm =x_{i,0}^\pm `$ and it follows from the definition of $`W(\lambda )`$ that $`e_i^+.w_\lambda =0`$. By Lemma 1.3(ii), $$(x_{i,0}^{})^{\lambda (h_i)+1}.w_\lambda =\left(X_{i,0}^{}(u)^{\lambda (h_i)+1}\mathrm{\Lambda }_i^+(u)\right)_{\lambda (h_i)+1}.w_\lambda =0.$$ In particular, this proves that $`𝐔^{fin}.w_\lambda `$ is a finite-dimensional $`𝔤`$-module. Turning to the case $`i=0`$, notice that $`e_0^{}=x_{\theta ,1}^+`$ is a linear combination of products of the $`x_{i,m}^+`$ ($`iI`$, $`m𝐙`$). Hence, $`x_{\theta ,1}^+.w_\lambda =0`$. For any $`m0`$, let $`w_m=(e_0^+)^{m+1}.w_\lambda `$. Suppose that $`w_m0`$. Since $`[e_0^+,x_{\beta ,0}^{}]=0`$, it follows that $$𝐔^{fin}(<).w_m=e_0^{m+1}𝐔^{fin}(<).w_\lambda $$ is finite-dimensional. Since $`W(\lambda )[\eta ]=0`$ for all but finitely many $`\eta Q^+`$, it follows that $$W_m=𝐔^{fin}.w_m=𝐔^{fin}(>)𝐔(𝔥)𝐔^{fin}(<).w_m$$ is a finite-dimensional $`𝔤`$-module. Hence, for all $`\sigma W`$ (the Weyl group of $`𝔤`$), we have $$(W_m)_{\sigma (\lambda (m+1)\theta +(m+1)\delta )}0.$$ Choosing $`\sigma `$ so that $`\sigma (\theta )=\alpha _i`$ for some $`iI`$, we get $$W(\lambda )_{\sigma (\lambda )+(m+1)\alpha _i+(m+1)\delta }0.$$ But this can only happen for finitely many values of $`m`$. This proves that $`w_m=0`$ for all but finitely many $`m`$. The Lie subalgebra of $`L(𝔤)`$ generated by $`e_0^\pm `$ and $`h_\theta `$ is isomorphic to $`sl_2`$, and we have just shown that the corresponding $`sl_2`$-submodule generated by $`w_\lambda `$ is finite-dimensional. It follows from standard results that $`(e_0^{})^{\lambda (h_\theta )+1}.w_\lambda =0`$. To prove (ii), it suffices now to notice (using Proposition 1.2(ii)) that $`W(𝝅)`$ is spanned by the elements $$x_{\beta _1,r_1}^{}x_{\beta _2,r_2}^{}\mathrm{}x_{\beta _s,r_s}^{}1$$ for $`s0`$, $`0r_t<\lambda (h_{\beta _t})`$, $`\beta _tR^+`$. ∎ The modules $`W(\lambda )`$ and $`W(𝝅)`$ have certain universal properties. ###### Proposition 2.1. 1. Let $`V`$ be any integrable $`𝐔^e`$-module generated by a non-zero element $`vV_\lambda ^+`$. Then, $`V`$ is a quotient of $`W(\lambda )`$. 2. Let $`V`$ be a finite-dimensional quotient $`𝐔`$-module of $`W(\lambda )`$, and assume that $`\text{dim}V_\lambda =1`$. Then, $`V`$ is a quotient of $`W(𝝅)`$ for some choice of $`𝛑`$. 3. Let $`V`$ be finite-dimensional $`𝐔`$-module generated by a vector $`vV_\lambda ^+`$ and such that $`\text{dim}V_\lambda =1`$. Then, $`V`$ is a quotient of $`W(𝝅)`$ for some $`𝛑`$. ###### Proof. Part (i) is immediate from Proposition 1.1 and the definition of $`W(\lambda )`$. To prove (ii), let $`v0`$ be the image of $`w_\lambda `$ in $`V`$. Notice that $`\text{dim}V_\lambda =1`$ implies that $$\mathrm{\Lambda }_{\beta ,m}.v=\pi _{\beta ,m}v,$$ for some scalars $`\pi _{\beta ,m}𝐂`$. By Proposition 1.1(i), it follows that $$\pi _{\beta ,m}=0(\|m\|>\lambda (h_\beta )).$$ For $`iI`$, set $`\pi _i(u)=_{k=0}^{\lambda (h_i)}\pi _{\alpha _i,k}u^k`$. The $`\pi _i(u)`$ are polynomials with constant term $`1`$ and Proposition 1.1(v) shows that $$\mathrm{\Lambda }_i^\pm (u).v=\pi _i^\pm (u)v.$$ This shows that $`V`$ is a quotient of $`W(𝝅)`$, where $`𝝅`$ is the $`n`$-tuple of polynomials defined above. The proof of (iii) is identical. ∎ ## 3. A tensor product theorem for $`W(𝝅)`$ The main result of this section is the following theorem. ###### Theorem 2. Let $`𝝅=(\pi _1(u),\pi _2(u),\mathrm{},\pi _n(u))`$ and $`\stackrel{\mathbf{~}}{𝝅}=(\stackrel{~}{\pi }_1(u),\stackrel{~}{\pi }_2(u),\mathrm{},\stackrel{~}{\pi }_n(u))`$ be $`n`$-tuples of polynomials in $`u`$ with constant term $`1`$, such that $`\pi _i`$ and $`\stackrel{~}{\pi }_j`$ are coprime for all $`1i,jn`$. Then, $$W(𝝅\stackrel{\mathbf{~}}{𝝅})W(𝝅)W(\stackrel{\mathbf{~}}{𝝅})$$ as $`𝐔`$-modules, where $`𝝅\stackrel{\mathbf{~}}{𝝅}=(\pi _1\stackrel{~}{\pi }_1,\pi _2\stackrel{~}{\pi }_2,\mathrm{},\pi _n\stackrel{~}{\pi }_n)`$. For $`\beta R^+`$ and $`𝝅`$ as in Theorem 2, define $`\pi _\beta (u)`$ by (3.1) $$\mathrm{\Lambda }_\beta ^+(u).w_𝝅=\pi _\beta (u)w_𝝅.$$ ###### Lemma 3.1. Let $`\beta =_{i=1}^nr_i\alpha _iR^+`$. Then, $$\pi _\beta (u)=\underset{i=1}{\overset{n}{}}\pi _i^{r_id_i/d_\beta }(u).$$ If $`\theta _sR^+`$ is the highest short root of $`𝔤`$, then $`\pi _\beta `$ divides $`\pi _{\theta _s}`$ for all $`\beta R^+`$. ###### Proof. The first statement is immediate from the formula for $`h_\beta `$ in terms of the $`h_i`$ given in Section 1 and the definition of the $`\mathrm{\Lambda }_\beta ^+`$. The second statement is proved case by case, using the explicit formula for $`\theta _s`$ (there is a simple uniform proof when $`𝔤`$ is simply-laced). ∎ From now on, we shall assume that, given $`\beta R^+`$ and an $`n`$-tuple of polynomials $`𝝅`$, we have defined a polynomial $`\pi _\beta (u)`$ by the formula given in (3.1). The polynomials $`\pi _\beta ^\pm (u)`$ are defined as in (2.1). ###### Definition 3.1. Let $`𝝅=\{\pi _1,\mathrm{},\pi _n\}`$ be an $`n`$-tuple of polynomials in $`u`$ with constant term 1. Define $`M(𝝅)`$ to be the left $`𝐔`$-module obtained by taking the quotient of $`𝐔`$ by the left ideal generated by the following: $`x_{i,k}^+,h_{i,0}\text{deg}\pi _i(iI,k𝐙),`$ $`\left(\mathrm{\Lambda }_i^\pm (u)\pi _i^\pm (u)\right)_s(iI,s0),`$ $`\left(\pi _\beta (u)\stackrel{~}{X}_\beta ^{}(u)\right)_s(\beta R^+,s𝐙).`$ Let $`m_𝝅`$ be the image of 1 in $`M(𝝅)`$. It is clear from equation (3.1) and Proposition 1.1(iii) that, for all $`\beta R^+`$ and all $`s𝐙`$, $`\left(\mathrm{\Lambda }_\beta ^\pm (u)\pi _\beta ^\pm (u)\right)_s.m_𝝅`$ $`=0,`$ $`\left(\pi _\beta (u)\stackrel{~}{H}_\beta (u)\right)_s.m_𝝅`$ $`=0.`$ Set $`\lambda _𝝅(h_i)=\text{deg}\pi _i`$. ###### Lemma 3.2. We have $$M(𝝅)=\underset{\eta Q^+}{}M(𝝅)_{\lambda _𝝅\eta },$$ and $`\text{dim}M(𝛑)_{\lambda \eta }<\mathrm{}`$. Further, for all $`\beta R^+`$, $`s𝐙`$, we have $$\left(\pi _{\theta _s}(u)\stackrel{~}{X}_\beta ^{}(u)\right)_s.M(𝝅)=0.$$ ###### Proof. The set $$\{x_{\beta ,r}^{}:\beta R^+,0r<\lambda (h_\beta )\}\{x_\beta ^{}t^m\pi _\beta (t):\beta R^+,m𝐙\}$$ is a basis of $$\left(\underset{\beta R^+}{}𝐂x_\beta ^{}\right)𝐂[t,t^1].$$ By the PBW basis theorem, we can write $$𝐔(<)=𝐔_𝝅𝐔^𝝅$$ where $`𝐔_𝝅`$ (resp. $`𝐔^𝝅`$) consists of ordered monomials from the first (resp. second) set. The relation $$\left(\pi _\beta (u)\stackrel{~}{X}_\beta ^{}(u)\right)_s.m_𝝅=0$$ for all $`s𝐙`$ implies that $`(𝐔^𝝅)_+.m_𝝅=0`$. A further use of the PBW theorem now shows that $$M(𝝅)𝐔_𝝅$$ as vector spaces. Moreover, this isomorphism takes $`M(𝝅)_{\lambda _𝝅\eta }`$ to $$𝐔_𝝅(\eta )=\{x𝐔_𝝅:[h,x]=\eta (h)xh𝔥\}.$$ Since this space is clearly finite–dimensional, the first statement of the lemma follows. For the second statement, we show by induction on $`\text{ht}\eta `$ that $$\left(\pi _\theta (u)\stackrel{~}{X}_\beta ^{}(u)\right)_s.M(𝝅)_{\lambda _𝝅\eta }=0(\beta R^+,s𝐙).$$ If $`\eta =0`$, the result is immediate from the definition of $`M(𝝅)`$ and the last part of Lemma 3.1. In general, let $`x_{\beta _1,r_1}^{}\mathrm{}x_{\beta _k,r_k}^{}.m_𝝅M(𝝅)_{\lambda _𝝅\eta }`$. Then, $`\left(\pi _{\theta _s}(u)\stackrel{~}{X}_\beta ^{}(u)\right)_sx_{\beta _1,r_1}^{}\mathrm{}x_{\beta _k,r_k}^{}.m_𝝅`$ $`=x_{\beta _1,r_1}^{}\left(\pi _\theta (u)\stackrel{~}{X}_\beta ^{}(u)\right)_sx_{\beta _2,r_2}^{}\mathrm{}x_{\beta _k,r_k}^{}.m_𝝅`$ $`+\left(\pi _\theta (u)\stackrel{~}{X}_{\beta +\beta _1}^{}(u)\right)_{s+r_1}x_{\beta _2,r_2}^{}\mathrm{}x_{\beta _k,r_k}^{}.m_𝝅,`$ where we understand that $`\stackrel{~}{X}_{\beta +\beta _1}^{}(u)=0`$ if $`\beta +\beta _1R^+`$. The right-hand side is zero by the induction hypothesis, and the inductive step is complete. ∎ ###### Lemma 3.3. The $`𝐔`$-module $`W(𝛑)`$ is a quotient of $`M(𝛑)`$, and any finite-dimensional quotient of $`M(𝛑)`$ is a quotient of $`W(𝛑)`$. ###### Proof. Let $`V`$ be a finite-dimensional quotient of $`M(𝝅)`$, let $`vV`$ be the image of $`m_𝝅`$, and let $`\lambda =\lambda _𝝅`$. Then, $`\text{dim}V_\lambda =\text{dim}M(𝝅)_\lambda =1`$, so by Proposition 2.1(ii), $`V`$ is a quotient of some $`W(\stackrel{~}{𝝅})`$ with $`\lambda =\lambda _{\stackrel{~}{𝝅}}`$. Since $`\text{dim}W(\stackrel{~}{𝝅})_\lambda =1`$, $`v`$ is a scalar multiple of the image of $`w_{\stackrel{~}{𝝅}}`$. But then by comparing the action of $`\mathrm{\Lambda }_i^\pm (u)`$ on $`w_{\stackrel{~}{𝝅}}`$ and on $`m_𝝅`$, we see that $`𝝅=\stackrel{~}{𝝅}`$. To show that $`W(𝝅)`$ is a quotient of $`M(𝝅)`$, it is clear from the definitions of these modules that we only need to show that $$\left(\pi _\beta (u)\stackrel{~}{X}_\beta ^{}(u)\right)_s.w_𝝅=0(s𝐙).$$ Since $`W(𝝅)`$ is a quotient of $`W(\lambda _{𝝅,0})`$, this follows from Proposition 1.1, thus completing the proof of the lemma. ∎ Denote by $`\mathrm{\Delta }:𝐔𝐔𝐔`$ the comultiplication of $`𝐔`$ defined by extending to an algebra homomorphism the assignment $$xx1+1x,$$ for all $`xL(𝔤)`$. The following lemma is proved in \[G\]. ###### Lemma 3.4. For all $`\beta R^+`$, $$\mathrm{\Delta }(\mathrm{\Lambda }_\beta ^\pm )=\mathrm{\Lambda }_\beta ^\pm \mathrm{\Lambda }_\beta ^\pm ,$$ where $$\mathrm{\Lambda }_\beta ^\pm \mathrm{\Lambda }_\beta ^\pm =\underset{k,m0}{}(\mathrm{\Lambda }_{\beta ,\pm k}\mathrm{\Lambda }_{\beta ,\pm m})u^{k+m}.$$ Theorem 2 is now clearly a consequence of the following proposition. ###### Proposition 3.1. Assume that $`𝛑=\{\pi _1,\pi _2,\mathrm{},\pi _n\}`$ and $`\stackrel{\mathbf{~}}{𝛑}=\{\stackrel{~}{\pi }_1,\stackrel{~}{\pi }_2,\mathrm{},\stackrel{~}{\pi }_n\}`$ are $`n`$-tuples of polynomials with constant term $`1`$, such that $`\pi _i`$ and $`\stackrel{~}{\pi }_j`$ are coprime for all $`1i,jn`$. Then: 1. $`M(𝝅\stackrel{\mathbf{~}}{𝝅})M(𝝅)M(\stackrel{\mathbf{~}}{𝝅})`$; 2. every finite-dimensional quotient $`𝐔`$-module of $`M(𝝅)M(\stackrel{\mathbf{~}}{𝝅})`$ is a quotient of $`W(𝝅)W(\stackrel{~}{𝝅})`$. ###### Proof. Set $`\lambda =\lambda _𝝅+\lambda _{\stackrel{~}{𝝅}}`$. It is clear from the proof of Lemma 3.2 that, for all $`\eta Q_+`$, we have $$M(𝝅\stackrel{\mathbf{~}}{𝝅})_{\lambda \eta }\left(M(𝝅)M(\stackrel{\mathbf{~}}{𝝅})\right)_{\lambda \eta }$$ as (finite-dimensional) vector spaces. To prove (i), it therefore suffices to prove that there exists a surjective homomorphism of $`𝐔`$-modules $`M(𝝅\stackrel{\mathbf{~}}{𝝅})M(𝝅)M(\stackrel{\mathbf{~}}{𝝅})`$. It is easy to see, using Lemma 3.4, that the element $`m_𝝅m_{\stackrel{\mathbf{~}}{𝝅}}`$ satisfies the defining relations of $`M(𝝅\stackrel{\mathbf{~}}{𝝅})`$, so there exists a $`𝐔`$-module map $`M(𝝅\stackrel{\mathbf{~}}{𝝅})M(𝝅)M(\stackrel{\mathbf{~}}{𝝅})`$ that sends $`m_{𝝅\stackrel{\mathbf{~}}{𝝅}}`$ to $`m_𝝅m_{\stackrel{\mathbf{~}}{𝝅}}`$. Thus, to prove (i), we must show that, if $`\pi _i`$ and $`\stackrel{~}{\pi }_j`$ have no roots in common, the element $`m_𝝅m_{\stackrel{\mathbf{~}}{𝝅}}`$ generates $`M(𝝅)M(\stackrel{\mathbf{~}}{𝝅})`$ as a $`𝐔`$-module. Set $$N=𝐔.(m_𝝅m_{\stackrel{\mathbf{~}}{𝝅}}).$$ Assume that, for all $`\eta =_ir_i\alpha _i`$, $`\stackrel{~}{\eta }=_i\stackrel{~}{r}_i\alpha _i`$, with $`\text{ht}\eta =_ir_i<s`$, $`\text{ht}\stackrel{~}{\eta }=_i\stackrel{~}{r}_i<s`$, we have $$M(𝝅)_{\lambda _𝝅\eta }M(\stackrel{\mathbf{~}}{𝝅})_{\lambda _{\stackrel{\mathbf{~}}{𝝅}}\stackrel{~}{\eta }}N.$$ We shall prove that $`(x_{i,m}^{}.M(𝝅)_{\lambda _𝝅\eta })M(\stackrel{\mathbf{~}}{𝝅})_{\lambda _{\stackrel{\mathbf{~}}{𝝅}}\stackrel{~}{\eta }}N,`$ $`M(𝝅)_{\lambda _𝝅\eta }x_{i,m}^{}.(M(\stackrel{\mathbf{~}}{𝝅})_{\lambda _{\stackrel{\mathbf{~}}{𝝅}}\stackrel{~}{\eta }})N`$ for all $`iI`$, $`m𝐙`$. This will prove that $$M(𝝅)_{\lambda _𝝅\eta }M(\stackrel{\mathbf{~}}{𝝅})_{\lambda _{\stackrel{\mathbf{~}}{𝝅}}\stackrel{~}{\eta }}N,$$ when $`\text{ht}\eta s`$, $`\text{ht}\stackrel{~}{\eta }s`$, and hence, by induction on $`s`$, that $`N=M(𝝅)M(\stackrel{\mathbf{~}}{𝝅})`$. Since $`\pi _{\theta _s}`$ and $`\stackrel{~}{\pi }_{\theta _s}`$ are coprime, we can choose polynomials $`R(u)`$, $`\stackrel{~}{R}(u)`$ such that $$R\pi _{\theta _s}+\stackrel{~}{R}\stackrel{~}{\pi }_{\theta _s}=1.$$ By the second part of Lemma 3.2, $`\left(R\pi _{\theta _s}\stackrel{~}{X}_i^{}(u)\right)_m.w`$ $`=0,`$ $`\left(\stackrel{~}{R}\stackrel{~}{\pi }_{\theta _s}\stackrel{~}{X}_i^{}(u)\right)_m.\stackrel{~}{w}`$ $`=0,`$ for all $`iI`$, $`m𝐙`$, $`wM(𝝅)`$, $`\stackrel{~}{w}M(\stackrel{~}{𝝅})`$. Hence, $`\left(R\pi _{\theta _s}\stackrel{~}{X}_i^{}(u)\right)_m`$ $`.(w\stackrel{~}{w})`$ $`=\left(R\pi _{\theta _s}\stackrel{~}{X}_i^{}(u)\right)_m.w\stackrel{~}{w}+w\left(R\pi _{\theta _s}\stackrel{~}{X}_i^{}(u)\right)_m.\stackrel{~}{w}`$ $`=w\left((1\stackrel{~}{R}\stackrel{~}{\pi }_{\theta _s})\stackrel{~}{X}_i^{}(u)\right)_m.\stackrel{~}{w}`$ $`=wx_{i,m}^{}.\stackrel{~}{w}.`$ Taking $`wM(𝝅)_{\lambda _𝝅\eta }`$ and $`\stackrel{~}{w}M(\stackrel{\mathbf{~}}{𝝅})_{\lambda _{\stackrel{\mathbf{~}}{𝝅}}\eta }`$, so that $`w\stackrel{~}{w}N`$, it follows that $`wx_{i,m}^{}.\stackrel{~}{w}N`$ for all $`iI`$, $`m𝐙`$. The other inclusion is proved similarly, and the proof of part (i) is complete. Suppose that $`V`$ is a finite-dimensional quotient of $`M(𝝅)M(\stackrel{\mathbf{~}}{𝝅})`$ with kernel $`K`$. We shall prove that, for all $`r1`$, $`iI`$, $`s\lambda _𝝅(h_i)+1`$, $`\stackrel{~}{s}\lambda _{\stackrel{\mathbf{~}}{𝝅}}(h_i)+1`$, (3.2) $`\left(\pi _i(u)X_i^{}(u)^r\right)_s.m_𝝅M(\stackrel{\mathbf{~}}{𝝅})K,`$ (3.3) $`M(𝝅)\left(\pi _i(u)X_i^{}(u)^r\right)_s.m_{\stackrel{\mathbf{~}}{𝝅}}K.`$ Since the sum of these subspaces is the kernel of the quotient map $$M(𝝅)M(\stackrel{\mathbf{~}}{𝝅})W(𝝅)W(\stackrel{\mathbf{~}}{𝝅}),$$ it follows that $`V`$ is a quotient of $`W(𝝅)W(\stackrel{\mathbf{~}}{𝝅})`$, which proves part (ii). To prove equation (3.2), it suffices (by the proof of Proposition 1.1(ii)) to prove that, for all $`iI`$, $`m𝐙`$, $$(x_{i,m}^{})^{r_i+1}.m_𝝅m_{\stackrel{\mathbf{~}}{𝝅}}K,$$ where $`r_i=\text{deg}\pi _i=\lambda _𝝅(h_i)`$. Since $`V`$ is finite-dimensional, the element $`(x_{i,m}^{})^r.m_𝝅m_{\stackrel{\mathbf{~}}{𝝅}}K`$ for some $`r0`$. Let $`r_0`$ be the smallest value of $`r`$ with this property. Since $$x_{i,m}^+(x_{i,m}^{})^{r_0}.m_𝝅m_{\stackrel{\mathbf{~}}{𝝅}}=(r_ir_0+1)(x_{i,m}^{})^{r_01}.m_𝝅m_{\stackrel{\mathbf{~}}{𝝅}},$$ it follows by the minimality of $`r_0`$ that $`r_i+1=r_0`$. Equation (3.3) is proved similarly, and we are done.∎ Note that, since $`\text{dim}W(𝝅)_{\lambda _𝝅}=1`$, it follows that $`W(𝝅)`$ has a unique irreducible quotient $`V(𝝅)`$. Write $`𝝅`$ as a product $$𝝅=𝝅^{(1)}𝝅^{(2)}\mathrm{}𝝅^{(k)},$$ where $`𝝅^{(j)}`$ is such that $$\pi _{\theta _s}^{(j)}=(1a_ju)^{m_j},$$ for some $`m_j>0`$ and $`a_j𝐂^\times `$ with $`a_ja_k`$ if $`jk`$. The following result was proved in \[CP1\]. ###### Proposition 3.2. With the above notation, $$V(𝝅)V(𝝅^{(1)})V(𝝅^{(2)})\mathrm{}V(𝝅^{(k)})$$ as $`L(𝔤)`$-modules. Further, $$V(𝝅^{(j)})V^{fin}(\lambda _{𝝅_j})$$ as $`𝔤`$-modules. ∎ ## 4. The quantum case In the remainder of this paper, we shall assume that $`𝔤`$ is simply-laced. Let $`q`$ be an indeterminate, let $`𝐂(q)`$ be the field of rational functions in $`q`$ with complex coefficients, and let $`𝐀=𝐂[q,q^1]`$ be the subring of Laurent polynomials. For $`r,m𝐍`$, $`mr`$, define $$[m]=\frac{q^mq^m}{qq^1},[m]!=[m][m1]\mathrm{}[2][1],\left[\begin{array}{c}m\\ r\end{array}\right]=\frac{[m]!}{[r]![mr]!}.$$ Then, $`\left[\begin{array}{c}m\\ r\end{array}\right]𝐀`$. Let $`𝐔_q^e`$ be the quantized enveloping algebra over $`𝐂(q)`$ associated to $`L^e(𝔤)`$. Thus, $`𝐔_q^e`$ is the quotient of the quantum affine algebra obtained by setting the central generator equal to $`1`$. It follows from \[Dr2\], \[B\], \[J\] that $`𝐔_q^e`$ is the algebra with generators $`𝐱_{i,r}^\pm `$ ($`iI`$, $`r𝐙`$), $`K_i^{\pm 1}`$ ($`iI`$), $`𝐡_{i,r}`$ ($`iI`$, $`r𝐙\backslash \{0\}`$), $`D^{\pm 1}`$, and the following defining relations: $`K_iK_i^1=K_i^1K_i=1,`$ $`DD^1=D^1D=1,`$ $`K_iK_j=K_jK_i,`$ $`K_i𝐡_{j,r}=𝐡_{j,r}K_i,`$ $`K_i𝐱_{j,r}^\pm K_i^1=q^{\pm a_{ij}}𝐱_{j,r}^\pm ,`$ $`D𝐱_{j,r}^\pm D^1=q^r𝐱_{j,r}^\pm ,`$ $`[𝐡_{i,r},𝐡_{j,s}]=0,`$ $`[𝐡_{i,r},𝐱_{j,s}^\pm ]=\pm {\displaystyle \frac{1}{r}}[ra_{ij}]𝐱_{j,r+s}^\pm ,`$ $`𝐱_{i,r+1}^\pm 𝐱_{j,s}^\pm q^{\pm a_{ij}}𝐱_{j,s}^\pm 𝐱_{i,r+1}^\pm `$ $`=q^{\pm a_{ij}}𝐱_{i,r}^\pm 𝐱_{j,s+1}^\pm 𝐱_{j,s+1}^\pm 𝐱_{i,r}^\pm ,`$ $`[𝐱_{i,r}^+,𝐱_{j,s}^{}]=\delta _{i,j}`$ $`{\displaystyle \frac{\psi _{i,r+s}^+\psi _{i,r+s}^{}}{qq^1}},`$ $`{\displaystyle \underset{\pi \mathrm{\Sigma }_m}{}}{\displaystyle \underset{k=0}{\overset{m}{}}}(1)^k\left[\begin{array}{c}m\\ k\end{array}\right]𝐱_{i,r_{\pi (1)}}^\pm \mathrm{}𝐱_{i,r_{\pi (k)}}^\pm `$ $`𝐱_{j,s}^\pm 𝐱_{i,r_{\pi (k+1)}}^\pm \mathrm{}𝐱_{i,r_{\pi (m)}}^\pm =0,\text{if }ij,`$ for all sequences of integers $`r_1,\mathrm{},r_m`$, where $`m=1a_{ij}`$, $`\mathrm{\Sigma }_m`$ is the symmetric group on $`m`$ letters, and the $`\psi _{i,r}^\pm `$ are determined by equating powers of $`u`$ in the formal power series $$\underset{r=0}{\overset{\mathrm{}}{}}\psi _{i,\pm r}^\pm u^{\pm r}=K_i^{\pm 1}exp\left(\pm (qq^1)\underset{s=1}{\overset{\mathrm{}}{}}𝐡_{i,\pm s}u^{\pm s}\right).$$ Define the $`q`$-divided powers $$(𝐱_{i,k}^\pm )^{(r)}=\frac{(𝐱_{i,k}^\pm )^r}{[r]!},$$ for all $`iI`$, $`k𝐙`$, $`r0`$. Suppose that $`a_{ij}=1`$. Then, a special case of the above relations is $$(𝐱_{i,s}^{})^2𝐱_{j,r}^{}(q+q^1)𝐱_{i,s}^{}𝐱_{j,r}^{}𝐱_{i,s}^{}+𝐱_{j,r}^{}(𝐱_{i,s}^{})^2=0.$$ Set $`\gamma _{s,r}^{i,j}=𝐱_{i,s}^{}𝐱_{j,r}^{}q𝐱_{j,r}^{}𝐱_{i,s}^{}`$. Again, the relations in $`𝐔_q^e`$ imply that $$\gamma _{s,r}^{i,j}=\gamma _{r+1,s1}^{j,i}.$$ The following result is proved in \[L2\]. ###### Proposition 4.1. For $`i,jI`$ with $`a_{ij}=1`$, $`r,s,l,m𝐙`$, and $`l,m0`$, there exist $`f_p𝐀`$, for $`0p\text{min}(l,m)`$, such that $$(𝐱_{i,s}^{})^{(l)}(𝐱_{j,r}^{})^{(m)}=\underset{p}{}f_p(𝐱_{j,r}^{})^{(mp)}(\gamma _{s,r}^{i,j})^{(p)}(𝐱_{i,s}^{})^{(lp)}.$$ Further, there exists $`g_p𝐀`$, for $`0pm`$, such that $$(\gamma _{s,r}^{i,j})^{(m)}=\underset{p}{}g_p(𝐱_{j,r}^{})^{(p)}(𝐱_{i,s}^{})^{(m)}(𝐱_{j,r}^{})^{(mp)}.$$ Define $$𝚲_i^\pm (u)=\underset{m=0}{\overset{\mathrm{}}{}}𝚲_{i,\pm m}u^m=\text{exp}\left(\underset{k=1}{\overset{\mathrm{}}{}}\frac{𝐡_{i,\pm k}}{[k]}u^k\right).$$ The subalgebras $`𝐔_q`$, $`𝐔_q^{fin}`$, $`𝐔_q(<)`$, etc., are defined in the obvious way. Let $`𝐔_q(𝔥^e)`$ be the subalgebra of $`𝐔_q^e`$ generated by $`K_i^{\pm 1}`$ ($`iI`$) and $`D^{\pm 1}`$. Let $`𝐔_q(0)`$ be the subalgebra of $`𝐔_q^e`$ generated by the elements $`𝚲_{i,m}`$ ($`iI`$, $`m𝐙`$). The following result is a simple corollary of the PBW theorem for $`𝐔_q^e`$ \[B\]. ###### Lemma 4.1. $`𝐔_q^e𝐔_q(<)𝐔_q(0)𝐔_q(𝔥^e)𝐔_q(>)`$. ∎ For any invertible element $`x𝐔_q^e`$, define $$\left[\begin{array}{c}x\\ r\end{array}\right]=\frac{xq^rx^1q^r}{qq^1}.$$ Let $`𝐔_𝐀^e`$ be the $`𝐀`$-subalgebra of $`𝐔_q^e`$ generated by the $`K_i^{\pm 1}`$, $`(𝐱_{i,k}^\pm )^{(r)}`$ ($`iI`$, $`k𝐙`$, $`r0`$), $`D^{\pm 1}`$, and $`\left[\begin{array}{c}D\\ r\end{array}\right]`$ ($`r1`$). Then, \[L2\], \[BCP\], $$𝐔_q^e=𝐂(q)_𝐀𝐔_𝐀^e.$$ Define $`𝐔_𝐀(<)`$, $`𝐔_𝐀(0)`$ and $`𝐔_𝐀(>)`$ in the obvious way. Let $`𝐔_𝐀(𝔥^e)`$ be the $`𝐀`$–subalgebra of $`𝐔_𝐀`$ generated by the elements $`K_i^{\pm 1}`$, $`D^{\pm 1}`$, $`\left[\begin{array}{c}K_i\\ r\end{array}\right]`$ and $`\left[\begin{array}{c}D\\ r\end{array}\right]`$ ($`iI`$, $`r𝐙`$). The following is proved as in Proposition 2.7 in \[BCP\]. ###### Proposition 4.2. $`𝐔_𝐀^e=𝐔_𝐀(<)𝐔_𝐀(0)𝐔_𝐀^e(𝔥)𝐔_𝐀(>)`$. ∎ The next lemma is easily checked. ###### Lemma 4.2. 1. There is a unique $`𝐂`$–linear anti–automorphism $`\mathrm{\Psi }`$ of $`𝐔_q^e`$ such that $`\mathrm{\Psi }(q)=q^1`$ and $`\mathrm{\Psi }(K_i)=K_i,`$ $`\mathrm{\Psi }(D)=D,`$ $`\mathrm{\Psi }(x_{i,r}^\pm )=x_{i,r}^\pm ,`$ $`\mathrm{\Psi }(h_{i,r})=h_{i,r},`$ for all $`iI`$, $`r𝐙`$. 2. There is a unique algebra automorphism $`\mathrm{\Phi }`$ of $`𝐔_q^e`$ over $`𝐂(q)`$ such that $`\mathrm{\Phi }(𝐱_{i,r}^\pm )=𝐱_{i,r}^\pm `$, $`\mathrm{\Phi }(\mathrm{\Lambda }_i^+(u))=\mathrm{\Lambda }_i^{}(u)`$. The first part of the following result is proved in \[BCP\], and the second part follows from it by applying $`\mathrm{\Psi }`$. ###### Lemma 4.3. 1. Let $`s>s^{}`$, $`m,m^{}0`$, $`iI`$. Then, $`(𝐱_{i,s}^{})^{(m)}(𝐱_{i,s^{}}^{})^{(m^{})}`$ is in the span of the elements $$(𝐱_{i,r_1}^{})^{(k_1)}(𝐱_{i,r_2}^{})^{(k_2)}\mathrm{}(𝐱_{i,r_l}^{})^{(k_l)},$$ where $`s^{}r_1<r_2<\mathrm{}<r_ls`$, $`_pk_p=m+m^{}`$ and $`_pk_pr_p=ms+m^{}s^{}`$. 2. Let $`s<s^{}`$, $`m,m^{}0`$, $`iI`$. Then, $`(𝐱_{i,s}^{})^{(m)}(𝐱_{i,s^{}}^{})^{(m^{})}`$ is in the span of the elements $$(𝐱_{i,r_1}^{})^{(k_1)}(𝐱_{i,r_2}^{})^{(k_2)}\mathrm{}(𝐱_{i,r_l}^{})^{(k_l)},$$ where $`s^{}r_1>r_2>\mathrm{}>r_ls`$, $`k_p=m+m^{}`$ and $`k_pr_p=ms+m^{}s^{}`$. For $`iI`$, let $`\stackrel{~}{𝐗}_i^{}(u)`$, $`𝐗_i^{}(u)`$, $`𝐗_{i,0}^{}(u)`$ be the formal power series with coefficients in $`𝐔_q`$ given by the same formulas as the $`\stackrel{~}{X}_i^{}(u)`$, etc., in Section 1. The next result is a reformulation of results in \[CP6, Section 5\]. ###### Lemma 4.4. Let $`sr1`$, $`iI`$. 1. $$(𝐱_{i,0}^+)^{(r)}(𝐱_{i,1}^{})^{(s)}=(1)^r(𝐗_i^{}(u)^{(sr)}𝚲_i^+(u))_smod𝐔_q𝐔_q(>)_+.$$ 2. $$(𝐱_{i,1}^+)^{(r)}(𝐱_{i,0}^{})^{(s)}=(1)^r(𝐗_{i,0}^{}(u)^{(sr)}𝚲_i^+(u))_smod𝐔_q𝐔_q(>)_+.$$ 3. $$\left(𝐗_i^{}(u)^{(r)}\right)_{s+r}=\mu (s_0,s_1,\mathrm{})(𝐱_{i,0}^{})^{(s_0)}(𝐱_{i,1}^{})^{(s_1)}\mathrm{},$$ where the sum is over those non-negative integers $`s_0,s_1,\mathrm{}`$ such that $`_ts_t=r`$ and $`_tts_t=s+r`$, and the coefficients $`\mu (s_0,s_1,\mathrm{})𝐀`$. In particular, the coefficient of $`(𝐱_{i,s+1}^{})^{(r)}`$ in $`\left(𝐗_i^{}(u)^{(r)}\right)_{(s+1)r}`$ is $`q^{sr(r1)}`$. ∎ ###### Definition 4.1. A $`𝐔_q^e`$–module $`V_q`$ is said to be of type 1 if $$V_q=\underset{\lambda P^e}{}(V_q)_\lambda ,$$ where $$(V_q)_\lambda =\{vV_q:K_i.v=q^{\lambda (h_i)}viI,D.v=q^{\lambda (d)}.v\}.$$ The subspaces $`(V_q)_\lambda ^+`$ are defined in the obvious way. We say that a type 1 module is integrable if the elements $`𝐱_{i,k}^\pm `$ act locally nilpotently on $`V_q`$ for all $`iI`$ and $`k𝐙`$. As in the classical case, one shows \[L2\] that, if $`V_q`$ is integrable, then $$(V_q)_\lambda 0(V_q)_{\sigma (\lambda )}0\sigma W.$$ The type 1 $`𝐔_q^{fin}`$-modules and their weight spaces are defined analogously. If $`\lambda P_+`$, there is a unique finite-dimensional irreducible $`𝐔_q^{fin}`$-module $`V_q^{fin}(\lambda )`$ generated by a vector $`v`$ such that $$k_i.v=q^{\lambda (h_i)}v,x_{i,0}^+.v=0,(x_{i,0}^{})^{\lambda (h_i)+1}.v=0,$$ for all $`iI`$. Further, $$\text{dim}_{𝐂(q)}(V_q^{fin}(\lambda )_\mu )=dim_𝐂(V^{fin}(\lambda )_\mu )$$ for all $`\mu P`$. From now on, we shall only consider modules of type 1. The next result is the quantum analogue of Proposition 1.1 and is proved in exactly the same way. ###### Proposition 4.3. Let $`V_q`$ be an integrable $`𝐔_q^e`$-module. Let $`\lambda P_+^e`$ and assume that $`0v(V_q)_\lambda ^+`$. 1. $`𝚲_{i,m}.v=0`$ for all $`iI`$ and $`\|m\|>\lambda (h_i)`$. 2. $`𝚲_{i,\lambda (h_i)}𝚲_{i,m}.v=𝚲_{i,\lambda (h_i)m}.v`$ for all $`iI`$ and $`0m\lambda (h_i)`$. 3. For $`r1`$, $`s>\lambda (h_i)`$, $`m𝐙`$, $`iI`$, $`\left(𝐗_i^{}(u)^r𝚲_i^+(u)\right)_s.v=0,`$ $`\left(𝐗_{i,0}^{}(u)^r𝚲_i^+(u)\right)_s.v=0,`$ $`\left(\stackrel{~}{𝐗}_i^{}(u)𝚲_i^+(u)\right)_m.v=0,`$ $`\left(\stackrel{~}{𝐇}_i(u)𝚲_i^+(u)\right)_m.v=0.`$ 4. $$\left(\mathrm{\Phi }(𝐗_i^{}(u)^r)𝚲_i^{}(u)\right)_s.v=0,\left(\mathrm{\Phi }(𝐗_{i,0}^{}(u)^r)𝚲_i^{}(u)\right)_s.v=0.$$ ###### Proposition 4.4. Let $`V_q`$ be an integrable type 1 $`𝐔_q^e`$-module and assume that $`\lambda P_+^e`$, $`0v(V_q)_\lambda `$ is such that $`V_q=𝐔_q^e.v`$ and $$𝐱_{i,k}^+.v=0iI,k𝐙.$$ Then, there exists $`s_\lambda 0`$ such that $`V_q`$ is spanned by the elements $$(𝐱_{i_1,s_1}^{})^{(l_1)}(𝐱_{i_2,s_2}^{})^{(l_2)}\mathrm{}(𝐱_{i_k,s_k}^{})^{(l_k)}𝐔_𝐀(0).v$$ for $`0jk`$, $`l_j0`$, $`i_jI`$, $`0s_js_\lambda `$. ###### Proof. For any $`N0`$, let $`V_N`$ be the $`𝐂(q)`$-subspace of $`V_q`$ spanned by the elements $$(𝐱_{i_1,s_1}^{})^{(l_1)}(𝐱_{i_2,s_2}^{})^{(l_2)}\mathrm{}(𝐱_{i_k,s_k}^{})^{(l_k)}𝐔_𝐀(0).v$$ for $`0jk`$, $`0s_jN`$, $`l_j0`$. By Lemma 4.1, we have $`V_q=𝐔_q(<)𝐔_q(0).v`$ and hence $$V_q=\underset{\eta ,m}{}(V_q)_{\lambda \eta +m\delta },$$ where $`\eta Q^+`$, $`m𝐙`$. The argument given in the proof of Proposition 1.2(i) (but replacing the modules by their quantum analogues) shows that $`(V_q)_{\lambda \eta +m\delta }0`$ for only finitely many $`\eta Q^+`$. Hence, it suffices to prove that: for each $`\eta Q^+`$, there exists $`N(\eta )0`$ such that $`(V_q)_{\lambda \eta +m\delta }V_{N(\eta )}`$ for all $`m𝐙`$. We proceed by induction on $`\text{ht}\eta `$. By Proposition 4.3, we see that, if $`s>\lambda (h_i)`$, $`p1`$, $$\left(𝐗_i^{}(u)^{(p)}𝚲_i^+(u)\right)_{ps}𝐔_𝐀(0).v_\lambda =0,$$ or equivalently that (4.1) $$\underset{l=0}{\overset{ps}{}}\left(𝐗_i^{}(u)^{(p)}\right)_l𝚲_i^+(u)_{psl}𝐔_𝐀(0).v_\lambda =0.$$ If $`p=1`$, it follows easily from equation (4.1) by induction on $`s`$ that $$𝐱_{i,s}^{}.(V_q)_{\lambda +m\delta }V_{\lambda (h_i)}$$ if $`s>\lambda (h_i)`$. To deal with the case, $`s0`$, we apply $`h_{i,s}`$ to both sides of equation (4.1), and as in the proof of Proposition 1.2 use the fact that $`\mathrm{\Lambda }_{i,\lambda (h_i)}`$ is invertible. Thus, the induction begins with $`N(\alpha _i)=\lambda (h_i)`$. Assume that we have proved the italicized statement above for all $`\eta Q^+`$ with $`\text{ht}\eta <p`$. We deal first with the case $`\eta =p\alpha _i`$, for some $`iI`$, $`p1`$. We show that we can take $`N(p\alpha _i)=\lambda (h_i)`$. For this, it suffices to prove that (4.2) $$(𝐱_{i,s}^{})^{(l)}(V_q)_{\lambda +(msl)\delta (pl)\alpha _i}V_{\lambda (h_i)}$$ for all $`m,s𝐙`$, $`pl>0`$. We prove this by induction on $`s`$ (for $`p`$ fixed), assuming first that $`s0`$. The induction begins since there is nothing to prove if $`0s\lambda (h_i)`$. Assume that $`s>\lambda (h_i)`$ and that the result holds for all smaller values of $`s0`$. If $`p>l>0`$ then, by the induction on $`p`$, we have $$(V_q)_{\lambda (pl)\alpha _i+(msl)\delta }\underset{s^{},l^{}}{}(𝐱_{i,s^{}}^{})^{(l^{})}V_{\lambda (pll^{})\alpha _i+(msls^{}l^{})\delta },$$ where $`0s^{}\lambda (h_i)`$, $`l^{}>0`$. By Lemma 4.3(i), we have $$(𝐱_{i,s}^{})^{(l)}(𝐱_{i,s^{}}^{})^{(l^{})}V_{\lambda (pll^{})\alpha _i+(msls^{}l^{})\delta }(𝐱_{i,s^{\prime \prime }}^{})^{(l^{\prime \prime })}V_{\lambda (pl^{\prime \prime })\alpha _i+(ml^{\prime \prime }s^{\prime \prime })\delta },$$ where $`s^{}s^{\prime \prime }<s`$. It follows by the induction hypothesis on $`s`$ that $$(𝐱_{i,s^{\prime \prime }}^{})^{(l^{\prime \prime })}V_{\lambda (pl^{\prime \prime })\alpha _i+(ml^{\prime \prime }s^{\prime \prime })\delta }V_{\lambda (h_i)},$$ which proves equation (4.2). In the case $`p=l`$, we must show that $$(𝐱_{i,s}^{})^{(p)}𝐔_𝐀(0).vV_{\lambda (h_i)}.$$ Since $`ps>\lambda (h_i)`$, by Proposition 4.3 we see that $$(𝐗_i^{}(u)^{(p)})_{ps}𝐔_𝐀(0).v+\underset{s^{}<ps}{}(𝐗_i^{}(u)^{(p)})_s^{}𝚲_{i,pss^{}}𝐔_𝐀(0).v=0.$$ Now, by Lemma 4.4(iii), we see that for $`s^{}<ps`$ the element $`\left(𝐗_i^{}(u)^{(p)}\right)_s^{}`$ is a sum of terms of the form (4.3) $$(𝐱_{i,s_1}^{})^{(r_1)}(𝐱_{i,s_2}^{})^{(r_2)}\mathrm{}$$ where $`0s_1s_2\mathrm{}`$, $`r_1,r_2,\mathrm{}>0`$, and $`s_1<s`$. The induction hypothesis on $`p`$ and $`s`$ proves that $`\left(𝐗_i^{}(u)^{(p)}\right)_s^{}𝐔_𝐀(0).vV_{\lambda (h_i)}`$. Finally, again by Lemma 4.4(iii), we have $$\left(𝐗_i^{}(u)^{(p)}\right)_{ps}=(𝐱_{i,s}^{})^{(p)}+A,$$ where $`A`$ is a sum of terms of the form in (4.3) where either $`s_1<s`$ or $`s_1=s`$ and $`r_1<p`$. If $`s_1<s`$ it follows as before that $`\left(𝐗_i^{}(u)^{(p)}\right)_s^{}𝐔_𝐀(0).vV_{\lambda (h_i)}`$, and if $`s_1=s`$ it follows by the case $`l<p`$ of equation (4.2) proved above. Thus, the induction on $`s`$ is completed in the case $`p=l`$ too. This completes the proof of (4.2) when $`s0`$. Next, consider the case when $`s0`$. The case $`p=1`$ was proved above. For $`p<l`$, the same method used for $`s0`$ works, this time using Lemma 4.2(ii). Finally, for the case $`p=l`$, we use the relation $$\left(\mathrm{\Phi }(𝐗_i^{}(u)^p)𝚲_i^{}(u)\right)_{ps}.v=0$$ and parts (i) and (ii) of Proposition 4.3, and proceed as in the case $`s0`$. We omit the details. This completes the proof of the italicized statement when $`\eta `$ is a multiple of $`\alpha _i`$ . We now turn to the case of arbitrary $`\eta =r_i\alpha _i`$ of height $`p`$. Choose $`M`$ so that if $`r_i<p`$ then $`(V_q)_{\lambda \eta +m\delta }V_M`$. As in the special case $`\eta =p\alpha _i`$, to complete the induction on $`p`$ it suffices to prove that there exists $`N0`$ such that $$(𝐱_{i,s}^{})^{(l)}(V_q)_{\lambda (\eta l\alpha _i)+(mls)\delta }V_N$$ for all $`iI`$, $`s𝐙`$, $`l>0`$. Since $`\text{ht}(\eta l\alpha _i)<p`$, it follows that $$(V_q)_{\lambda (\eta l\alpha _i)+(mls)\delta }\underset{j,s_j,l_j}{}(𝐱_{j,s_j}^{})^{(l_j)}(V_q)_{\lambda (\eta l\alpha _il_j\alpha _j)+(mlsl_js_j)\delta }$$ where $`0s_jM`$. Thus, we must show that there exists $`N`$ such that $$(𝐱_{i,s}^{})^{(l)}(𝐱_{j,s_j}^{})^{(l_j)}(V_q)_{\lambda (\eta l\alpha _il_j\alpha _j)+(mlsl_js_j)\delta }V_N,$$ for all $`i,jI`$, $`s,s_j𝐙`$, $`0s_jM`$. Assume that $`s0`$ (the case $`s0`$ is similar). If $`sM`$, there is nothing to prove. Assume that we know the result for all $`i`$, and all smaller positive values of $`s`$. If $`i=j`$, then we prove exactly as in the case $`\eta =r\alpha _i`$ that we can take $`N=M`$. If $`a_{ij}=0`$, the result is obvious. Assume now that $`a_{ij}=1`$. Then, with the notation in Proposition 4.1, we see that $$(\gamma _{s,s_j}^{i,j})^{(m)}=(1)^m(\gamma _{s_j+1,s1}^{j,i})^{(m)}=\underset{p^{}=0}{\overset{m}{}}g_p^{}(𝐱_{i,s1}^{})^{(p^{})}(𝐱_{j,s_j+1}^{})^{(m)}(𝐱_{i,s1}^{})^{(mp^{})},$$ where the $`g_p^{}𝐀`$. Using the induction hypothesis on $`s`$, we get that $$(\gamma _{s,s_j}^{i,j})^{(m)}.(V_q)_{\lambda (\eta m\alpha _im\alpha _j)+m^{}\delta }V_{M+1}$$ for all $`m^{}𝐙`$. Now, using Proposition 4.1 again, we see that $$(𝐱_{i,s}^{})^{(l)}(𝐱_{j,s_j}^{})^{(l_j)}(V_q)_{\lambda (\eta l\alpha _il_j\alpha _j)+(mlsl_js_j)\delta }V_{M+1}.$$ This proves the result.∎ Let $`I_q(0)`$ be the left ideal in $`𝐔_q^e`$ generated by the elements $`𝚲_{i,m}`$ ($`iI`$, $`\|m\|>\lambda (h_i)`$) and $`𝚲_{i,\lambda (h_i)}𝚲_{i,m}𝚲_{i,\lambda (h_i)m}`$ ($`iI`$, $`0m\lambda (h_i)`$). Let $`I_q^e(\lambda )`$ be the left ideal in $`𝐔_q^e`$ generated by $`I_q(0)`$, $`𝐔_q(>)`$, the elements $`K_i^{\pm 1}q^{\pm \lambda (h_i)}`$ ($`iI`$), and $`D^{\pm 1}q^{\pm \lambda (d)}`$. The ideal $`I_q(\lambda )`$ in $`𝐔_q(\lambda )`$ is defined in the obvious way. Let $$W_q(\lambda )=𝐔_q^e/I_q^e(\lambda )=𝐔_q/I_q(\lambda )$$ be the corresponding left $`𝐔_q^e`$-module. Let $`w_\lambda `$ be the image of $`1`$ in $`W_q(\lambda )`$. We have $$𝐔_q(0).w_\lambda 𝐔_q(0)/I_q(0)$$ as $`𝐔_q(0)`$-modules. ###### Proposition 4.5. The $`𝐔_q^e`$-module $`W_q(\lambda )`$ is integrable. ###### Proof. In \[K\], Kashiwara defines an integrable $`𝐔_q^e`$-module $`V^{max}(\lambda )`$ for all $`\lambda P_+^e`$. In fact, according to unpublished work of Kashiwara, $$V^{max}(\lambda )W_q(\lambda ),$$ and hence $`W_q(\lambda )`$ is an integrable $`𝐔_q^e`$-module. One can also give a direct proof that $`W_q(\lambda )`$ is integrable along the lines of the proof of Theorem 1. One works with the presentation of $`𝐔_q^e`$ in terms of Chevalley generators and the quantum version of Lemma 2.1 proved in \[K\], \[L2\]. We omit the details. ∎ Given any $`𝐔_q^e`$-module $`V_q`$ and a $`𝐔_𝐀`$-submodule $`V_𝐀`$ of $`V_q`$ such that $$V_q𝐂(q)_𝐀V_𝐀,$$ we set $$\overline{V_q}𝐂_1_𝐀V_𝐀,$$ where we regard $`𝐂_1`$ as an $`𝐀`$-module by letting $`q`$ act as 1. The algebra $`𝐂_1_𝐀𝐔_𝐀`$ is a quotient of $`𝐔^e`$, so $`\overline{V_q}`$ is a $`𝐔^e`$-module. Similar results hold for $`𝐔`$-modules. Let $`𝝅_q`$ be an $`n`$-tuple of polynomials with constant term 1 and coefficients in $`𝐂(q)`$. Define $`\lambda _{𝝅_q}𝒫_+^e`$ and $`𝝅_q^\pm (u)`$ as in Section 2. Let $`I_q(𝝅_q)`$ be the left ideal in $`𝐔_q`$ generated by $`I_q(\lambda _{𝝅_q})`$ and the elements $$\left(𝚲_i^\pm (u)\pi _i^\pm (u)\right)_s(iI,s0).$$ Set $`W_q(𝝅_q)=𝐔_q/I_q(𝝅_q)`$. ###### Lemma 4.5. $`W_q(𝝅_q)`$ is a finite-dimensional $`𝐔_q`$-module. ###### Proof. This is proved in the same way as the corresponding result for $`𝐔`$. We use Proposition 4.4 instead of Proposition 1.2. ∎ One has the following analogue of Proposition 2.1 for the modules $`W_q(\lambda )`$ and $`W_q(𝝅_q)`$. We omit the proof, which is entirely similar to that of Proposition 2.1. ###### Proposition 4.6. 1. Let $`V_q`$ be any integrable $`𝐔_q^e`$-module generated by an element of $`(V_q)_\lambda ^+`$. Then, $`V_q`$ is a quotient of $`W_q(\lambda )`$. 2. Let $`V_q`$ be a finite-dimensional quotient $`𝐔_q`$-module of $`W_q(\lambda )`$, and assume that $`\text{dim}(V_q)_\lambda =1`$. Then, $`V_q`$ is a quotient of $`W_q(𝝅_q)`$ for some choice of $`𝝅_q`$. 3. Let $`V_q`$ be finite-dimensional $`𝐔_q`$-module generated by an element of $`(V_q)_\lambda ^+`$ and such that $`\text{dim}(V_q)_\lambda =1`$. Then, $`V_q`$ is a quotient of $`W_q(𝝅_q)`$ for some $`𝝅_q`$. ###### Definition 4.2. We call $`𝝅_q`$ integral if the polynomials $`\pi _i^\pm (u)`$ have coefficients in $`𝐀`$ for all $`iI`$. Equivalently, for all $`iI`$, $`\pi _i(u)`$ has coefficients in $`𝐀`$ and the coefficient of the highest power of $`u`$ should lie in $`𝐂^\times q^𝐙`$. Let $`\overline{𝝅_q}`$ be the $`n`$-tuple of polynomials with coefficients in $`𝐂`$ and constant term one obtained from $`𝝅_q`$ by evaluating its coefficients at $`q=1`$. For any $`𝝅_q`$, $`W_𝐀(𝝅_q)=𝐔_𝐀.w_{𝝅_q}`$ is a $`𝐔_𝐀`$-module. ###### Lemma 4.6. Assume that $`𝛑_q`$ is integral. 1. $`W_𝐀(𝝅_q)`$ is a free $`𝐀`$-module and we have $$W_q(𝝅_q)𝐂(q)_𝐀W_𝐀(𝝅_q).$$ 2. The $`𝐔`$-module $`\overline{W_q(𝝅_q)}`$ is a quotient of $`W(\overline{𝝅_q})`$. ###### Proof. Since $`W_𝐀(𝝅_q)`$ is a quotient of $`W_q(\lambda _{𝝅_q})`$, it follows from Proposition 4.4 that $`W_𝐀(𝝅_q)`$ is a finitely-generated $`𝐀`$-module. Since it is clearly a torsion-free $`𝐀`$-module, part (i) follows. To prove (ii), observe that the defining relations of $`W_q(𝝅_q)`$ specialize to those of $`W(\overline{𝝅})`$. The result now follows from Proposition 2.1. ∎ The $`𝐔_q`$-module $`W_q(𝝅_q)`$ has a unique irreducible quotient $`V_q(𝝅_q)`$. Let $`v_{𝝅_q}`$ be the image of $`w_{𝝅_q}`$ and set $$V_𝐀(𝝅_q)=𝐔_𝐀.v_{𝝅_q}.$$ If $`𝝅_q`$ is integral, $`V_𝐀(𝝅_q)`$ is a $`𝐔_𝐀`$-module and is free as an $`𝐀`$-module (since it is torsion-free and the quotient of a finitely-generated $`𝐀`$-module). Let $`V(𝝅)`$ be the unique irreducible quotient of the $`𝐔`$-module $`W(𝝅)`$. ###### Lemma 4.7. The $`𝐔`$-module $`V(\overline{𝛑_q})`$ is a quotient of $`\overline{V_q(𝛑_q)}`$. ###### Proof. The module $`\overline{V_q(𝝅_q)}`$ has an unique irreducible quotient $`V`$. By Lemma 4.6(ii), $`V`$ is a quotient of $`W(\overline{𝝅_q})`$ and hence by uniqueness $`VV(\overline{𝝅_q})`$. ∎ We have thus proved the following statement. Assume that $`𝝅_q`$ is integral. Then, we have a commutative diagram of surjective $`𝐔`$-module homomorphisms $`W(\overline{𝝅_q})`$ $``$ $`\overline{W_q(𝝅_q)}`$ $``$ $``$ $`V(\overline{𝝅_q})`$ $``$ $`\overline{V_q(𝝅_q)}.`$ Conjecture. If $`𝝅_q=𝝅`$ has coefficients in $`𝐂`$, the natural map $`W(𝝅)\overline{V_q(𝝅)}`$ is an isomorphism of $`𝐔`$-modules, and hence $`W_q(𝝅)V_q(𝝅)`$. ∎In Section 6, we prove this conjecture when $`𝔤=sl_2`$. ## 5. An irreducibility criterion. In this section we return to the classical case and obtain a criterion for the irreducibility of the modules $`W(𝝅)`$. For any $`a𝐂^\times `$, and any $`𝐔^{fin}`$-module $`V`$, define a $`𝐔`$-module structure on $`V`$ by $$(xt^r).v=a^rx.v$$ for $`x𝔤`$, $`r𝐙`$, $`vV`$. Let $`V(a)`$ denote the corresponding $`𝐔`$-module. For $`iI`$, $`a𝐂^\times `$, we denote by $`W(i,a)`$ the $`𝐔`$-module corresponding to the $`n`$-tuple $`𝝅`$ of polynomials defined by $$\pi _j(u)=1\text{if}ji,\pi _i(u)=1au,$$ and denote $`w_𝝅`$ by $`w_{i,a}`$. Clearly, $`V^{fin}(\omega _i)(a)`$ is the irreducible quotient of $`W(i,a)`$. We set $`V^{fin}(\omega _i)(a)=V(i,a)`$. For $`iI`$ and $`a𝐂(q)^\times `$, the $`𝐔_q`$-modules $`W_q(i,a)`$ and $`V_q(i,a)`$ are defined similarly. We need the following result, due to \[CP3\] for $`𝔤`$ of type $`sl_2`$ and due to \[K2\] and \[FM\] in general. ###### Proposition 5.1. Let $`r1`$, $`a_1,\mathrm{},a_r𝐂(q)^\times `$, $`i_1,i_2,\mathrm{},i_rI`$. There is a finite set $`S𝐂(q)^\times `$ (depending on $`i_1,\mathrm{},i_r`$) such that the tensor product $$V_q(i_1,a_1)V_q(i_2,a_2)\mathrm{}V_q(i_r,a_r),$$ is irreducible if $`a_l/a_mS`$ for all $`l,m=1,2,\mathrm{},r`$. If $`𝔤=sl_2`$, $`S=\{q^{\pm 2}\}`$. ∎ ###### Proposition 5.2. For $`iI`$, $`a𝐂^\times `$, $`W(i,a)V(i,a)`$ if and only if $`r_i=1`$. ###### Proof. The proof rests on the following fact, which can be established by a case by case check: $`r_i=1`$ if and only if there exists $`\mu P_+`$ with $`0\omega _i\mu Q_+`$ such that $`V^{fin}(\mu )`$ occurs as a component of $`V^{fin}(\theta )V^{fin}(\omega _i)`$. Suppose first that $`r_i>1`$. Let $`\mu P_+`$ have the above property. For $`x𝔤`$, $`m𝐙`$, $`vV^{fin}(\omega _i)`$, $`v^{}V^{fin}(\mu )`$, define $$x_m.(v,v^{})=a^m(x.v,mpr(xv)+x.v^{}),$$ where $`pr:V^{fin}(\theta )V^{fin}(\omega _i)V^{fin}(\mu )`$ is the $`𝔤`$-module projection. It is straightforward to check that this defines a $`𝐔`$-module structure on $`V^{fin}(\omega _i)V^{fin}(\mu )`$, and that this $`𝐔`$-module is generated by the highest weight vector in $`V^{fin}(\omega _i)`$. It is therefore a quotient of $`W(i,a)`$. To prove the converse, notice that, as a $`𝔤`$-module, $`W(i,a)`$ is completely reducible, and hence $$W(i,a)V^{fin}(\omega _i)\underset{\mu <\omega _i}{}V^{fin}(\mu )^{m_\mu },$$ where $`m_\mu `$ is the multiplicity with which $`V^{fin}(\mu )`$ occurs in $`W(i,a)`$. Consider the map $`L(𝔤)W(i,a)W(i,a)`$ given by $$x_nvx_n.v.$$ This is clearly a map of $`𝔤`$-modules, where we regard $`L(𝔤)`$ as a module for $`𝔤`$ through the adjoint representation. For each $`m𝐙`$, consider the restriction of this map to $`(𝔤t^m)V(\omega _i)`$. Since $`𝔤t^mV^{fin}(\theta )`$ as $`𝔤`$-modules, we have a $`𝔤`$-module map $`V^{fin}(\theta )V(i,a)W(i,a)`$. By the fact stated above, this map takes $`(𝔤t^m)V(\omega _i)`$ into the $`𝐔(𝔤)`$-submodule $`V^{fin}(\omega _i)`$ of $`W(i,a)`$ for all $`m𝐙`$. This proves that $`V^{fin}(\omega _i)`$ is a $`𝐔`$-submodule of $`W(i,a)`$ and hence (since $`w_{i,a}V^{fin}(\omega _i)`$) is equal to $`W(i,a)`$. ∎ ###### Remark 5.1. The same criterion $`r_i=1`$ occurs, for the same reason, in Drinfeld’s work on finite-dimensional representations of Yangians, \[Dr1\]. See also \[CP4, Proposition 12.1.17\]. We can now state the main result of this section. ###### Theorem 3. Let $`𝝅=(\pi _1,\mathrm{},\pi _n)`$ be an $`n`$-tuple of polynomials in $`𝐂[u]`$ with constant coefficient one. Then, the $`𝐔`$-module $`W(𝝅)`$ is irreducible if and only if $`\pi _\theta `$ has distinct roots. ###### Proof. Assume that $`\pi _\theta `$ has distinct roots. By Lemma 3.1, it follows that $`\pi _i=1`$ if $`r_i1`$. Let $$I^{}=\{iI:r_i=1\}.$$ If $`iI^{}`$ then $`\pi _i`$ must have distinct roots, and for any $`i,jI^{}`$, $`ij`$, the polynomials $`\pi _i`$ and $`\pi _j`$ must be relatively prime. Hence, by Theorem 2 and Proposition 5.2, it follows that $$W(𝝅)\underset{iI^{},a_{ij}𝐂^\times }{}W(i,a_{ij})\underset{iI^{},a_{ij}𝐂^\times }{}V(i,a_{ij}),$$ where the $`a_{ij}`$ are all distinct. By Proposition 3.2, we see that the second tensor product in the preceding equation is an irreducible $`L(𝔤)`$-module. For the converse, suppose that $`\pi _\theta `$ has repeated roots. By Theorem 2, it follows that $`W(𝝅)`$ is isomorphic to a tensor product of modules $`W(𝝅^a)`$, where $`𝝅^a`$ is an $`n`$-tuple of polynomials such that $`\pi _\theta ^a=(1au)^m`$ for some $`a𝐂^\times `$ and $`m1`$, and where $`m>1`$ for at least one value of $`a`$. Thus, it suffices to prove the theorem in the case where $`\pi _\theta (u)=(1au)^m`$ with $`a𝐂^\times `$ and $`m>1`$. From now on, we shall assume that we are in this case. To prove that $`W(𝝅)`$ is not isomorphic to $`V(𝝅)`$ as $`L(𝔤)`$-modules, recall that by Proposition 3.2, we have $`V(𝝅)V^{fin}(\lambda _𝝅)`$ as $`𝔤`$-modules. Hence, it suffices to prove that $`W(𝝅)`$ is reducible as a $`𝔤`$-module. Let $`𝝅_q`$ be an $`n`$-tuple of polynomials with constant term 1 such that $$\pi _i(u)=(1a_{i,1}u)(1a_{i,2}u)\mathrm{}(1a_{i,m_i}u),$$ where $`a_{i,j}=aq^{l_{ij}}`$, for some $`l_{ij}𝐙`$. Let $`v_{i,j}`$ be the highest weight vector in $`V_q(i,a_j)`$ and consider $$𝐕=V_q(1,a_{1,1})V_q(1,a_{1,2})\mathrm{}V(1,a_{1,m_1})\mathrm{}V_q(n,a_{n,1})\mathrm{}V_q(n,a_{n,m_n}).$$ Let $`𝐯=v_{1,1}v_{1,2}\mathrm{}v_{n,m_n}`$ and set $$Z_q(𝝅_q)=𝐔_q.𝐯𝐕,Z_𝐀(𝝅_q)=𝐔_𝐀.𝐯.$$ Since $`Z_q(𝝅_q)`$ is a quotient of $`W_q(𝝅_q)`$, and $`𝝅_q`$ is integral, it follows that $$Z_q(𝝅_q)Z_𝐀(𝝅_q)_𝐀𝐂(q),$$ so we can define the $`𝐔`$-module $`\overline{Z_q(𝝅_q)}=Z_𝐀(𝝅_q)_𝐀𝐂_1`$. Clearly, $`\overline{Z_q(𝝅_q)}`$ is a quotient of $`W(𝝅)`$ and hence it suffices to show that $`\overline{Z_q(𝝅_q)}`$ is reducible as a $`𝔤`$-module. Suppose first that $`m_{i_0}2`$ for some $`i_0I`$. Take $`l_{ij}=0`$ for all $`iI`$ and $`j=1,\mathrm{},m_i`$. Let $`𝐔_q^{i_0}`$ be the subalgebra of $`𝐔_q`$ generated by $`K_{i_0}^{\pm 1}`$ and $`x_{i_0,k}^\pm `$ for $`k𝐙`$. Consider the $`𝐔_q^{i_0}`$-module $$Z_q^{i_0}(𝝅_q)=𝐔_q^{i_0}.𝐯.$$ Let $`\omega `$ be the fundamental weight of $`sl_2`$. Then, by Proposition 5.1, we know that $`V_q(\omega ,a)^{m_{i_0}}`$ is irreducible and hence is a quotient of $`Z_q^{i_0}(𝝅_q)`$. Clearly, $$dim(Z_q(𝝅_q)_{\lambda \alpha _{i_0}})dim(V(\omega ,a)^{m_{i_0}})_{m_{i_0}\omega \alpha }=m_{i_0}$$ hence $$dim(\overline{Z_q(𝝅_q)}_{\lambda \alpha _{i_0}}m_{i_0}>1.$$ On the other hand, $`V(𝝅)`$ is a quotient $`𝐔`$-module of $`\overline{Z_q(𝝅_q)}`$, since $`𝝅_q=𝝅`$, and $`V(𝝅)V^{fin}(\lambda _𝝅)`$ as $`𝐔^{fin}`$-modules from above. But $$dim(V^{fin}(\lambda _𝝅)_{\lambda _𝝅\alpha _{i_0}})=1.$$ Hence, $`\overline{Z_q(𝝅_q)}`$ is reducible as a $`𝐔^{fin}`$-module. We can therefore assume that each $`m_i=0\text{or}1`$, and that at least one $`m_i=1`$. Consider first the case $`m_{i_0}=m_{i_1}=1`$ with $`i_0<i_1`$, and $`m_j=0`$ for all $`i_0<j<i_1`$. Set $`J=\{i_0,i_0+1,\mathrm{},i_1\}`$. By a suitable choice of the numbering, we can assume that the corresponding diagram subalgebra $`𝔤^J`$ of $`𝔤`$ is of type $`A_{\|J\|}`$. Let $`𝐔_q^J`$ be the subalgebra of $`𝐔_q`$ generated by $`K_i^{\pm 1}`$ and $`x_{i,k}^\pm `$ for $`k𝐙`$ and $`iJ`$. Define $$Z_q^J(𝝅_q)=𝐔_q^J.𝐯.$$ By Proposition 5.2, the $`𝐔_q^J`$-module $`V_q^J(i_0,a_{i_0})V_q^J(i_1,a_{i_1})`$ is irreducible except for finitely many values of the ratio $`a_{i_0}/a_{i_1}`$. Since each $`a_i`$ can be chosen from the infinite set $`\{aq^m:m𝐙\}`$, we can assume that $`a_{i_0}`$ and $`a_{i_1}`$ are chosen so that $`V_q^J(i_0,a_{i_0})V_q^J(i_1,a_{i_1})`$ is an irreducible $`𝐔_q^J`$-module and hence a quotient of $`Z_q^J(𝝅_q)`$. If $`\theta _J`$ is the highest root of the subdiagram $`J`$, then $`\text{dim}(Z_q(𝝅_q)_{\lambda _{𝝅_q}\theta _J})`$ $`\text{dim}(Z_q^J(𝝅_q)_{\lambda _{𝝅_q}\theta _J})`$ $`\text{dim}((V_q^J(i_0,a_{i_0})V_q^J(i_1,a_{i_1}))_{\lambda _{𝝅_q}\theta _J})`$ $`dim((V_q^{fin,J}(\omega _{i_0})V_q^{fin,J}(\omega _{i_1}))_0)=\|J\|+1`$ (in an obvious notation). On the other hand, $`V(𝝅)`$ is a quotient $`𝐔`$-module of $`\overline{Z_q(𝝅_q)}`$ and $$V(𝝅)_{\lambda _𝝅\theta _J}=V^J(𝝅)_{\lambda _𝝅\theta _J}=V^{fin,J}(\lambda _𝝅)_{\lambda _𝝅\theta _J},$$ which has dimension $`\|J\|`$. Hence, $`V(𝝅)`$ is a proper quotient of $`\overline{Z_q(𝝅_q)}`$. This shows that $`Z_q(𝝅_q)`$ is not isomorphic to $`V_q(\lambda _{𝝅_q})`$ and hence is reducible as a $`𝐔_q(𝔤)`$-module. It remains to consider the case when exactly one $`m_i=1`$, say $`m_{i_0}=1`$ and all other $`m_i=0`$. Since $`\pi _\theta `$ has repeated roots, this means that $`r_{i_0}1`$. By Proposition 5.2 we know that $`W(𝝅)=W(i,a)`$ is reducible. This completes the proof of the theorem.∎ ## 6. The $`sl_2`$ case. In this section, $`𝔤=sl_2`$. Let $`\omega `$ be the fundamental weight, $`\alpha `$ the positive root, and set $`x^\pm =x_{\pm \alpha }`$, $`h=h_\alpha `$. Let $`\pi `$ be a polynomial with coefficients in $`𝐂`$ and constant term 1. When $`\pi =1au`$, denote $`V_q(\pi )`$ by $`V_q(a)`$, and define $`V(a)`$ similarly. Note that these modules are two-dimensional over $`𝐂(q)`$ and $`𝐂`$, respectively. Set $`V=V^{fin}(\omega )`$ and let $`L(V)=V𝐂[t,t^1]`$ be the obvious $`𝐔^e`$-module, given by $$x_r.(vt^s)=x.vt^{r+s},d.(vt^r)=rvt^r,$$ for $`r,s𝐙`$, $`x𝔤`$ and $`vV`$. Set $`𝒫_𝓂=𝐂[𝓉_\mathcal{1}^{\pm \mathcal{1}},𝓉_\mathcal{2}^{\pm \mathcal{1}},\mathrm{},𝓉_𝓂^{\pm \mathcal{1}}]`$. Let $`\mathrm{\Sigma }_m`$ be the symmetric group on $`m`$ letters and let $`𝒫^{\mathcal{\Sigma }_𝓂}`$ be the subalgebra of polynomials invariant under the obvious action of $`\mathrm{\Sigma }_m`$. Let $`S^m(L(V))`$ be the symmetric part of the $`m`$–fold tensor product $`T^m(L(V))`$ of $`L(V)`$. Then, $`T^m(L(V))`$ is a $`𝐔^e`$–module in the obvious way, and $`S^m(L(V))`$ is a $`𝐔^e`$–submodule. Moreover, as vector spaces, $$T^m(L(V))V^m𝒫_𝓂,$$ and so $`T^m(L(V))`$ is a right module for $`𝒫`$ by right multiplication. This induces a right $`𝒫_𝓂^{\mathcal{\Sigma }_𝓂}`$-action on $`S^m(L(V))`$. The left $`𝐔^e`$–module $`W(m\omega )`$ is also a right $`𝒫^{\mathcal{\Sigma }_𝓂}`$–module. In fact, by equation (2), $`W(m\omega )`$ is a right $`𝐔(0)/I_{m\omega }(0)`$–module, i.e., a right module for the algebra $`𝐂[\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_m,\mathrm{\Lambda }_m^1]`$. But this algebra is isomorphic to $`𝒫^{\mathcal{\Sigma }_𝓂}`$ by taking $`\mathrm{\Lambda }_r`$ to the $`r^{th}`$ elementary symmetric function of $`t_1,\mathrm{},t_m`$. In this section, we shall prove the following two theorems. ###### Theorem 4. As left $`𝐔^e`$-modules and as right $`𝒫_𝓂^{\mathcal{\Sigma }_𝓂}`$-modules, we have $$W(m\omega )S^m(L(V)).$$ To prove Theorem 4 we shall need ###### Theorem 5. Let $`\pi (u)`$ be a polynomial with coefficients in $`𝐂`$ and constant term 1. Then, the dimension of $`W(\pi )`$ is $`2^{\text{deg}\pi }`$. In fact, $$W(\pi )\overline{V_q(\pi )},$$ as $`𝐔`$-modules, and $$V_q(\pi )\underset{a}{}V_q(a)$$ where $`a^1`$ runs over the set of roots of $`\pi `$ counted with multiplicity. It does not follow from this result that $`W(\pi )_aV(a)`$. In fact, this is false except when $`\text{deg}\pi =1`$. The point is that the $`𝐀`$-form of $`_aV_q(a)`$ is not the tensor product of the $`𝐀`$-forms of $`V_q(a)`$ (in fact, the former is a proper subset of the latter, unless $`\text{deg}\pi =1`$). We note the following corollary. ###### Corollary 6.1. For any $`\pi (u)`$ as in Theorem 5, we have $`W_q(\pi )V_q(\pi )`$ as $`𝐔_q`$–modules. ###### Proof. Since $`V_q(\pi )`$ is a quotient of $`W_q(\pi )`$ it suffices to prove that $$\text{dim}_{𝐂(q)}W_q(\pi )2^{\text{deg}\pi }.$$ But this is clear from Theorem 5, since $`\overline{W_q(\pi )}`$ is a quotient of $`W(\pi )`$, so $$\text{dim}_{𝐂(q)}W_q(\pi )=\text{dim}_𝐂\overline{W_q(\pi )}\text{dim}_𝐂W(\pi )=2^{\text{deg}\pi }.$$ Assume Theorem 5 for the moment. To prove Theorem 4, we begin with the following trivial lemma. ###### Lemma 6.1. The $`𝐔^e`$-module $`S^m(L(V))`$ is integrable. ∎ Let $`\{v_+,v_{}\}`$ be the usual basis of $`V`$, so that $$x^\pm .v_\pm =0,x^\pm .v_{}=v_\pm ,hv_\pm =\pm v_\pm .$$ For $`0rm`$, $`l_1,\mathrm{},l_m𝐙`$, define $$𝐯_{(l_1,l_2,\mathrm{},l_r),(l_{r+1},\mathrm{},l_m)}=\underset{\sigma \mathrm{\Sigma }_m}{}v_{\sigma (1)}t^{l_{\sigma (1)}}\mathrm{}v_{\sigma (m)}t^{l_{\sigma (m)}},$$ where we set $`v_s`$ $`=v_{},\text{if}1sr,`$ $`v_s`$ $`=v_+,\text{if}r+1sm.`$ Clearly the set $$\{𝐯_{(l_1,l_2,\mathrm{},l_r),(l_{r+1},\mathrm{},l_m)}:0rm,l_1,\mathrm{},l_m𝐙\},$$ is a $`𝐂`$–basis of $`S^m(L(V))`$. ###### Lemma 6.2. The $`𝒫_𝓂^{\mathcal{\Sigma }_𝓂}`$–module $`S^m(L(V))`$ is free of rank $`2^m`$. ###### Proof. For $`0rm`$, let $`S^m(L(V))_r`$ be the subspace of $`S^m(L(V))`$ spanned by the elements $$\{𝐯_{(l_1,l_2,\mathrm{},l_r),(l_{r+1},\mathrm{},l_m)}:l_1,\mathrm{},l_m𝐙\}.$$ Clearly, $`S^m(L(V))_r`$ is a right $`𝒫_𝓂^{\mathcal{\Sigma }_𝓂}`$–submodule of $`S^m(L(V))`$. It is easy to see that it is isomorphic to the $`𝒫_𝓂^{\mathcal{\Sigma }_𝓂}`$–module $`𝒫_𝓂^{\mathcal{\Sigma }_𝓇\times \mathcal{\Sigma }_{𝓂𝓇}}`$, consisting of the elements in $`𝒫_𝓂`$ invariant under permutation of the first $`r`$ and the last $`mr`$ variables.. But it is well known that the latter module is free of rank $`\left(\begin{array}{c}m\\ r\end{array}\right)`$. This proves the lemma. ∎ ###### Lemma 6.3. The assignment $`w_{m\omega }𝐯_+^m`$ extends to a well defined surjective homomorphism $`W(m\omega )S^m(L(V))`$ of left $`𝐔^e`$–modules and right $`𝒫_𝓂^{\mathcal{\Sigma }_𝓂}`$–modules. ###### Proof. It follows by Proposition 2.1(i) that there exists a $`𝐔^e`$-module homomorphism $`\varphi :W(m\omega )S^m(L(V))`$ that takes $`w_{m\omega }`$ to $`𝐯=v_+^m`$. It is trivial to check that $`\varphi `$ is also a map of right $`𝒫_𝓂^{\mathcal{\Sigma }_𝓂}`$–modules. To show that $`\varphi `$ is surjective, it is enough to prove that (6.1) $$S^m(L(V))=𝐔^e.𝐯.$$ We prove by induction on $`r`$ that (6.2) $$𝐯_{(l_1,l_2,\mathrm{},l_r),(l_{r+1},\mathrm{},l_m)}𝐔^e.𝐯$$ for all $`l_1,\mathrm{},l_m𝐙`$. Consider the case $`r=0`$. For any $`k_1,k_2,\mathrm{},k_m𝐙`$, we have $$(x_0^+)^mx_{k_1}^{}x_{k_2}^{}\mathrm{}x_{k_m}^{}.𝐯=\underset{\sigma \mathrm{\Sigma }_m}{}v_+t^{k_{\sigma (1)}}v_+t^{k_{\sigma (2)}}\mathrm{}v_+t^{k_{\sigma (m)}},$$ which proves (6.2) in this case. The case $`r=m`$ can be done similarly, since the element $`𝐯^{}=v_{}^m=\frac{1}{m!}(x^{})^m.𝐯𝐔^e`$. Assuming the result for $`r`$, we prove it for $`r+1`$. For this we shall proceed by an induction on $$N=\mathrm{\#}\{j:jr+1,l_j0\}.$$ Now, $$x_k^{}.𝐯_{(l_1,l_2,\mathrm{},l_r),(l_{r+1},\mathrm{},l_m)}=\underset{s=r+1}{\overset{m}{}}𝐯_{(l_1,l_2,\mathrm{},l_r,l_s+k),(l_{r+1},\mathrm{},\widehat{l}_s,\mathrm{},l_m)}.$$ Taking $`l_s=0`$ for all $`s>r`$, we get that $$𝐯_{(l_1,l_2,\mathrm{},l_r,k),(0,\mathrm{},0)}𝐔^e.𝐯,$$ for all $`k𝐙`$, proving our assertion when $`N=0`$. Assume the result for $`N1`$. We have to show that $$𝐯_{(l_1,l_2,\mathrm{},l_r,l_{r+1}),(l_{r+2},\mathrm{},l_{r+N+1},0,\mathrm{},0)}𝐔^e.𝐯.$$ Now $$x_{l_{r+1}}^{}.𝐯_{(l_1,l_2,\mathrm{},l_r),(l_{r+2},\mathrm{},l_{r+N+1},0,\mathrm{},0)}$$ is in $`𝐔^e.𝐯`$ by the induction hypothesis on $`r`$ , and is a sum of the term $$(mr)𝐯_{(l_1,\mathrm{},l_{r+1}),(l_{r+2},\mathrm{},l_{r+N+1},0,\mathrm{},0)}$$ and terms of type $$𝐯_{(l_1,\mathrm{},l_r,l_s+l_{r+1}),(l_{r+2},\mathrm{}\widehat{l}_s,\mathrm{},l_{r+N+1},0,\mathrm{},0)},$$ for $`r+2sr+N+1`$. Since, by the induction hypothesis on $`N`$, all the terms of the second type are in $`𝐔^e.𝐯`$, it follows that $$𝐯_{(l_1,\mathrm{},l_{r+1}),(l_{r+2},\mathrm{},l_{r+N+1},0,\mathrm{},0)}𝐔^e.𝐯.$$ This completes the proof that $`\varphi `$ is surjective.∎ Proof of Theorem 4. Let $`K`$ be the kernel of the homomorphism $`W(m\omega )S^m(L(V))`$ given by Lemma 6.3. Since $`S^m(L(V))`$ is a free, hence projective, right $`𝒫_𝓂^{\mathcal{\Sigma }_𝓂}`$–module by Lemma 6.2, it follows that $$W(m\omega )=S^m(L(V))K,$$ as right $`𝒫_𝓂^{\mathcal{\Sigma }_𝓂}`$–modules. Let $`𝔪`$ be any maximal ideal in $`𝒫_𝓂^{\mathcal{\Sigma }_𝓂}`$. Identifying $`𝒫_𝓂^{\mathcal{\Sigma }_𝓂}`$ with $`𝐔(0)/I_{m\omega }(0)`$ as described earlier in this section, it is clear that $$𝔪=I_\pi (0)/I_{m\omega }(0),$$ for some polynomial $`\pi `$ with constant term 1 such that $`\text{deg}\pi =m`$. It follows that $$W(m\omega )/W(m\omega )𝔪W(\pi )$$ as vector spaces over $`𝐂`$, and hence has dimension $`2^m`$. On the other hand, by Lemmma 6.2, $`S^m(L(V))/S^m(L(V))𝔪`$ also has dimension $`2^m`$ over $`𝐂`$. It follows that $`K/K𝔪=0`$. Since this holds for all maximal ideals $`𝔪`$, Nakayama’s lemma implies that $`K=0`$, proving the theorem. The rest of the section is devoted to proving Theorem 5. First, observe that, in view of Theorem 2, it suffices to consider the case when $`𝝅(u)=(1au)^m`$ for some $`a𝐂^\times `$. Since we have a surjective map $`W(𝝅)\overline{V_q(𝝅)}`$, it suffices to prove that (6.3) $$\text{dim}_𝐂W(𝝅)\text{dim}_𝐂\overline{V_q(\pi )}=2^m.$$ For $`a𝐂^\times `$, let $`\tau _a:𝔤𝐂[t]𝔤𝐂[t]`$ be the Lie algebra automorphism obtained by extending the assignment $$xt^kx(ta)^k,x𝔤,k0.$$ Set $$X_a^{}(u)=\tau _a(X_{\alpha ,0}^{}(u)),\mathrm{\Lambda }_a^+(u)=\tau _a(\mathrm{\Lambda }_\alpha ^+(u))=\text{exp}\left(\underset{k=1}{\overset{\mathrm{}}{}}\frac{h(ta)^k}{k}u^k\right).$$ It is easy to see, using the relation between the $`\mathrm{\Lambda }_m`$ and $`h_m`$, that (6.4) $$h_k.w_𝝅=ma^kw_𝝅,$$ or equivalently that $$h(ta)^k.w_𝝅=0,$$ for all $`k0`$. It follows that (6.5) $$(\mathrm{\Lambda }_a^+)_k.w_𝝅=0,k>0.$$ Further, using Lemma 1.3(ii) and observing that the identity there is actually an identity in $`𝐔(𝔤𝐂[t])`$, we get by applying $`\tau _a`$ that $$\tau _a(x_1^+)^{(r)}(x_0^{})^{(s)}=(1)^r\left(X_a^{}(u)^{(sr)}\mathrm{\Lambda }_a^+(u)\right)_smod\mathrm{𝐔𝐔}(>)_+,$$ for $`sr1`$. Together with (6.5), it follows that (6.6) $$\left(X_a^{}(u)^{(sr)}\right)_s.w_𝝅=0r1,sm+1.$$ In particular, this means that $$(x^{}(ta)^s).w_𝝅=0sm.$$ Let $`𝐔_a^+(<)`$ be the commutative subalgebra of $`𝐔`$ generated by the elements $`\tau _a(x_k^{})`$ for all $`k0`$, and let $`I_a(m)`$ be the ideal in $`𝐔_a^+(<)`$ generated by the elements $`\left(X_a^{}(u)^{(sr)}\right)_s`$ for all $`r1`$, $`sm+1`$. ###### Lemma 6.4. The assignment $`uu.w_𝛑`$ induces a surjective map of vector spaces $`𝐔_a^+(<)/I_a(m)W(𝛑)`$. ###### Proof. The map is well defined by (6.6). It is obviously surjective because the polynomials $`(ta)^k`$ for $`k0`$ are a basis of $`𝐂[t]`$.∎ Thus, to prove (6.3) it suffices to show that the dimension of $`𝐔_a^+(<)/I_a(m)`$ is at most $`2^m`$. It is convenient to reformulate the problem as follows. Let $`R_m=𝐂[z_0,\mathrm{},z_{m1}]`$ be the polynomial algebra in $`m`$ variables. For $`0j<m`$, set $$Z_j(u)=\underset{i=j}{\overset{m1}{}}z_iu^{ij+1}R_m[u].$$ Let $`J_m`$ be the ideal in $`R_m`$ generated by the elements $`\left(Z_0(u)^r\right)_s`$, for $`r1`$, $`sm+1`$. It is trivial to see that $$R_m/J_m𝐔_a^+(<)/I_a(m),$$ via the map $`\tau (x_k^{})z_k`$. It is clear that (6.3) is now a consequence of the following proposition and Lemma 6.4. ###### Proposition 6.1. For $`mr>0`$, let $$_{𝓂,𝓇}=\{𝓏_{𝒾_\mathcal{1}}𝓏_{𝒾_\mathcal{2}}\mathrm{}𝓏_{𝒾_𝓇}:\mathcal{0}𝒾_\mathcal{1}𝒾_\mathcal{2}\mathrm{}𝒾_𝓇𝓂𝓇\}.$$ Let $`_{𝓂,\mathcal{0}}=\{\mathcal{1}\}`$. The set $$_𝓂=\underset{𝓇=\mathcal{0}}{\overset{𝓂}{}}_{𝓂,𝓇}$$ spans $`R_m/J_m`$. We prove Proposition 6.1 by induction on $`m`$. The case $`m=1`$ is trivial, but for the inductive step, we need the following lemmas. Set $`J_m=J_{m,0}`$ and, for $`0<j<m`$, define ideals $`J_{m,j}`$ in $`R_m`$ inductively by $$J_{m,j}=J_{m,j1}+\underset{r=1}{\overset{j}{}}R_m(Z_1(u)^r)_{mj}=J_m+\underset{smj}{}\underset{1rms}{}R_m\left(Z_1(u)^r\right)_s.$$ ###### Lemma 6.5. If $`j1`$, there is a unique homomorphism of algebras $$R_{mj}/J_{mj}R_m/J_{m,j1}$$ such that $`z_iz_{i+1}`$ for $`0i<mj`$. ###### Proof. To establish the lemma, we must prove that (6.7) $$\left(\left(\underset{i=1}{\overset{mj}{}}z_iu^i\right)^r\right)_sJ_{m,j1},$$ for all $`r1`$ and $`smj+1`$. We proceed by induction on $`j`$. For $`j=1`$, we must show that $$\left(Z_1(u)^r\right)_sJ_msm.$$ If $`r=1`$, this is trivial from the definition of $`Z_1(u)`$. Assume the result for smaller values of $`r`$. Writing $$Z_0(u)=u(z_0+Z_1(u)),$$ we have $$\left(Z_0(u)^r\right)_s=\left(\underset{t=0}{\overset{r}{}}\left(\begin{array}{c}r\\ t\end{array}\right)z_0^t(Z_1(u)^{rt})_{sr}\right).$$ Take $`s=n+r`$ with $`nm`$. Then, $`\left(Z_0(u)^r\right)_{n+r}J_m`$ by the definition of $`J_m`$, so $$\underset{t=0}{\overset{r}{}}\left(\begin{array}{c}r\\ t\end{array}\right)z_0^t(Z_1(u)^{rt})_nJ_m.$$ But, if $`t>0`$, then $`(Z_1(u)^{rt})_nJ_m`$ for all $`nm`$ by the induction hypothesis on $`r`$, so $`(Z_1(u)^r)_nJ_m`$ for all $`nm`$, thus completing the induction on $`r`$, and establishing (6.7) when $`j=1`$. So now assume that we know (6.7) for $`j1`$. Write $$\underset{i=1}{\overset{mj+1}{}}z_iu^i=z_{mj+1}u^{mj+1}+\underset{i=1}{\overset{mj}{}}z_iu^i.$$ By the induction hypothesis on $`j`$, $$\left(\left(\underset{i=1}{\overset{mj+1}{}}z_iu^i\right)^r\right)_sJ_{m,j2}$$ for all $`r1`$, $`smj+2`$. Since $`z_{mj+1}J_{m,j1}`$, we can conclude, by using the binomial expansion, that $$\left(\left(\underset{i=1}{\overset{mj}{}}z_iu^i\right)^r\right)_sJ_{m,j1},$$ if $`r1`$, $`s>mj+1`$. Thus, it suffices to prove that (6.8) $$\left(\left(\underset{i=1}{\overset{mj}{}}z_iu^i\right)^r\right)_{mj+1}J_{m,j1}$$ for all $`r1`$. If $`rj1`$, we have $`(Z_1(u)^r)_{mj+1}J_{m,j1}`$ by definition. Further, the elements $`z_{mj+1},\mathrm{},z_{m1}J_{m,j1}`$. Thus, writing $$Z_1(u)=\underset{i=1}{\overset{mj}{}}z_iu^i+\underset{i=mj+1}{\overset{m1}{}}z_iu^i,$$ and using the binomial expansion, we see that (6.8) follows. If $`r>j1`$, then $`r+mj+1>m`$, so $$\left(Z_0(u)^r\right)_{r+mj+1}J_mJ_{m,j1}$$ by the definition of $`J_m`$, we have $$\left((z_0+Z_1(u))^r\right)_{mj+1}J_{m,j1}.$$ Now, $`Z_1(u)_{mj+1}=z_{mj+1}J_{m,j1}`$ by definition, and so using the binomial expansion again and an induction on $`r`$, we conclude that $$\left(Z_1(u)^r\right)_{mj+1}J_{m,j1}.$$ But now the proof is completed as in the case $`r<j`$.∎ The proof of the following lemma is elementary. ###### Lemma 6.6. Let $`rk0`$. Then, the matrix $$\left[\begin{array}{cccc}\left(\begin{array}{c}k+1\\ 1\end{array}\right)& \left(\begin{array}{c}k+1\\ 2\end{array}\right)& \mathrm{}& \left(\begin{array}{c}k+1\\ rk+1\end{array}\right)\\ & & & & \\ \left(\begin{array}{c}k+2\\ 1\end{array}\right)& \left(\begin{array}{c}k+2\\ 2\end{array}\right)& \mathrm{}& \left(\begin{array}{c}k+2\\ rk+1\end{array}\right)\\ .& .& \mathrm{}& .\\ .& .& \mathrm{}& .\\ .& .& \mathrm{}& .\\ .& .& \mathrm{}& .\\ \left(\begin{array}{c}r+1\\ 1\end{array}\right)& \left(\begin{array}{c}r+1\\ 2\end{array}\right)& \mathrm{}& \left(\begin{array}{c}r+1\\ rk+1\end{array}\right)\end{array}\right]$$ has determinant $`\left(\begin{array}{c}r+1\\ k\end{array}\right)`$ and hence is invertible. ∎ ###### Lemma 6.7. For $`rt>0`$, $`0sr`$, the element $`z_0^{rs}(Z_1(u)^{s+1})_{mt}`$ belongs to the span of $$\{z_0^{tj}(Z_1(u)^{rt+j+1})_{mt}:1jt\}.$$ ###### Proof. We assume that $`m>1`$, otherwise there is nothing to prove. We consider the following equations in $`R_m/J_m`$: $$z_0^{rj}(Z_0(u)^{j+1})_{m+j+1t}=0,0<tj,$$ i.e., $$z_0^{rj}((z_0+Z_1(u))^{j+1})_{mt}=0,0<tj,$$ i.e., $$\underset{i=0}{\overset{j}{}}\left(\begin{array}{c}j+1\\ i\end{array}\right)z_0^{rj+i}(Z_1(u)^{j+1i})_{mt}=0.$$ We must show that these equations, for $`j=t,t+1,\mathrm{},r`$, can be solved for the elements $`z_0^s(Z_1(u)^{r+1s})_{mt}`$ with $`tsr`$ in terms of those with $`s<t`$. But this follows from the preceding lemma. ∎ Proof of Proposition 6.1. The proposition is trivially true if $`m=1`$. Assume now that we know the result for $`m1`$. For $`0k<rm`$, set $$_{𝓂,𝓇,𝓀}=\{𝓏_{𝒾_\mathcal{1}}𝓏_{𝒾_\mathcal{2}}\mathrm{}𝓏_{𝒾_𝓇}_{𝓂,𝓇}:𝒾_{𝓀+\mathcal{1}}\mathcal{1}\}.$$ The proposition obviously follows from Claim Let $`rk0`$, and let $`gR_m`$ be a homogenous polynomial of degree $`rk`$ in $`z_1,z_2,\mathrm{},z_{m1}`$. Then, $`z_0^kg`$ is in the span of $`_{𝓂,𝓇,𝓀}`$ modulo $`J_m`$. We proceed by induction on $`k`$. If $`k=0`$, then by Lemma 6.5 we have a homomorphism $`R_{m1}/J_{m1}R_m/J_m`$ which sends $`z_iz_{i+1}`$. Clearly, $`g`$ is in the image of this homomorphism and the induction hypothesis on $`m`$ implies that $`g_{𝓂,𝓇,\mathcal{0}}`$. Assume the result for $`k1`$. Write $$g=g_0+g_1z_{mk}+g_2z_{mk+1}+\mathrm{}+g_kz_{m1},$$ where for $`0jk`$, $`g_j`$ is a polynomial in $`z_1,z_2,\mathrm{},z_{mk+j1}`$. Now, for $`j0`$, we see by Lemma 6.7 that the element $`z_0^kz_{mk+j}`$ is in the span of the sets $`_{𝓂,𝓀+\mathcal{1},𝓈}`$ with $`s<k`$. Thus, the element $`z_0^kz_{mk+j}g_{j+1}`$ can be written as a sum $`_{s<k}z_0^sh_{sj}`$, where the $`h_{sj}`$ are polynomials in $`z_1,\mathrm{},z_{m1}`$. Hence, by the induction on $`k`$, $$z_0^kz_{mk+j}g_{j+1}_{𝓂,𝓇,𝓀},$$ for $`j0`$. Finally observe that by Lemma 6.5, $`g_0`$ is in the image of the map $`R_{mk1}/J_{mk+j1}R_m/J_m`$ and hence, by using the induction on $`m`$, we get that $$g_0\text{span}(_{𝓂,𝓇𝓀,\mathcal{0}})mod𝒥_{𝓂,𝓀+\mathcal{1}}.$$ Thus, $`z_0^kg_0`$ is in the span of $`_{𝓂,𝓇,𝓀}`$ provided that $`z_0^kJ_{m,k}`$ is also in the span of $`_{𝓂,𝓇,𝓀}`$, i.e if $`z_0^k(Z_1(u)^r)_{mk1}`$ is in the span of $`_{𝓂,𝓇,𝓀}`$ for all $`smk`$, $`1rms`$. But this follows from Lemma 6.7 again, and the proof of the proposition is complete. ∎
warning/0004/math0004104.html	ar5iv	text	# Invariant Subspaces of Voiculescu’s Circular Operator ## 1. Introduction The invariant subspace problem relative to a von Neumann algebra $`McB()`$ asks whether every operator $`TMc`$ has a proper, nontrivial invariant subspace $`_0`$ such that the orthogonal projection $`p`$ onto $`_0`$ is an element of $`Mc`$; equivalently, it asks whether there is a projection $`pMc`$, $`p\{0,1\}`$, such that $`Tp=pTp`$. Even when $`Mc`$ is a II<sub>1</sub>–factor, this invariant subspace problem remains open. In this paper we show that the circular operator and each circular free Poisson operator (defined below) has a continuous family of invariant subspaces relative to the von Neumann algebra it generates. These operators arise naturally in free probability theory, (see the book ), and each generates the von Neumann algebra II<sub>1</sub>–factor $`L(F_2)`$ associated to the nonabelian free group on two generators. Given a von Neumann algebra $`Mc`$ with normal faithful state $`\varphi `$, a circular operator is $`y=(x_1+ix_2)/\sqrt{2}Mc`$, where $`x_1`$ and $`x_2`$ are centered semicircular elements having the same second moments and that are free with respect to $`\varphi `$. For specificity, we will always take cicular elements to have the normalization $`\varphi (y^{}y)=1`$, which is equivalent to $`\varphi (x_i^2)=1`$. Voiculescu found a matrix model for a circular element, showing that if $`X(n)`$ is a random matrix whose entries are i.i.d. complex $`(0,1/n)`$–Gaussian random variables then $`X(n)`$ converges in $``$–moments as $`n\mathrm{}`$ to a circular element, meaning that $$\underset{n\mathrm{}}{lim}\tau _n(X(n)^{ϵ(1)}X(n)^{ϵ(2)}\mathrm{}X(n)^{ϵ(k)})=\varphi (y^{ϵ(1)}y^{ϵ(2)}\mathrm{}y^{ϵ(k)})$$ for every $`k𝐍`$ and for every choice of $`ϵ(j)`$ being “$``$” or no symbol, where $`\tau _n`$ is the expectation of the normalized trace and where $`y`$ is a circular element. Using the matrix model, Voiculescu showed that if $`(y_{ij})_{1i,jN}`$ is a $``$–free family of circular elements in a von Neumann algebra $`Mc`$ with respect to a normal faithful state $`\varphi `$, then the matrix $`y=\frac{1}{\sqrt{N}}(y_{ij})_{1i,jN}M_N(Mc)`$ is circular with respect to the state $`\varphi _N`$ given by $`\varphi _N\left((x_{ij})_{1i,jN}\right)=\frac{1}{N}_{i=1}^N\varphi (x_{ii})`$. Furthermore, he showed that the polar decomposition of a circular operator $`y`$ is $`y=ub`$ where $`u`$ is a Haar unitary (i.e. a unitary satisfying $`\varphi (u^k)=0`$ for every integer $`k>0`$), where $`b`$ is a quarter–circular element, (i.e. having moments $`\varphi (b^k)=\frac{1}{\pi }_0^2t^k\sqrt{4t^2}`$) and where $`u`$ and $`b`$ are $``$–free. These and results of a similar nature have been instrumental in applications of free probability to the study of the free group factors $`L(F_n)`$ and related factors; some of the first of these were , , , , . Voiculescu’s matrix model for the circular element is the starting point for finding invariant subspaces. Combined with a result of Dyson, it leads to upper triangular matrix models for the circular operator, namely, a sequence $`Y(n)`$ of upper triangluar random matrices whose $``$–moments converge to those of a circular operator. In these models, the elements of $`Y(n)`$ that are above the diagonal are complex $`(0,1/n)`$–Gaussian random variables, and we show that a number of different choices are possible for the diagonal entries. From these matrix models we show for any $`N2`$ that a circular operator can be realized as an $`N\times N`$ upper triangular matrix $$\frac{1}{\sqrt{N}}\left(\begin{array}{cccccc}a_1& b_{12}& b_{13}& \mathrm{}& b_{1,N1}& b_{1N}\\ 0& a_2& b_{23}& \mathrm{}& b_{2,N1}& b_{2N}\\ 0& 0& a_3& \mathrm{}& \mathrm{}& b_{3N}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& \mathrm{}& 0& a_{N1}& b_{N1,N}\\ 0& 0& \mathrm{}& 0& 0& a_N\end{array}\right)$$ (1) with entries in some W–noncommutative probability space, where the collection of $`N(N+1)/2`$ nonzero entries is $``$–free, where the entries $`b_{ij}`$ lying strictly above the diagonal are circular elements and where the entries $`a_j`$ on the diagonal are circular free Poisson elements, ($`a_j`$ being circular free Poisson of parameter $`j`$). These latter are generalizations of the circular operator (in the family of R–diagonal elements introduced by Nica and Speicher ) which are quite natural from the perspective of free probability theory. ###### Definition 1.1. Let $`(A,\psi )`$ be a W–noncommutative probability space with $`\psi `$ faithful and let $`c1`$. A circular free Poisson element of parameter $`c`$ in $`(A,\psi )`$ is an element of the form $`UH_c`$ where $`U,H_cA`$, $`U`$ is a Haar unitary, $`H_c0`$, the pair $`\{U,H_c\}`$ is $``$–free and $`H_c^2`$ has moments equal to those of a free Poisson distribution<sup>1</sup><sup>1</sup>1We must point out that the formula found on \[17, p. 35\] for the free Poisson distribution has some errors. The formula for $`_\mu `$ found there is correct, but the formulae for $`G_\mu (z)`$ and for the density are incorrect. with parameter $`c`$. Thus, letting $`a=(1\sqrt{c})^2`$ and $`b=(1+\sqrt{c})^2`$, the moments of $`H_c^2`$ are equal to those of the measure $`\nu _c`$ that is supported on $`[a,b]`$, is absolutely continuous with respect to Lebesgue measure and has density $$\frac{\text{d}\nu _c}{\text{d}\lambda }(t)=\frac{\sqrt{(bt)(ta)}}{2\pi t}1_{[a,b]}(t).$$ We hasten to point out that a circular element $`z`$ with normalization $`\psi (z^{}z)=1`$ is nothing other than a circular free Poisson element of parameter $`c=1`$. The spectrum of a circular free Poission element of parameter $`c`$ has been found by Haagerup and Larsen to be the annulus centered at zero with radii $`\sqrt{c1}`$ and $`\sqrt{c}`$, if $`c>1`$, whereas the spectrum of the circular operator is the disk of radius $`1`$. In the realization (1) of the circular operator, we have that the diagonal entry $`a_j`$ is circular free Poisson of parameter $`j`$. Therefore, the spectrum of the diagonal entry $`a_j`$ increases in modulus as $`j`$ increases, and the spectra of $`a_j`$ and $`a_k`$ overlap only if $`\|jk\|1`$. These properties of the realization (1) allow general techniques for upper triangular operators to be applied in order to find invariant subspaces of the circular operator $`y`$. It turns out that for every $`0<r<1`$ there is a unique projection $`pMc=\{y\}^{\prime \prime }`$ such that 1. $`yp=pyp`$ 2. $`\sigma (yp)=\{z𝐂\|z\|r\}`$ 3. $`\sigma ((1p)y)=\{z𝐂r\|z\|1\}`$ where in (ii) (respectively (iii)), the spectrum is computed relative to the algebra $`pMcp`$, (respectively $`(1p)Mc(1p)`$). In fact, the techniques outlined above can be employed with very little extra effort to find invariant subspaces for every circular free Poisson operator, and the proof is presented in this generality throughout. In §2, some theory is developed proving the existence of invariant subspaces, relative to the generated von Neumann algebras, of upper triangular operators, the spectra of whose diagonal entries satisfy certain conditions. In §3, we consider upper triangular random matrices whose entries strictly above the diagonal are i.i.d. complex Gaussian random variables which are independent of the diagonal entries. The general theme of the results in §3 is that the diagonal entries may be changed in certain ways but that as matrix size tends to infinity, the limit $``$–moments remain the same. In §4, we generalize asymptotic freeness results which Voiculescu originally proved for Gaussian random matrices and constant diagonal matrices; we allow the diagonal matrices to be random, subject to certain conditions. In §5, the random matrix results of the previous two sections together with results of Dyson and others are used to find various upper triangular matrix models for circular free Poisson elements, and these are in turn used to find an upper triangular realization of the same, as in (1). In §6, this upper triangular realization of the circular free Poisson element is fed into the machinery of §2 to find invariant subspaces. ## 2. An invariant subspace for an upper triangular operator Suppose $``$ is a Hilbert space and $`T:`$ is a bounded operator. In this section we are concerned with invariant subspaces $`_0`$ for $`T`$ such such that the projection from $``$ onto $`_0`$ lies in the von Neumann algebra generated by $`T`$. It is easy to see (Lemma 2.1) that for every $`r0`$ the set $$_r(T)=\overline{\{\xi \underset{k\mathrm{}}{lim\; sup}T^k\xi ^{1/k}r\}}$$ (2) is such an invariant subspace $`_0`$. The question is, for any given operator, whether these subspaces can be other than $`\{0\}`$ or $``$. We will show (Proposition 2.2) that the answer is yes if $`T`$ can be written as an upper triangular operator, $$T=\left(\begin{array}{ccc}& & \\ 0& & \\ 0& 0& \end{array}\right)$$ with respect to a decomposition $`=_1_2_3`$ under a condition on the spectra of the elements in the upper left–hand and lower right–hand corners of the above matrix. ###### Lemma and Definition 2.1. Let $`T:`$ be a bounded operator on a Hilbert space $``$, let $`r0`$ and let $`_r=_r(T)`$ be the subspace (2) defined above. Then $`_r`$ is an invariant subspace for $`T`$ such that the orthogonal projection $`p_r=p_r(T)`$ from $``$ onto $`_r`$ lies in the von Neumann algebra $`\{T\}^{\prime \prime }`$ generated by $`T`$. ###### Proof. Consider the subset $$E_r=E_r(T)=\{\xi \underset{k\mathrm{}}{lim\; sup}T^k\xi ^{1/k}r\}$$ (3) which is dense in $`_r`$. To see that $`_r`$ is a closed subspace of $``$, it will suffice to show that $`E_r`$ is a subspace. Let $`a𝐂`$ and $`\xi _1,\xi _2E_r`$. Since $`T^k(a\xi _1)=\|a\|T^k\xi _1`$ and $`T^k(\xi _1+\xi _2)2\mathrm{max}(T^k\xi _1,T^k\xi _2)`$, it is clear that $`a\xi _1E_r`$ and $`\xi _1+\xi _2E_r`$. Moreover, since $`T^k(T\xi )TT^k(\xi )`$ it is clear that $`E_r`$, and hence also $`_r`$, is invariant for $`T`$. To show that $`p_r\{T\}^{\prime \prime }`$ it will be enough to show that $`Up_r=p_rU`$ whenever $`U`$ is a unitary operator on $``$ such that $`UT=TU`$. Moreover, $`Up_r=p_rU`$ will follow once we show that $`U\xi E_r`$ for every $`\xi E_r`$. But this holds because $$T^kU\xi =UT^k\xi =T^k\xi .$$ ###### Proposition 2.2. Let $`T:`$ be a bounded operator on a Hilbert space $``$. Suppose that $`e_1,e_2,e_3`$ are orthogonal projections in $``$ with $`e_1+e_2+e_3=1`$ and that $`e_1`$ and $`e_1+e_2`$ are invariant for $`T`$. This means that $`T`$ is upper triangular with respect to this decomposition of $``$: $$T=\left(\begin{array}{ccc}e_1Te_1& & \\ 0& e_2Te_2& \\ 0& 0& e_3Te_3\end{array}\right).$$ Let $`r0`$ and suppose that $$sup\{\|\lambda \|\lambda \sigma (e_1Te_1)\}r<inf\{\|\lambda \|\lambda \sigma (e_3Te_3)\},$$ where $`\sigma (e_jTe_j)`$ denotes the spectrum of $`e_jTe_j`$ acting on $`e_j`$. Then the invariant subspace $`_r(T)`$ and its projection $`p_r=p_r(T)`$ defined in Lemma and Definition 2.1 satisfy $$e_1p_re_1+e_2.$$ ###### Proof. If $`\xi =e_1\xi `$ then $`T^k\xi =(e_1Te_1)^k\xi `$ so $`T^k\xi ^{1/k}\left((e_1Te_1)^k\xi \right)^{1/k}`$ while $`lim_k\mathrm{}(e_1Te_1)^k^{1/k}r`$. This shows $`e_1p_r`$. Suppose $`\xi `$ and $`e_3\xi 0`$. Then $`e_3T^k\xi =(e_3Te_3)^k\xi `$. As an operator on $`e_3`$, $`e_3Te_3`$ is invertible and its inverse has spectral radius $`<r^1`$. Therefore $`T^k\xi (e_3Te_3)^k\xi (e_3Te_3)^k^1e_3\xi `$ and hence $`lim\; sup_k\mathrm{}T^k\xi ^{1/k}lim_k\mathrm{}(e_3Te_3)^k^{1/k}>r`$, so $`\xi E_r(T)`$. This shows that $`E_r(T)e_3`$, and therefore that $`p_re_1+e_2`$. ∎ Invariant projections $`p_r(T)`$ and the dense subspaces $`E_r(T)`$ for some specific operators $`T`$ are described in §6. Now we show that for an element $`x`$ of an abstract W–algebra, the projection $`p_r(x)`$ is defined independently of how the W–algebra is represented as a von Neumann algebra acting on a Hilbert space. ###### Lemma 2.3. Let $``$ and $`^{}`$ be Hilbert spaces, let $`TB()`$ and take $`r>0`$. Then $$p_r(T)1_{^{}}=p_r(T1_{^{}}).$$ (4) ###### Proof. Let $`E_r(T)`$ be given by (3) and let $$E_r(T1)=\{w^{}\underset{k\mathrm{}}{lim\; sup}(T^k1_{^{}})w^{1/k}r\}.$$ Letting $``$ denote the algebraic tensor product of vector spaces, we clearly have $`E_r(T)^{}E_r(T1)`$, so the inequality $``$ holds in (4). Given a unit vector $`\eta ^{}`$ let $`V_\eta :^{}`$ be $`V_\eta (\xi )=\xi \eta `$. Then $`TV_\eta ^{}=V_\eta ^{}(T1_{^{}})`$, so if $`wE_r(T1)`$ then $`V_\eta ^{}wE_r(T)`$ for every $`\eta `$. Therefore $`E_r(T1)_r(T)^{}`$ and $``$ holds in (4). ∎ ###### Lemma and Definition 2.4. If $`Mc`$ is a von Neumann algebra, if $``$ is a Hilbert space and if $`\pi :McB()`$ is a normal, faithful $``$–representation then given $`xMc`$ we have, for $`r0`$, the projection $$p_r(\pi (x))\{\pi (x)\}^{\prime \prime }\pi (Mc).$$ Let us denote by $`p_r(x)Mc`$ the element so that $`\pi (p_r(x))=p_r(\pi (x))`$. Then $`p_r(x)`$ is independent of the choice of $``$ and $`\pi `$. ###### Proof. Let $`^{}`$ be a Hilbert space and $`\pi ^{}:McB(^{})`$ a normal, faithful $``$–representation. Using \[3, Theorem 1.4.3\], one finds a Hilbert space $`^{\prime \prime }`$ such that the representations $`\pi 1_{^{\prime \prime }}`$ and $`\pi ^{}1_{^{\prime \prime }}`$ are unitarily equivalent, via a unitary $`U:^{\prime \prime }^{}^{\prime \prime }`$. Now applying Lemma 2.3 twice, we have $`U^{}\left(\pi ^{}(p_r(x))1\right)U`$ $`=\pi (p_r(x))1=p_r(\pi (x))1=p_r(\pi (x)1)=`$ $`=p_r\left(U^{}(\pi ^{}(x)1)U\right)=U^{}p_r\left(\pi ^{}(x)1\right)U=U^{}\left(p_r(\pi ^{}(x))1\right)U.`$ Hence $`\pi ^{}(p_r(x))=p_r(\pi ^{}(x))`$. ∎ ## 3. Upper triangular random matrices In this section, we consider upper triangular random matrices whose entries strictly above the diagonal are i.i.d. Gaussian, and we prove that the diagonal entries can be modified in various ways without affecting the limiting $``$–moments as matrix size increases without bound. Throughout this paper, we consider random matrices whose entries have moments of all orders. Thus, let $`(\mathrm{\Omega },\mu )`$ be a usual probability space, let $`=_{1p<\mathrm{}}L^p(\mu )`$ and consider the expectation $`E(f)=f𝑑\mu `$. If $`S_1,S_2`$ are sets of random variables, we say that $`S_1`$ and $`S_2`$ are independent sets if the two $`\sigma `$–algebras generated by $`\{f^1(A)fS_i,A\text{ Borel subset of }𝐂\}`$ ($`i=1,2`$) are independent with respect to $`\mu `$, and similarly for families of sets of random variables. The $``$–algebra of $`n\times n`$ random matrices is $`_n=M_n(𝐂)`$ and has the trace $`\tau _n=E\mathrm{tr}_n`$, where $`\mathrm{tr}_n`$ is the trace on $`M_n(𝐂)`$ normalized so that $`\mathrm{tr}_n(1)=1`$. We fix a system of matrix units in $`M_n(𝐂)`$, denoted by $`(e(i,j;n))_{1i,jn}`$. ###### Notation 3.1. Let $`\sigma ^2>0`$ and $`n𝐍`$. 1. On $`𝐂`$, by Lebesgue measure we shall mean $`\text{d}(\mathrm{Re}z)\text{d}(\mathrm{Im}z)`$, i.e. normalized so that the unit disk has measure $`\pi `$. On the space $`M_n(𝐂)`$ of $`n\times n`$ complex matrices, Lebesgue measure shall mean the product of Lebesgue measure on each of the $`n^2`$ complex entries. On the space $`M_n^{s.a.}`$ of self–adjoint complex $`n\times n`$ matrices, Lebesgue measure shall mean the product of Lebesgue measure on each of the $`n(n1)/2`$ complex entries strictly above the diagonal and Lebesgue measure on each of the $`n`$ real diagonal entries. 2. We say that a random variable $`a`$ is a complex $`(0,\sigma ^2)`$–Gaussian if $`\mathrm{Re}a`$ and $`\mathrm{Im}a`$ are independent real Gaussian random variables each having first moment $`0`$ and second moment $`\sigma ^2/2`$. Thus $`E(a)=0`$, $`E(\|a\|^2)=\sigma ^2`$ and $`a`$ has density $`(\pi \sigma ^2)^1e^{\|z\|^2/\sigma ^2}`$ with respect to Lebesgue measure. 3. Given $`T_n`$, we will write $`T\text{UTGRM}(n,\sigma ^2)`$, (the acronym is for“upper triangular Gaussian random matrix”) if the entries $`t_{ij}`$ of $`T`$ ($`1i,jn`$) satisfy that $`t_{ij}=0`$ whenever $`ij`$ and that $`(t_{ij})_{1i<jn}`$ is an independent family of random variables, each of which is complex $`(0,\sigma ^2)`$–Gaussian. Our first result is that if $`T(n)\text{UTGRM}(n,1/n)`$, if $`D(n)_n`$ is a diagonal random matrix that is independent from $`T(n)`$ and if the joint distribution of the diagonal entries of $`D(n)`$ is permutation invariant then in the limit as $`n\mathrm{}`$, the $``$–moments of $`T(n)+D(n)`$ depend only on the marginal $``$–distributions of finite sets of the diagonal entries of $`D(n)`$. ###### Theorem 3.2. For every $`n𝐍`$ let $`T(n)\text{UTGRM}(n,\frac{1}{n})`$ and let $`D_1(n)`$ $`={\displaystyle \underset{i=1}{\overset{n}{}}}a(i;n)e(i,i;n)_n`$ $`D_2(n)`$ $`={\displaystyle \underset{i=1}{\overset{n}{}}}b(i;n)e(i,i;n)_n`$ be diagonal random matrices such that $`T(n)`$ and $`D_1(n)`$ are independent matrix–valued random variables and $`T(n)`$ and $`D_2(n)`$ are independent matrix–valued random variables. Let $`d\nu _\iota (\lambda _1,\mathrm{},\lambda _n)`$ be the joint distribution of the diagonal entries of $`D_\iota `$. Assume that $`d\nu _\iota `$ is invariant under all permutations of its $`n`$ variables, ($`\iota =1,2`$). Suppose that the marginal $``$–distributions of the diagonal entries of $`D_1(n)`$ are arbitrarily close to those of $`D_2(n)`$ as $`n\mathrm{}`$; namely suppose that $$\begin{array}{cc}\hfill p𝐍r_1,s_1,\mathrm{},r_p,& s_p𝐍\{0\},\hfill \\ \hfill \underset{n\mathrm{}}{lim}(& E(a(1;n)^{r_1}\overline{a(1;n)}^{s_1}\mathrm{}a(p;n)^{r_p}\overline{a(p;n)}^{s_p})\hfill \\ & E(b(1;n)^{r_1}\overline{b(1;n)}^{s_1}\mathrm{}b(p;n)^{r_p}\overline{b(p;n)}^{s_p}))=0.\hfill \end{array}$$ (5) Let $`Z_\iota (n)=D_\iota (n)+T(n)`$. Then $`m𝐍`$ $`ϵ(1),\mathrm{},ϵ(m)\{,\},`$ $`\underset{n\mathrm{}}{lim}\left(\tau _n(Z_1(n)^{ϵ(1)}\mathrm{}Z_1(n)^{ϵ(m)})\tau _n(Z_2(n)^{ϵ(1)}\mathrm{}Z_2(n)^{ϵ(m)})\right)=0,`$ where we take $`ϵ(j)=`$ to mean $`ϵ(j)`$ is “not $``$” or “no symbol.” Therefore if $`Z_1(n)`$ converges in $``$–moments as $`n\mathrm{}`$ then so does $`Z_2(n)`$ and their limit $``$–moments coincide. ###### Proof. Write $$T(n)=\underset{1i<jn}{}t(i,j;n)e(i,j;n).$$ Let us first fix $`n𝐍`$, $`\iota \{1,2\}`$ and let us denote $`D_\iota (n)`$ simply by $`D`$, $`T(n)`$ by $`T`$, $`Z_\iota (n)`$ by $`Z`$ and the diagonal entries of $`D`$ by $`d(i;n)`$, ($`1in`$). Now each word in $`Z`$ and $`Z^{}`$ is a sum of words in $`D`$, $`D^{}`$, $`T`$ and $`T^{}`$; hence we will investigate words of the form $$\begin{array}{c}W=\left(T^{ϵ(1)}\mathrm{}T^{ϵ(l(1))}\right)D^{\kappa (1)}\left(T^{ϵ(l(1)+1)}\mathrm{}T^{ϵ(l(2))}\right)D^{\kappa (2)}\mathrm{}\hfill \\ \hfill \mathrm{}\left(T^{ϵ(l(p1)+1)}\mathrm{}T^{ϵ(l(p))}\right)D^{\kappa (p)}\left(T^{ϵ(l(p)+1)}\mathrm{}T^{ϵ(q)}\right)\end{array}$$ (6) for arbitrary $$p,q𝐍\{0\},$$ $$ϵ(1),\mathrm{},ϵ(q),\kappa (1),\mathrm{},\kappa (p)\{,\},$$ $$0l(1)l(2)\mathrm{}l(p)q.$$ Now, using the independence of T and D, we see that $$\begin{array}{c}\tau _n(W)=\hfill \\ \hfill =n^1\underset{i_1,\mathrm{},i_q\{1,\mathrm{},n\}}{}\left(\underset{j=1}{\overset{q}{}}G^{ϵ(j)}(i_j,i_{j+1})\right)E\left(t^{ϵ(1)}(i_1,i_2;n)t^{ϵ(2)}(i_2,i_3;n)\mathrm{}t^{ϵ(q)}(i_q,i_{q+1};n)\right)\\ \hfill E\left(d^{\kappa (1)}(i_{l(1)};n)d^{\kappa (2)}(i_{l(2)};n)\mathrm{}d^{\kappa (p)}(i_{l(p)};n)\right),\end{array}$$ (7) where we have used the convention $`i_{q+1}=i_1`$, where $`G^ϵ(,)`$ is defined by $`G(i_j,i_{j+1})`$ $`=\{\begin{array}{cc}1\hfill & \text{if }i_j<i_{j+1}\hfill \\ 0\hfill & \text{otherwise}\hfill \end{array}`$ $`G^{}(i_j,i_{j+1})`$ $`=\{\begin{array}{cc}1\hfill & \text{if }i_j>i_{j+1}\hfill \\ 0\hfill & \text{otherwise}\hfill \end{array}`$ and where $`t^{ϵ(j)}(i_j,i_{j+1};n)`$ $`=\{\begin{array}{cc}t(i_j,i_{j+1};n)\hfill & \text{if }ϵ(j)=\hfill \\ [1.5ex]\overline{t(i_{j+1},i_j;n)}\hfill & \text{if }ϵ(j)=\hfill \end{array}`$ $`d^{\kappa (j)}(i_j;n)`$ $`=\{\begin{array}{cc}d(i_j;n)\hfill & \text{if }\kappa (j)=\hfill \\ [1.5ex]\overline{d(i_j;n)}\hfill & \text{if }\kappa (j)=\hfill \end{array}`$ Using that the joint distribution of $`d(1;n),\mathrm{},d(n;n)`$ is invariant under permutation of the $`n`$ variables, we see that each moment $$E\left(d^{\kappa (1)}(i_{l(1)};n)d^{\kappa (2)}(i_{l(2)};n)\mathrm{}d^{\kappa (p)}(i_{l(p)};n)\right)$$ (8) appearing in (7) is equal to a moment $$E\left(d(1;n)^{r(1)}\overline{d(1;n)}^{s(1)}d(2;n)^{r(2)}\overline{d(2;n)}^{s(2)}\mathrm{}d(p;n)^{r(p)}\overline{d(p;n)}^{s(p)}\right)$$ (9) where $`r(1),s(1),\mathrm{},r(p),s(p)`$ $`𝐍\{0\},`$ (10) $`r(1)+s(1)+\mathrm{}+r(p)+s(p)`$ $`=p`$ and by further permutation the moment (8) corresponds to a unique moment of the form (9) if we make the additional stipulation that $$\begin{array}{c}r(1)+s(1)r(2)+s(2)\mathrm{}r(p)+s(p)\\ \text{and, if }r(j)+s(j)=r(j+1)+s(j+1)\text{ then }r(j)r(j+1).\end{array}$$ (11) Hence, rearranging the sum in (7) we get $$\begin{array}{c}\tau _n(W)=n^1\underset{r(1),s(1),\mathrm{},r(p),s(p)}{}E\left(d(1;n)^{r(1)}\overline{d(1;n)}^{s(1)}\mathrm{}d(p;n)^{r(p)}\overline{d(p;n)}^{s(p)}\right)\hfill \\ \hfill \underset{i_1,\mathrm{},i_q}{}\left(\underset{j=1}{\overset{q}{}}G^{ϵ(j)}(i_j,i_{j+1})\right)E\left(t^{ϵ(1)}(i_1,i_2;n)\mathrm{}t^{ϵ(q)}(i_q,i_{q+1};n)\right),\end{array}$$ (12) where the first sum is over all $`r(1),s(1),\mathrm{},r(p),s(p)`$ satisfying (10) and (11), and the second sum is over all $`i_1,\mathrm{},i_q\{1,\mathrm{},n\}`$ such that there is a permutation, $`\sigma `$, of $`\{1,\mathrm{},n\}`$ for which $$\begin{array}{cc}\hfill d(\sigma (i_1);n)^{\kappa (1)}d(\sigma (i_2);n)^{\kappa (2)}& \mathrm{}d(\sigma (i_p);n)^{\kappa (p)}=\hfill \\ & =d(1;n)^{r(1)}\overline{d(1;n)}^{s(1)}\mathrm{}d(p;n)^{r(p)}\overline{d(p;n)}^{s(p)}.\hfill \end{array}$$ (13) Let $`W_i`$ ($`\iota \{1,2\}`$) denote the word on the right–hand–side of (6) where $`D`$ is taken to be $`D_\iota `$. Then $$\begin{array}{c}\tau _n(W_1)\tau _n(W_2)=\hfill \\ \hfill \begin{array}{cc}\hfill =n^1\underset{r(1),s(1),\mathrm{},r(p),s(p)}{}(& E\left(a(1;n)^{r(1)}\overline{a(1;n)}^{s(1)}\mathrm{}a(p;n)^{r(p)}\overline{a(p;n)}^{s(p)}\right)\hfill \\ & E\left(b(1;n)^{r(1)}\overline{b(1;n)}^{s(1)}\mathrm{}b(p;n)^{r(p)}\overline{b(p;n)}^{s(p)}\right))\hfill \end{array}\\ \hfill \underset{i_1,\mathrm{},i_q}{}\left(\underset{j=1}{\overset{q}{}}G^{ϵ(j)}(i_j,i_{j+1})\right)E\left(t^{ϵ(1)}(i_1,i_2;n)\mathrm{}t^{ϵ(q)}(i_q,i_{q+1};n)\right).\end{array}$$ Now in order to prove the lemma it will suffice to show that $`lim_n\mathrm{}\tau _n(W_1)\tau _n(W_2)=0`$. Since by the hypothesis (5), the difference of moments of $`a`$ and $`b`$ tends to zero as $`n\mathrm{}`$, it will suffice to show that for every $`p,q𝐍`$ and every $`r(1),s(1),\mathrm{},r(p),s(p)𝐍\{0\}`$, the quantity $$n^1\underset{i_1,\mathrm{},i_q}{}\left(\underset{j=1}{\overset{q}{}}G^{ϵ(j)}(i_j,i_{j+1})\right)E\left(t^{ϵ(1)}(i_1,i_2;n)\mathrm{}t^{ϵ(q)}(i_q,i_{q+1};n)\right)$$ (14) remains bounded as $`n\mathrm{}`$, where the sum is over all $`i_1,\mathrm{},i_q\{1,\mathrm{},n\}`$ such that there is a permutation, $`\sigma `$, of $`\{1,\mathrm{},n\}`$ for which (13) holds. But this follows from the sort of counting arguments used by Voiculescu in . Indeed, for any $`1i<i^{}n`$ and for any $`m,m^{}𝐍\{0\}`$, $$E\left(t(i,i^{};n)^m\overline{t(i,i^{};n)}^m^{}\right)0$$ implies $`m=m^{}`$. Taken together with the independence assumtion on the entries of $`T`$, this shows that for any $`i_1,\mathrm{},i_q\{1,\mathrm{},n\}`$, a necessary condition so that $$E\left(t^{ϵ(1)}(i_1,i_2;n)\mathrm{}t^{ϵ(q)}(i_q,i_{q+1};n)\right)0$$ is that there be a bijection, $`\gamma `$, from $`\{1,\mathrm{},q\}`$ to itself, without fixed points, such that $`\gamma ^2=\mathrm{id}`$ and $$j\{1,\mathrm{},q\},i_{\gamma (j)}=i_{j+1},\text{ and }i_{\gamma (j)+1}=i_j.$$ (15) Moreover, there is a constant, $`c_1`$, depending only on $`q`$, such that for all $`n𝐍`$ and all choices of $`i_1,\mathrm{},i_q`$, $$\left\|E\left(t^{ϵ(1)}(i_1,i_2;n)\mathrm{}t^{ϵ(q)}(i_q,i_{q+1};n)\right)\right\|c_1n^{q/2}.$$ (16) If $`\gamma `$ is a bijection of $`\{1,\mathrm{},q\}`$, let $`N(\gamma ,n)`$ be the number of choices of $`i_1,\mathrm{},i_q\{1,\mathrm{},n\}`$ such that (15) holds. There are only finitely many bijections, $`\gamma `$, of $`\{1,\mathrm{},q\}`$. Hence, in light of the bound (16), in order to show that (14) is bounded as $`n\mathrm{}`$, it will suffice to show that for each bijection $`\gamma `$ of $`\{1,\mathrm{},q\}`$ without fixed points such that $`\gamma ^2=\mathrm{id}`$, the quantity $$n^{1(q/2)}N(\gamma ,n)$$ (17) remains bounded as $`n\mathrm{}`$. However, $`N(\gamma ,n)n^{d(\gamma )}`$, where $`d(\gamma )`$ is the number of vertices in the quotient graph, $`G_\gamma `$, which is obtained from the $`q`$–gon graph by identifying the $`j`$th and $`\gamma (j)`$th edges with opposite orientations, for every $`j`$. The graph $`G_\gamma `$ has $`q/2`$ edges, hence at most $`1+(q/2)`$ vertices, which shows that $`d(\gamma )1+(q/2)`$ and hence that (17) remains bounded as $`n\mathrm{}`$. ∎ Now we work on results that let us dispense with the permutation invariance supposed for the diagonal matrices of the previous theorem. Let $`𝒰_2`$ denote the group of unitary $`2\times 2`$ complex matrices. ###### Lemma 3.3. There is a Borel function, $`U:𝐂^3𝒰_2`$ such that for all $`a,b,c𝐂`$, $$U(a,b,c)^{}\left(\begin{array}{cc}a& b\\ 0& c\end{array}\right)U(a,b,c)=\left(\begin{array}{cc}c& b\\ 0& a\end{array}\right).$$ ###### Proof. If $`a=c`$ then let $`U(a,b,c)=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)`$. If $`ac`$ but $`b=0`$ then let $`U(a,b,c)=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)`$. If $`ac`$ and $`b0`$ then let $$U(a,b,c)=\frac{1}{(\|ac\|^2+\|b\|^2)^{1/2}}\left(\begin{array}{cc}b& \overline{(ac)}b/\overline{b}\\ ca& b\end{array}\right).$$ ###### Lemma 3.4. Fix $`n𝐍`$, let $`T\text{UTGRM}(n,\frac{1}{n})`$ and let $$D=\underset{i=1}{\overset{n}{}}d(i)e(i,i;n)_n$$ be a diagonal random matrix. Suppose that $`T`$ and $`D`$ are independent matrix–valued random variables. Let $`\pi `$ be a permutation of $`\{1,\mathrm{},n\}`$ and let $$D_\pi =\underset{i=1}{\overset{n}{}}d(\pi (i))e(i,i;n)_n.$$ (18) Let $`Z=D+T`$ and $`Z_\pi =D_\pi +T`$. Then $`Z`$ and $`Z_\pi `$ have the same $``$–moments with respect to $`\tau _n`$. ###### Proof. We may without loss of generality assume that $`\pi `$ is a transposition of neighbors: $$\pi (k)=k+1$$ $$\pi (k+1)=k$$ $$\pi (j)=j\text{if }j\{k,k+1\}.$$ We will use Lemma 3.3 to show that there is a unitary random matrix, $`V_n`$, such that $`V^{}ZV`$ has the same distribution as $`Z_\pi `$, and this will prove the theorem. Recall that $`\mathrm{\Omega }`$ is the usual probability space underlying our random matrices $`_n`$. For $`\omega \mathrm{\Omega }`$ let $`V(\omega )`$ be the block diagonal matrix $$V(\omega )=I_{k1}U(d(k)(\omega ),t(k,k+1)(\omega ),d(k+1)(\omega ))I_{nk1}.$$ By this we mean that $`V(\omega )`$ has $`k1`$ ones down the diagonal, then a $`2\times 2`$ block that is the matrix $`U`$ from Lemma 3.3, then $`nk1`$ ones. Let $`x(i,j)`$ denote the $`(i,j)`$th entry of the random matrix $`V^{}ZV`$. If we write $$T=\underset{1i<jn}{}t(i,j)e(i,j;n)$$ then $`x(i,j)`$ $`=0`$ $`\text{if }i>j`$ $`x(i,i)`$ $`=d(i)`$ $`\text{if }i\{k,k+1\}`$ $`x(k,k)`$ $`=d(k+1)`$ $`x(k+1,k+1)`$ $`=d(k)`$ $`x(k,k+1)`$ $`=t(k,k+1).`$ Let $$\left(\begin{array}{cc}u_{11}& u_{12}\\ u_{21}& u_{22}\end{array}\right)=U(d(k)(\omega ),t(k,k+1)(\omega ),d(k+1)(\omega )).$$ If $`i<k`$ then $`x(i,k)`$ $`=t(i,k)u_{11}+t(i,k+1)u_{21}`$ $`x(i,k+1)`$ $`=t(i,k)u_{12}+t(i,k+1)u_{22}`$ and if $`j>k+1`$ then $`x(k,j)`$ $`=\overline{u_{11}}t(k,j)+\overline{u_{21}}t(k+1,j)`$ $`x(k+1,j)`$ $`=\overline{u_{12}}t(k,j)+\overline{u_{22}}t(k+1,j).`$ In order to show that $`V^{}ZV`$ and $`Z_\pi `$ have the same distribution, it is thus enough to show 1. $`(x(i,j))_{1i<jn}`$ is an independent family of complex $`(0,1/n)`$–Gaussian random variables; 2. $`\{x(i,i)1in\}`$ and $`\{x(i,j)1i<jn\}`$ are independent sets of random variables. From the facts that $$\{d(i)1in\}\{t(k,k+1)\}\text{and}\{t(i,j)1i<jn,(i,j)(k,k+1)\}$$ (19) are independent sets of random variables, each $`t(i,j)`$ is complex $`(0,1/n)`$–Gaussian and $`U`$ is everywhere unitary and is independent from the right–hand set in (19), we see that $$(x(i,j))_{1i<jn,(i,j)(k,k+1)}$$ (20) is an independent family of complex $`(0,1/n)`$–Gaussian random variables. Moreover, the joint distribution of the family (20) is not changed by conditioning on the values of $`d(1),\mathrm{},d(n),t(k,k+1)`$. Hence $$\{d(i)1in\}\{t(k,k+1)\}\text{and}\{x(i,j)1i<jn,(i,j)(k,k+1)\}$$ are independent sets of random variables. But this implies that (a) and (b) hold. ∎ ###### Lemma 3.5. Let $`D_n`$ be a diagonal random matrix, let $`T\text{UTGRM}(n,\frac{1}{n})`$ and let $`Z=D+T`$. Let $`\mu `$ be the joint distribution of the $`n`$ random variables, $`d(1),d(2),\mathrm{},d(n)`$, in that order. For every permutation $`\pi `$ of $`\{1,\mathrm{},n\}`$ let $`\mu _\pi `$ be the joint distribution of the random variables, $`d(\pi (1)),d(\pi (2)),\mathrm{},d(\pi (n))`$, in that order. Let $`A`$ be a nonempty set of permutations of $`\{1,\mathrm{},n\}`$ and let $`\|A\|`$ denote the cardinality of $`A`$. Consider the measure on $`𝐂^n`$, $$\stackrel{~}{\mu }=\frac{1}{\|A\|}\underset{\pi A}{}\mu _\pi .$$ Let $`\stackrel{~}{d}(1),\stackrel{~}{d}(2),\mathrm{},\stackrel{~}{d}(n)`$ be random variables whose joint distribution is $`\stackrel{~}{\mu }`$ and such that $$\{\stackrel{~}{d}(i)1in\}\text{and}\{t(i,j)1i<jn\}$$ are independent sets of random variables. Let $$\stackrel{~}{D}=\underset{i=1}{\overset{n}{}}\stackrel{~}{d}(i)e(i,i;n)_n$$ and let $`\stackrel{~}{Z}=\stackrel{~}{D}+T(n)`$. Then $`Z`$ and $`\stackrel{~}{Z}`$ have the same $``$–moments with respect to $`\tau _n`$. ###### Proof. We may introduce a uniformly distributed $`A`$–valued random variable, $`\sigma `$, that is independent from $`D`$. Then $`\stackrel{~}{D}`$ has the same distribution as the random matrix, $`D_\sigma `$, which at a point $`\omega \mathrm{\Omega }`$ takes the value $$D_\sigma (\omega )=\underset{i=1}{\overset{n}{}}d(\sigma _{(\omega )}(i))_{(\omega )}e(i,i;n).$$ Now each $``$–moment of $`D_\sigma `$ is the average over $`\pi A`$ of the corresponding $``$–moments of the matrices $`D_\pi `$ described in (18). By Lemma 3.4, each of these is in turn equal to the corresponding $``$–moment of $`D`$. ∎ Now we combine Theorem 3.2 and Lemma 3.5 to obtain this section’s main result. ###### Theorem 3.6. For each $`n𝐍`$ let $`D_1(n),D_2(n)_n`$ be diagonal random matrices, let $`T(n)\text{UTGRM}(n,\frac{1}{n})`$ and let $`Z_\iota (n)=D_\iota (n)+T(n)`$, ($`\iota =1,2`$). Suppose that $`Z_1(n)`$ converges in $``$–moments as $`n\mathrm{}`$. For $`\iota \{1,2\}`$ and $`n𝐍`$ let $`\nu _\iota ^{(n)}`$ be the measure on $`𝐂^n`$ that is the joint distribution of the $`n`$ diagonal entries of $`D_\iota (n)`$, and let $`\stackrel{~}{\nu }_\iota ^{(n)}`$ be the symmetrization of $`\nu _\iota ^{(n)}`$. Suppose that for every $`p𝐍`$ and every $`r(1),s(1),\mathrm{},r(p),s(p)𝐍\{0\}`$, $$\begin{array}{cc}\hfill \underset{n\mathrm{}}{lim}(_{𝐂^n}\lambda _1^{r(1)}\overline{\lambda _1}^{s(1)}& \lambda _2^{r(2)}\overline{\lambda _2}^{s(2)}\mathrm{}\lambda _p^{r(p)}\overline{\lambda _p}^{s(p)}\text{d}\stackrel{~}{\nu }_1^{(n)}(\lambda _1,\mathrm{},\lambda _n)\hfill \\ & _{𝐂^n}\lambda _1^{r(1)}\overline{\lambda _1}^{s(1)}\lambda _2^{r(2)}\overline{\lambda _2}^{s(2)}\mathrm{}\lambda _p^{r(p)}\overline{\lambda _p}^{s(p)}\text{d}\stackrel{~}{\nu }_2^{(n)}(\lambda _1,\mathrm{},\lambda _n))=0.\hfill \end{array}$$ Then also $`Z_2(n)`$ converges in $``$–moments as $`n\mathrm{}`$, and its limit $``$–moments are the same as those of $`Z_1(n)`$. ## 4. Asymptotically free random matrices This section concerns asymptotic freeness of self–adjoint i.i.d. Gaussian random matrices $`Y(t,n)`$ and certain diagonal random matrices that are independent from the $`Y(t,n)`$, this being a generalization of Voiculescu’s pioneering result , which concerned constant diagonal matrices. (See and for some other generalizations.) Just as, using a technique based on the polar decomposition, Voiculescu parlayed his asymptotic freeness results for Gaussian random matrices into asymptotic freeness results for Haar distributed random unitary matrices, so in this section do we prove asymptotic freeness of Haar distributed random unitary matrices $`U(r,n)`$ and certain random diagonal matrices that are independent from the $`U(r,n)`$. Finally, this section culminates in a result (Theorem 4.6) about matrix models for ($``$–free families of) $`R`$–diagonal elements. See Notation 3.1 for details of some terms used below. ###### Notation 4.1. 1. For a random matrix $`X_n`$, we will write $`X\text{GRM}(n,\sigma ^2)`$, (the acronym is for “Gaussian random matrix”) if the entries $`x_{ij}`$ of $`X`$ ($`1i,jn`$) satisfy that $`(x_{ij})_{1i,jn}`$ is an independent family of random variables, each of which is complex $`(0,\sigma ^2)`$–Gaussian. Thus $`X\text{GRM}(n,\sigma ^2)`$ if and only if $`X`$ has density $`(\pi \sigma ^2)^{n^2}\mathrm{exp}(\frac{1}{\sigma ^2}\mathrm{Tr}(X^{}X))`$ with respect to Lebesgue measure on $`M_n(𝐂)`$. 2. Given $`Y_n`$ and $`\sigma ^2>0`$, we will write $`Y\text{SGRM}(n,\sigma ^2)`$ (“self–adjoint Gaussian random matrix”) if the entries $`y_{ij}`$ of $`Y`$ ($`1i,jn`$) satisfy that $`y_{ij}=\overline{y_{ji}}`$ for all $`i`$ and $`j`$, that $`y_{ij}`$ is complex $`(0,\sigma ^2)`$–Gaussian if $`ij`$ and is real $`(0,\sigma ^2)`$–Gaussian if $`i=j`$ and that $`(y_{ij})_{1ijn}`$ is an independent family of random variables. Thus $`Y\text{SGRM}(n,\sigma ^2)`$ if and only if $`Y`$ has density $`(\pi \sigma ^2)^{n^2/2}\mathrm{exp}(\frac{1}{\sigma ^2}\mathrm{Tr}(Y^{}Y))`$ with respect to Lebesgue measure on $`M_n^{s.a.}`$. 3. Given $`U_n`$, we will write $`U\mathrm{HURM}(n)`$, (“Haar unitary random matrix”) if $`U`$ is a random unitary matrix distributed according to Haar measure on the $`n\times n`$ unitary matrices. We begin with a preliminary result that is essentially just a combination of Theorems 2.1 and 2.2 of , in the case of random diagonal matrices. ###### Lemma 4.2. Let $`S`$ and $`T`$ be sets. For any $`n𝐍`$ let $`Y(s,n)\text{SGRM}(n,1/n)`$ ($`sS`$). and consider diagonal random matrices $`D(t,n)_n`$ ($`tT`$). Suppose that for some $`t_0T`$, $`D(t_0,n)=I_n`$ ($`n𝐍`$), that $`\{D(t,n)tT\}`$ is closed under multiplication ($`n𝐍`$) and that $`\{D(t,n)tT\}`$ converges in moments as $`n\mathrm{}`$. Suppose that for every $`n𝐍`$ $$(\{D(t,n)tT\},\left(\{Y(s,n)\}\right)_{sS})$$ (21) is an independent family of sets of matrix–valued random variables. Finally, suppose that for every $`m𝐍`$, every $`s_1,\mathrm{},s_mS`$ and every $`t_1,\mathrm{},t_mT`$, the quantity $$\left\|\tau _n\left(Y(s_1,n)D(t_1,n)\mathrm{}Y(s_m,n)D(t_m,n)\right)\right\|$$ (22) remains bounded as $`n\mathrm{}`$. Then the following are equivalent: 1. The family (21) of sets of noncommutative random variables is asymptotically free as $`n\mathrm{}`$. 2. Whenever $`m𝐍`$ is even, $`t_1,\mathrm{},t_mT`$ are fixed and $`\alpha :\{1,\mathrm{},m\}S`$ is such that for every $`sS`$, $`\alpha ^1(s)`$ has either two or zero elements, 1. if $`\alpha (1)=\alpha (2)`$ then $$\begin{array}{c}\underset{n\mathrm{}}{lim}(\tau _n\left(Y(\alpha (1),n)D(t_1,n)Y(\alpha (2),n)D(t_2,n)\mathrm{}Y(\alpha (m),n)D(t_m,n)\right)\hfill \\ \hfill \tau _n\left(D(t_1,n)\right)\tau _n\left(D(t_2,n)Y(\alpha (3),n)D(t_3,n)\mathrm{}Y(\alpha (m),n)D(t_m,n)\right))=0\end{array}$$ 2. if $`\alpha (p)\alpha (p+1)`$ for every $`1pm1`$ and if $`\alpha (m)\alpha (1)`$ then $$\underset{n\mathrm{}}{lim}\tau _n\left(Y(\alpha (1),n)D(t_1,n)Y(\alpha (2),n)D(t_2,n)\mathrm{}Y(\alpha (m),n)D(t_m,n)\right)=0.$$ ###### Proof. We may without loss of generality suppose $`S=𝐍`$ and $`T=𝐍`$. That (1)$``$(2) follows from the last paragraph of \[15, 2.1\] and the fact that the limit moments of each $`Y(s,n)`$ are those of a centered semicircle law with second moment $`1`$. To show (2)$``$(1) we will use \[15, 2.1\] and an idea from the proof of \[15, 2.2\]. Let $`\omega `$ be a nontrivial ultrafilter on $`𝐍`$. On the algebra, $`𝐂(T_s)_{s𝐍},(A_t)_{t𝐍}`$, of polynomials in noncommuting variables $`(T_s)_{s𝐍}`$ and $`(A_t)_{t𝐍}`$, let $`\varphi _\omega `$ be the tracial linear functional defined by $`\varphi _\omega (p)=lim_{n\omega }\tau _n(\pi _n(p))`$, where $`\pi _n:𝐂(T_s)_{s𝐍},(A_t)_{t𝐍}_n`$ is the algebra homomorphism defined by $`\pi _n(T_s)=Y(s,n)`$ and $`\pi _n(A_t)=D(t,n)`$. Let $`\mathrm{\Delta }`$ denote the subalgebra of $`𝐂(T_s)_{s𝐍},(A_t)_{t𝐍}`$ generated by $`1`$ and $`\{A_tt𝐍\}`$. We will check that the conditions 1 and 2 of \[15, 2.1\] hold for the sequence $`(T_s)_{s𝐍}`$ and the subalgebra $`\mathrm{\Delta }`$ with respect to $`\varphi _\omega `$. Note that every element of $`\mathrm{\Delta }`$ is a linear combination of words of the form $`A_{t_1}A_{t_2}\mathrm{}A_{t_k}`$ and that $`\pi _n(A_{t_1}A_{t_2}\mathrm{}A_{t_k})=D(t,n)`$ for some $`t𝐍`$. Moreover, $`\pi _n(1)=D(1,n)`$. Therefore $`\pi _n(\mathrm{\Delta })=\mathrm{span}\{D(t,n)t𝐍\}`$ and hence condition 1 of \[15, 2.1\] follows from the boundedness of (22) as $`n\mathrm{}`$. To see that condition 2a of \[15, 2.1\] holds, it suffices to see that if $`m𝐍`$ and if $`\alpha :\{1,\mathrm{},m\}𝐍`$ is such that $`\alpha ^1(\{\alpha (1)\})`$ has only one element and if $`t_1,\mathrm{},t_m𝐍`$ then $$n𝐍\tau _n\left(Y(\alpha (1),n)D(t_1,n)Y(\alpha (2),n)D(t_2,n)\mathrm{}Y(\alpha (m),n)D(t_m,n)\right)=0.$$ (23) But this follows from the independence of $`Y(\alpha (1),n)`$ from all the other matrices appearing in (23) and the fact that all entries of $`Y(\alpha (1),n)`$ have expectation zero. Now conditions 2b and 2c of \[15, 2.1\] follow from the hypotheses (2a) and (2b). Therefore, by \[15, 2.1\], given an injection $`\beta :𝐍\times 𝐍𝐍`$ and defining $$X_{m,k}=k^{1/2}\underset{j=1}{\overset{k}{}}T_{\beta (m,j)},$$ the family of sets of noncommutative random variables, $$(\mathrm{\Delta },\left(\{X_{m,k}\}\right)_{m=1}^{\mathrm{}})$$ is asymptotically free with respect to $`\varphi _\omega `$ as $`k\mathrm{}`$. However, using the Gaussianity of the entries of the $`Y(s,n)`$, we see that for every $`k𝐍`$, $`(\mathrm{\Delta },(\{X_{m,k}\})_{m=1}^{\mathrm{}})`$ has the same moments as $`(\mathrm{\Delta },(\{T_m\})_{m=1}^{\mathrm{}})`$. Hence $`(\mathrm{\Delta },(\{T_m\})_{m=1}^{\mathrm{}})`$ is free with respect to $`\varphi _\omega `$. Since $`\omega `$ was arbitrary, and since each $`Y(s,n)`$ converges in moments as $`n\mathrm{}`$, this implies that $$(\{D(t,n)t𝐍\},\left(\{Y(s,n)\}\right)_{s𝐍})$$ is asymptotically free as $`n\mathrm{}`$. ∎ ###### Theorem 4.3. Let $`S`$ and $`T`$ be sets. For $`sS`$ and $`n𝐍`$ let $`Y(s,n)\text{SGRM}(n,\frac{1}{n})`$. For $`tT`$ and $`n𝐍`$ let $$D(t,n)=\underset{i=1}{\overset{n}{}}d(i;n,t)e(i,i;n)_n$$ be a diagonal random matrix, and suppose that for some $`t`$ and every $`n`$, $`D(1,n)=I_n`$, that $`\{D(t,n)tT\}`$ is closed under multiplication and that $`\{D(t,n)tT\}`$ converges in moments as $`n\mathrm{}`$. Suppose that for every $`n𝐍`$ $$(\{D(t,n)tT\},\left(\{Y(s,n)\}\right)_{sS})$$ is an independent family of sets of matrix–valued random variables. Assume further that 1. for every $`tT`$ and every $`1p<\mathrm{}`$, $$\underset{\begin{array}{c}n𝐍\\ 1in\end{array}}{sup}d(i;n,t)_p<\mathrm{};$$ 2. for every $`m,n𝐍`$, $`mn`$, every $`t_1,\mathrm{},t_mT`$ and every permutation, $`\sigma `$, of $`\{1,\mathrm{},n\}`$, the joint distribution of $$(d(1;n,t_1),d(2;n,t_2),\mathrm{},d(m;n,t_m))$$ is equal to the joint distribution of $$(d(\sigma (1);n,t_1),d(\sigma (2);n,t_2),\mathrm{},d(\sigma (m);n,t_m));$$ 3. for every $`p𝐍`$, every $`t_1,\mathrm{},t_pT`$ and every $`p`$–tuple, $`(i_1,\mathrm{},i_p)`$, of distinct, positive integers, we have $$\underset{n\mathrm{}}{lim}\left(E\left(d(i_1;n,t_1)d(i_2;n,t_2)\mathrm{}d(i_p;n,t_p)\right)\underset{j=1}{\overset{p}{}}E(d(i_j;n,t_j))\right)=0.$$ Then the family $$(\{D(t,n)tT\},\left(\{Y(s,n)\}\right)_{sS})$$ of sets of random variables converges in moments and is asymptotically free as $`n\mathrm{}`$. ###### Proof. We will apply Lemma 4.2. Let us first show that the quantity (22) remains bounded as $`n\mathrm{}`$. Write $`a(i,j;n,s)`$ for the $`(i,j)`$th entry of $`Y(s,n)`$. We have $`\tau _n(Y(s_1,n)D(t_1,n)\mathrm{}Y(s_m,n)`$ $`D(t_m,n))=`$ (24) $`=n^1{\displaystyle \underset{i_1,\mathrm{},i_m\{1,\mathrm{},n\}}{}}`$ $`E\left(d(i_1;n,t_1)d(i_2;n,t_2)\mathrm{}d(i_m;n,t_m)\right)`$ $`E\left(a(i_0,i_1;n,s_1)a(i_1,i_2;n,s_2)\mathrm{}a(i_{m1},i_m;n,s_m)\right),`$ where $`i_0=i_m`$. Using the generalized Hölder inequality, we have $$\|E(d(i_1;n,t_1)\mathrm{}d(i_m;n,t_m))\|d(i_1;n,t_1)\mathrm{}d(i_m;n,t_m)_1\underset{j=1}{\overset{m}{}}d(i_j;n,t_j)_m.$$ (25) But by the assumption (i), there is $`c_2>0`$ such that $$n𝐍i_1,\mathrm{},i_m\{1,\mathrm{},n\},\underset{j=1}{\overset{m}{}}d(i_j;n,t_j)_mc_2.$$ (26) Now consider $$E\left(a(i_0,i_1;n,s_1)a(i_1,i_2;n,s_2)\mathrm{}a(i_{m1},i_m;n,s_m)\right).$$ (27) From the nature of the entries $`a(i,j;n,s)`$, we see that the quantity (27) can be nonzero only if there is a bijection $`\gamma :\{1,\mathrm{},m\}\{1,\mathrm{},m\}`$ such that $`\gamma ^2=\mathrm{id}`$, $`\gamma `$ has no fixed points and $$j\{1,\mathrm{},m\},s_j=s_{\gamma (j)},i_j=i_{\gamma (j)1},i_{j1}=i_{\gamma (j)}.$$ One also calculates $`\left\|E\left(a(i_0,i_1;n,s_1)a(i_1,i_2;n,s_2)\mathrm{}a(i_{m1},i_m;n,s_m)\right)\right\|`$ $`{\displaystyle \underset{j=1}{\overset{m}{}}}a(i_{j1},i_j;n,s_j)_m`$ $`n^{m/2}\left(\frac{m}{2}\right)!`$ Let us call a bijection, $`\gamma `$, of $`\{1,\mathrm{},m\}`$ without fixed points and such that $`\gamma ^2=\mathrm{id}`$, a pairing of $`\{1,\mathrm{},m\}`$ and for every pairing $`\gamma `$ let $$\mathrm{\Theta }(\gamma )=\{(i_1,\mathrm{},i_m)\{1,\mathrm{},n\}^mj\{1,\mathrm{},m\},i_j=i_{\gamma (j)1},i_{j1}=i_{\gamma (j)}\}.$$ From the above estimates we obtain $$\left\|\tau _n\left(Y(s_1,n)D(t_1,n)\mathrm{}Y(s_m,n)D(t_m,n)\right)\right\|c_2\left(\frac{m}{2}\right)!n^{(\frac{m}{2}+1)}\underset{\gamma }{}\|\mathrm{\Theta }(\gamma )\|,$$ (28) where the sum is over all pairings, $`\gamma `$ of $`\{1,\mathrm{},m\}`$. To each pairing we associate the quotient graph, $`G_\gamma `$, of the clockwise oriented $`m`$–gon graph obtained by identifying with opposite orientation each pair of $`j`$th and $`\gamma (j)`$th edges. The resulting graph has $`m/2`$ edges, hence at most $`\frac{m}{2}+1`$ vertices. Consequently $`\|\mathrm{\Theta }(\gamma )\|n^{\frac{m}{2}+1}`$, and the quantity in (28) is bounded by $`c_2\left(\frac{m}{2}\right)!`$ times the number of pairings, which is finite and independent of $`n`$. This shows that the quantities (22) remain bounded as $`n\mathrm{}`$. We now show that 2a and 2b of Lemma 4.2 are satisfied. Let $`\alpha `$ be as described there. Then $`\tau _n(Y(\alpha (1),n)D(t_1,n)\mathrm{}Y(\alpha `$ $`(m),n)D(t_m,n))=`$ (29) $`=n^1{\displaystyle \underset{i_1,\mathrm{},i_m\{1,\mathrm{},n\}}{}}E(`$ $`d(i_1;n,t_1)d(i_2;n,t_2)\mathrm{}d(i_m;n,t_m))`$ $`E(`$ $`a(i_0,i_1;n,\alpha (1))a(i_1,i_2;n,\alpha (2))\mathrm{}a(i_{m1},i_m;n,\alpha (m))),`$ where we let $`i_0=i_m`$. Consider the clockwise oriented $`m`$–gon graph, label the edges consecutively $`e_1,e_2,\mathrm{},e_m`$ and the vertices $`v_1,v_2,\mathrm{},v_m`$ so that the vertices of the edge $`e_j`$ are $`v_{j1}`$ and $`v_j`$, (mod $`m`$). Let $`G`$ be the quotient of the $`m`$–gon graph obtained by identifying edges $`j`$ and $`\alpha (j)`$ with opposite orientation, ($`1jm`$). The resulting identification of vertices of the $`m`$–gon graph gives an equivalence relation $``$ on $`\{v_1,\mathrm{},v_m\}`$ whose equivalence classes $`F_1,F_2,\mathrm{},F_{k(G)}`$ are precisely the lists of indices labeling the $`k(G)`$ vertices of $`G`$. The expression $$E\left(a(i_0,i_1;n,\alpha (1))a(i_1,i_2;n,\alpha (2))\mathrm{}a(i_{m1},i_m;n,\alpha (m))\right)$$ (30) in (29) is nonzero if and only if whenever $`1p,qm`$ and $`v_pv_q`$ then $`i_p=i_q`$, and then the value of (30) is $`n^{m/2}`$. For each equivalence class $`F_j=\{v_{p(1)},v_{p(2)},\mathrm{},v_{p(r_j)}\}`$ of $``$ there is $`t_j^{}T`$ so that $`D(t_j^{},n)=D(t_{p(1)},n)D(t_{p(2)},n)\mathrm{}D(t_{p(r_j)},n)`$; for $`p\{1,\mathrm{},n\}`$ let $`d(p;n,t_j^{})`$ be the $`p`$th diagonal entry of $`D(t_j^{},n)`$. Thus $`\tau _n(Y(\alpha `$ $`(1),n)D(t_1,n)\mathrm{}Y(\alpha (m),n)D(t_m,n))=`$ (31) $`=n^{(\frac{m}{2}+1)}{\displaystyle \underset{p_1,\mathrm{},p_{k(G)}\{1,\mathrm{},n\}}{}}E\left(d(p_1;n,t_1^{})d(p_2;n,t_2^{})\mathrm{}d(p_{k(G)};n,t_{k(G)}^{})\right).`$ Using the Hölder inequality estimate (25) and (26), we see that the terms $$E\left(d(p_1;n,t_1^{})d(p_2;n,t_2^{})\mathrm{}d(p_{k(G)};n,t_{k(G)}^{})\right)$$ in (31) are uniformly bounded in modulus. Moreover, because $`G`$ has $`m/2`$ edges, it follows that $`k(G)\frac{m}{2}+1`$. If $`\alpha `$ satisfies the hypothesis in 2b of Lemma 4.2, then every vertex of the $`m`$–gon graph is equivalent to at least one other vertex, so $`k(G)m/2`$ and the quantity (31) tends to zero as $`n\mathrm{}`$, as required. We have proved that 2b of Lemma 4.2 is satisfied. (In fact, a similar analysis shows that the limit moment is zero unless the pairing of $`\{1,\mathrm{},m\}`$ given by $`\alpha `$ is non–crossing — this was examined in a slightly different context in , but unfortunately without the benefit of the idea of a non–crossing pairing.) Suppose that $`\alpha `$ satisfies the hypothesis in 2a of Lemma 4.2, namely that $`\alpha (1)=\alpha (2)`$. We are interested in the limit of the moment (31) as $`n\mathrm{}`$. As the number of terms in the sum (31) where $`p_i=p_j`$ for some $`ij`$ becomes negligably small compared to $`n^{\frac{m}{2}+1}`$ as $`n\mathrm{}`$, we may in (31) sum over only all distinct choices of $`p_1,\mathrm{},p_{k(G)}\{1,\mathrm{},n\}`$. The assumption $`\alpha (1)=\alpha (2)`$ implies that $`v_1`$ is not equivalent to any other vertex under $``$. Therefore, renumbering if necessary, we may take $`F_1=\{v_1\}`$ and hence $`t_1^{}=t_1`$. By hypotheses (ii) and (iii), we have that $`\delta _n\stackrel{\text{def}}{=}E(d(p_1;n,t_1^{})`$ $`d(p_2;n,t_2^{})\mathrm{}d(p_{k(G)};n,t_{k(G)}^{}))`$ $`E\left(d(p_1;n,t_1)\right)E\left(d(p_2;n,t_2^{})\mathrm{}d(p_{k(G)};n,t_{k(G)}^{})\right)`$ is independent of the choice of distinct $`p_1,\mathrm{},p_{k(G)}\{1,\mathrm{},n\}`$ and that $`\delta _n0`$ as $`n\mathrm{}`$. Moreover, an analysis of the quotient graph $`G`$ similar to that used to obtain (31), and keeping the same notation as in (31), shows that $`\tau _n(D(t_2,n)Y(`$ $`\alpha (3),n)D(t_3,n)\mathrm{}Y(\alpha (m),n)D(t_m,n))=`$ $`=n^{\frac{m}{2}}{\displaystyle \underset{p_2,\mathrm{},p_{k(G)}\{1,\mathrm{},n\}}{}}E\left(d(p_2;n,t_2^{})\mathrm{}d(p_{k(G)};n,t_{k(G)}^{})\right).`$ But then $`n^{(\frac{m}{2}+1)}{\displaystyle \underset{p_1,\mathrm{},p_{k(G)}\{1,\mathrm{},n\}}{}}E\left(d(p_1;n,t_1)\right)E\left(d(p_2;n,t_2^{})\mathrm{}d(p_{k(G)};n,t_{k(G)}^{})\right)`$ $`=\left(n^1{\displaystyle \underset{p_1=1}{\overset{n}{}}}E\left(d(p_1;n,t_1)\right)\right)\left(n^{\frac{m}{2}}{\displaystyle \underset{p_2,\mathrm{},p_{k(G)}\{1,\mathrm{},n\}}{}}E\left(d(p_2;n,t_2^{})\mathrm{}d(p_{k(G)};n,t_{k(G)}^{})\right)\right)`$ $`=\tau _n\left(D(t_1,n)\right)\tau _n\left(D(t_2,n)Y(\alpha (3),n)D(t_3,n)\mathrm{}Y(\alpha (m),n)D(t_m,n)\right).`$ Taking the limit as $`n\mathrm{}`$, we can at will require $`p_j`$’s to be distinct and then relax this requirement; using that $`\delta _n0`$, we obtain the conclusion of 2a of Lemma 4.2. ∎ Following Voiculescu’s proof \[15, Theorem 3.8\], we will use polar decomposition to extend the asymptotic freeness result of Theorem 4.3 to include also Haar distributed random unitary matrices. ###### Theorem 4.4. Let $`R`$, $`S`$ and $`T`$ be sets. For every $`n𝐍`$ and $`sS`$ let $`Z(s,n)\text{GRM}(n,\frac{1}{n})`$ and for every $`rR`$ let $`U(r,n)\mathrm{HURM}(n)`$; furthermore, for every $`tT`$ let $$D(t,n)=\underset{i=1}{\overset{n}{}}d(i;n,t)e(i,i;n)_n$$ be diagnoal random matrices such that the family $`(D(t,n))_{tT}`$ is closed under multiplication and converges in moments as $`n\mathrm{}`$; assume further that the entries $`d(i;n,t)`$ satisfy the conditions (i), (ii) and (iii) of the statement of Theorem 4.3. Suppose that $$(\left(\{Z(s,n)\}\right)_{sS},\left(\{U(r,n)\}\right)_{rR},\{D(t,n)tT\})$$ (32) is an independent family of sets of matrix–valued random variables. Then the family $$(\left(\{Z(s,n)^{},Z(s,n)\}\right)_{sS},\left(\{U(r,n)^{},U(r,n)\}\right)_{rR},\{D(t,n)tT\})$$ (33) of sets of noncommutative random variables converges in moments and is asymptotically free as $`n\mathrm{}`$. Furthermore, each $`Z(s,n)`$ converges in $``$–moments to a circular element and each $`U(r,n)`$ converges in $``$–moments to a Haar unitary, as $`n\mathrm{}`$. ###### Proof. We shall use Theorem 4.3 and adapt the proof of \[15, Theorem 3.8\] to our situation. ###### Claim 4.4.1. The family $$(\left(\{Z(s,n)^{},Z(s,n)\}\right)_{sS},\{D(t,n)tT\})$$ of sets of random variables is asymptotically free as $`n\mathrm{}`$. ###### Proof. With $`\mathrm{Re}Z(s,n)`$ $`=(Z(s,n)+Z(s,n)^{})/2`$ $`\mathrm{Im}Z(s,n)`$ $`=(Z(s,n)Z(s,n)^{})/2i,`$ each of $`\mathrm{Re}Z(s,n)`$ and $`\mathrm{Im}Z(s,n)`$ is in $`\text{SGRM}(n,\frac{1}{2n})`$ and $$(\left(\{\mathrm{Re}Z(s,n)\}\right)_{sS},\left(\{\mathrm{Im}Z(s,n)\}\right)_{sS},\{D(t,n)tT\})$$ (34) is an independent family of sets of matrix–valued random variables. Thus, by Theorem 4.3, each $`\mathrm{Re}Z(s,n)`$ and each $`\mathrm{Im}Z(s,n)`$ converges in moments to a semicircular element and the family (34) is asymptotically free as $`n\mathrm{}`$. This proves Claim 4.4.1. ∎ Now take $`W(r,n)\text{GRM}(n,1/n)`$ so that $$(\left(\{W(r,n)\}\right)_{rR},\left(\{Z(s,n)\}\right)_{sS},\{D(t,n)tT\})$$ is an independent family of sets of matrix–valued random variables. By Claim 4.4.1, $$\left(\right(\{W(r,n)^{},W(r,n)\})_{rR},(\{Z(s,n)^{},Z(s,n)\})_{sS},\{D(t,n)tT\})$$ (35) is asymptotically free as $`n\mathrm{}`$. If $`W(r,n)=V\left(W(r,n)^{}W(r,n)\right)^{1/2}`$ is the polar decomposition of $`W(r,n)`$, then $`V`$ is almost everywhere a unitary, and is distributed according to Haar measure on the group of $`n\times n`$ unitaries. Therefore, letting $`U(r,n)`$ be the polar part, $`V`$, of $`W(r,n)`$, these random unitary matrices satisfy the hypotheses of the theorem. We will follow the proof of \[15, Theorem 3.8\] to show the asymptotic freeness of (33). For $`ϵ>0`$ let $`Y_ϵ(r,n)=W(r,n)\left(ϵ+W(r,n)^{}W(r,n)\right)^{1/2}`$. ###### Claim 4.4.2. For every $`ϵ>0`$, the family $$\left(\right(\{Y_ϵ(r,n)^{},Y_ϵ(r,n)\})_{rR},(\{Z(s,n)^{},Z(s,n)\})_{sS},\{D(t,n)tT\})$$ (36) is asymptotically free as $`n\mathrm{}`$. ###### Proof. Given $`A_n`$ and $`1p<\mathrm{}`$, let $`\|A\|_p=\left(\tau _n(A^{}A)^{p/2}\right)^{1/p}`$; moreover, let $`\|A\|_{\mathrm{}}`$ be the essential supremum of the operator norm of $`A`$ evaluated at points of the underlying probability space. Let $`q`$ be a noncommutative monomial in $`d=2a+2b+c`$ variables (for nonegative integers $`a,b,c`$), with coefficient equal to $`1`$. Let $`0<\delta 1`$. By Step I of the proof of \[15, 3.8\], letting $`f`$ be the function $`f(t)=(ϵ+t)^{1/2}`$, there is a polynomial $`Q_\delta `$ such that, letting $$A_\delta (r,n)=W(r,n)Q_\delta \left(W(r,n)^{}W(r,n)\right),$$ we have $$\underset{n\mathrm{}}{lim\; sup}\|A_\delta (r,n)Y_ϵ(r,n)\|_d<\delta .$$ Because $`\|Y_ϵ(r,n)\|_d\|Y_ϵ(r,n)\|_{\mathrm{}}1`$, it follows that $`\|A_\delta (r,n)\|_d<1+\delta `$. The assumption (i) of Theorem 4.3 on the entries of $`D(t,n)`$ implies that for all $`p1`$ and for all $`tT`$, $`sup_{n𝐍}\|D(t,n)\|_p<\mathrm{}`$. Moreover, the convergence in $``$–moments as $`n\mathrm{}`$ of $`Z(s,n)`$ implies that $`sup_{n1}\|Z(s,n)\|_p<\mathrm{}`$ whenvever $`p`$ is an even integer; however, as $`\|Z(s,n)\|_p`$ is increasing in $`p`$, this holds for all $`1p<\mathrm{}`$. Fix $`r_1,\mathrm{},r_aR`$, $`s_1,\mathrm{},s_bS`$, $`t_1,\mathrm{},t_cT`$ and let $`R_1(n,ϵ)`$ $`=q(\left(Y_ϵ(r_i,n)^{}\right)_{i=1}^a,\left(Y_ϵ(r_i,n)\right)_{i=1}^a,\left(Z(s_i,n)^{}\right)_{i=1}^b,\left(Z(s_i,n)\right)_{i=1}^b,\left(D(t_i,n)\right)_{i=1}^c)`$ $`R_2(n,ϵ,\delta )`$ $`=q(\left(A_\delta (r_i,n)^{}\right)_{i=1}^a,\left(A_\delta (r_i,n)\right)_{i=1}^a,\left(Z(s_i,n)^{}\right)_{i=1}^b,\left(Z(s_i,n)\right)_{i=1}^b,\left(D(t_i,n)\right)_{i=1}^c).`$ We may chose a constant $`K`$ indepent of $`\delta `$ and large enough so that $`i\{1,\mathrm{},b\}`$ $`\underset{n\mathrm{}}{lim\; sup}\|Z(s_i,n)\|_d<K`$ (37) $`i\{1,\mathrm{},c\}`$ $`\underset{n\mathrm{}}{lim\; sup}\|D(t_i,n)\|_d<K`$ (38) Using Hölder’s inequality we find $$\underset{n\mathrm{}}{lim\; sup}\|R_1(n,ϵ)R_2(n,ϵ,\delta )\|_12aK^{2b+c}(1+\delta )^{2a1}\delta ,$$ and therefore $$\underset{\delta 0}{lim}\underset{n\mathrm{}}{lim\; sup}\left\|\tau _n\left(R_1(n,ϵ)\right)\tau _n\left(R_2(n,ϵ,\delta )\right)\right\|=0.$$ (39) The asymptotic freeness of $$\left(\right(\{A_\delta (r,n)^{},A_\delta (r,n)\})_{rR},(\{Z(s,n)^{},Z(s,n)\})_{sS},\{D(t,n)tT\})$$ follows from that of (35); this together with (39) implies the asymptotic freeness of (36), and claim 4.4.2 is proved. ∎ Step III of the proof of \[15, 3.8\] shows that for every $`\theta >0`$ there is $`ϵ_0>0`$ such that $$\underset{n\mathrm{}}{lim\; sup}\|Y_ϵ(r,n)U(r,n)\|_d<\theta $$ (40) whenever $`0<ϵϵ_0`$. Let again $`q`$ be a noncommutative monomial having coefficient equal to $`1`$ and with degree $`d=2a+2b+c`$, and let $`r_1,\mathrm{},r_aR`$, $`s_1,\mathrm{},s_bS`$, $`t_1,\mathrm{},t_cT`$. Let $$R_3(n)=q(\left(U(r_i,n)^{}\right)_{i=1}^a,\left(U(r_i,n)\right)_{i=1}^a,\left(Z(s_i,n)^{}\right)_{i=1}^b,\left(Z(s_i,n)\right)_{i=1}^b,\left(D(t_i,n)\right)_{i=1}^c)$$ Letting $`K`$ be a constant so that (37) and (38) hold, we easily see using (40) and Hölder’s inequality that if $`0<ϵϵ_0`$ then $$\underset{n\mathrm{}}{lim\; sup}\|R_1(n,ϵ)R_3(n)\|_12aK^{2b+c}\theta .$$ Therefore $$\underset{ϵ0}{lim}\underset{n\mathrm{}}{lim\; sup}\left\|\tau _n\left(R_1(n,ϵ)\right)\tau _n\left(R_3(n)\right)\right\|=0.$$ This, together with Claim 4.4.2 shows that the family (33) is asympototically free as $`n\mathrm{}`$ and finishes the proof of the theorem. ∎ The following sort of result is standard, but a proof is provided here for completeness. ###### Lemma 4.5. Let $`(A,\varphi )`$ be a W–noncommutative probability space, let $`B`$ be a unital subalgebra of $`A`$, let $`S`$ be a set and for every $`sS`$ let $`v_sA`$ be a unitary with $`\varphi (v_s)=0`$. Suppose the family $$(B,\left(\{v_s^{},v_s\}\right)_{sS})$$ (41) of $`\|S\|+1`$ sets of noncommutative random vairables is free. Then the family $$(v_sBv_s^{})_{sS}$$ (42) of unital subalgebras of $`A`$ is free. ###### Proof. From $``$–freeness of $`B`$ and $`v_s`$ we get $`\varphi (v_sbv_s^{})=\varphi (b)`$ for every $`bB`$. Let $`n𝐍`$ and let $`s_1,\mathrm{},s_nS`$ be so that $`s_js_{j+1}`$ for every $`j\{1,\mathrm{},n1\}`$. For every $`j\{1,\mathrm{},n\}`$ let $`b_jB`$ be such that $`\varphi (b_j)=0`$. In order for freeness of (42) to hold, it will suffice that $$\varphi \left((v_{s_1}b_1v_{s_1}^{})(v_{s_2}b_2v_{s_2}^{})\mathrm{}(v_{s_n}b_nv_{s_n}^{})\right)=0.$$ But the above equality follows directly from freeness of (41). ∎ Now we apply the asymptotic freeness results proved previously in this section to give some matrix models for ($``$–free families of) $`R`$–diagonal elements. ###### Theorem 4.6. Let $`S`$ be a set and for every $`sS`$ let $`X_s(n)_n`$ and let $`\sigma _{s,n}`$ be the symmetrized joint distribution of the eigenvalues of $`\left(X_s^{}(n)X_s(n)\right)^{1/2}`$. Given $`p\{1,\mathrm{},n\}`$, let $`\sigma _{s,n}^{(p)}`$ be the marginal distribution of $`\sigma _{s,n}`$ corresponding to $`p`$ of the variables. Suppose that for a compactly supported measure $`\rho _s`$ on $`𝐑_+`$ and for every $`p𝐍`$, $`\sigma _{s,n}^{(p)}`$ converges in moments as $`n\mathrm{}`$ to the product measure $`\underset{1}{\overset{p}{\times }}\rho _s`$. Suppose also that for any non–random $`n\times n`$ unitary matrix $`U`$, the distributions of $`UX_s(n)`$ and of $`X_s(n)`$ are the same. 1. Fix $`sS`$. Then $`X_s(n)`$ converges in $``$–moments to an element $`u_sh_s`$ of a noncommutative probability space, where $`u_s`$ is a Haar unitary, $`h_s0`$, $`h_s`$ has the same moments as the measure $`\rho _s`$ and where the pair $`(\{u_s,u_s^{}\},\{h_s\})`$ of sets of noncommutative random variables is free. 2. Suppose in addition that $$\left(X_s(n)\right)_{sS}$$ (43) is a mutually independent family of matrix–valued random variables and that the joint $``$–moments of (43) are the same as the joint $``$–moments of $$\left(U_s^{(1)}X_s(n)U_s^{(2)}\right)_{sS}$$ (44) whenever $`U_s^{(1)}`$ and $`U_s^{(2)}`$ are non–random unitary matrices, ($`sS`$). Then the family (43) is asymptotically $``$–free as $`n\mathrm{}`$. ###### Proof. For brevity we shall prove parts (i) and (ii) simultaneously; while proving (i), we may from the outset assume that the stronger hypotheses of (ii) hold, because if we require $`S`$ to be a single element then they will in any case be satisfied. We may write $`X_s(n)=V_s(n)H_s(n)`$ where $`V_s(n)`$ is a random unitary matrix and $`H_s(n)=\left(X_s(n)^{}X_s(n)\right)^{1/2}`$. For every $`sS`$ let $`W_s(n)`$ be a random unitary matrix so that $`D_s(n)=W_s(n)^{}H_s(n)W_s(n)`$ is diagonal, so that the joint distribution of the diagonal entries of $`D_s(n)`$ is invariant under all permutations of the $`n`$ variables and so that $$\left(\{D_s(n),V_s(n),W_s(n)\}\right)_{sS}$$ is a mutually independent family of sets of matrix–valued random variables. Let $`U_s^{(1)}(n),U_s^{(2)}(n)\mathrm{HURM}(n)`$ be so that $$(\left(\{H_s(n),V_s(n),W_s(n)\}\right)_{sS},\left(\{U_s^{(1)}(n)\}\right)_{sS},\left(\{U_s^{(2)}(n)\}\right)_{sS})$$ is a mutually independent family of sets of matrix–valued random variables. It follows from the hypotheses of (ii) that $$\left(U_s^{(1)}(n)X_s(n)U_s^{(2)}(n)\right)_{sS}$$ has the same joint $``$–moments as the family (43). We have $$U_s^{(1)}(n)X_s(n)U_s^{(2)}(n)=\left(U_s^{(1)}(n)V_s(n)U_s^{(2)}(n)\right)\left(U_s^{(2)}(n)^{}W_s(n)D_s(n)W_s(n)^{}U_s^{(2)}(n)\right).$$ Let $`\stackrel{~}{V}_s(n)`$ $`=U_s^{(1)}(n)V_s(n)U_s^{(2)}(n)`$ $`\stackrel{~}{W}_s(n)`$ $`=U_s^{(2)}(n)^{}W_s(n).`$ Then $`\stackrel{~}{V}_s(n),\stackrel{~}{W}_s(n)\mathrm{HURM}(n)`$ and $$(\left(D_s(n)\right)_{sS},\left(\stackrel{~}{V}_s(n)\right)_{sS},\left(\stackrel{~}{W}_s(n)\right)_{sS})$$ (45) is an independent family of matrix–valued random variables. Let $$\mathrm{\Delta }(n)=\{I_n\}\{D_{s_1}(n)D_{s_2}(n)\mathrm{}D_{s_q}(n)q𝐍,s_1,\mathrm{},s_qS\}.$$ By hypothesis, each $`D_s(n)`$ converges in moments as $`n\mathrm{}`$. Since $`\mathrm{\Delta }(n)`$ forms a commuting family of self–adjoint random matrices, and since the family $$\left(D_s(n)\right)_{sS,}$$ (46) is independent, it follows that $`\mathrm{\Delta }(n)`$ converges in moments as $`n\mathrm{}`$; moreover, the subfamily (46) converges in $``$–moments to a family $`(d_s)_{sS}`$ is some W–noncommutative probability space $`(A,\varphi )`$, where $`d_s`$ is positive and has the same moments as the measure $`\sigma _s`$, and where for distinct $`s_1,\mathrm{},s_mS`$ and any $`k_1,\mathrm{},k_m𝐍`$, $$\varphi (d_{s_1}^{k_1}d_{s_2}^{k_2}\mathrm{}d_{s_m}^{k_m})=\underset{j=1}{\overset{m}{}}\varphi (d_{s_j}^{k_j}).$$ (47) We shall show that the entries of the set $`\mathrm{\Delta }(n)`$ of diagonal random matrices satisfy the properties (i), (ii) and (iii) in the statement of Theorem 4.3. Let $`d_s(i,n)`$ denote the $`i`$th diagonal entry of $`D_s(n)`$. Note that $`E(d_s(i,n)^k)`$ stays bounded (in fact converges) as $`n\mathrm{}`$, for every $`k𝐍`$ and $`sS`$. This, together with the independence of the family (46), implies condition (i). Condition (ii) follows from the independence of (46) and the fact that the joint distribution of the diagonal entries of each $`D_s(n)`$ is invariant under permutations of the $`n`$ variables. Because $`\sigma _{s,n}^{(p)}`$ converges to a product measure, we have for every $`sS`$, $`m𝐍`$, $`k_1,\mathrm{},k_m𝐍`$ and every $`p`$–tuple $`(i_1,i_2,\mathrm{},i_m)`$ of distinct, positive integers, that $$\underset{n\mathrm{}}{lim}\left(E\left(d_s(i_1,n)^{k_1}d_s(i_2,n)^{k_2}\mathrm{}d_s(i_m,n)^{k_m}\right)\underset{j=1}{\overset{m}{}}E(d_s(i_j,n)^{k_j})\right)=0.$$ This together with the independence of (46) implies condition (iii). Hence we conclude from Theorem 4.4 that $$(\{D_s(n)sS\},\left(\{\stackrel{~}{V}_s(n)^{},\stackrel{~}{V}_s(n)\}\right)_{sS},\left(\{\stackrel{~}{W}_s(n)^{},\stackrel{~}{W}_s(n)\}\right)_{sS})$$ is asymptotically free as $`n\mathrm{}`$. Therefore the family (45) converges in $``$–moments to a family $$((d_s)_{sS},(v_s)_{sS},(w_s)_{sS})$$ in some W–noncommutative probability space, where the joint $``$–moments of $`(d_s)_{sS}`$ are as described above, where each $`v_s`$ and each $`w_s`$ is a Haar unitary and where $$(\{d_ssS\},\left(\{v_s^{},v_s\}\right)_{sS},\left(\{w_s^{},w_s\}\right)_{sS})$$ (48) is free. Therefore, the family (43) converges in $``$–moments as $`n\mathrm{}`$ to the family $$\left(v_s(w_sd_sw_s^{})\right)_{sS}.$$ It is clear that $`w_sd_sw_s^{}`$ has the same moments as $`d_s`$, namely the same moments as the measure $`\sigma _s`$. From the freeness of (48) and Lemma 4.5, it follows that the family $$((v_s)_{sS},(w_sd_sw_s^{})_{sS})$$ is $``$–free, and the theorem is proved. ∎ ## 5. Upper triangular representations of circular free Poisson elements In this section, the random matrix results of §3 and §4 are used, together with results of Dyson and others, to give upper triangular matrix models of circular free Poisson elements, and finally to give an upper triangular realization of a circular free Poisson element. An outline of the contents of this section is as follows: a first intermediate goal is a unitarily invariant matrix model for a circular free Poisson element (Theorem 5.4); next, a result of Dyson is quoted (Theorem 5.5) and used to convert the unitarily invariant matrix model to an upper triangular matrix model for a circular free Poisson element (Corollary 5.6); then the diagonal elements of this upper triangular matrix model are decoupled and desymmetrized so as to yield, in the limit as matrix size increases without bound, a triangular realization of a circular free Poisson element (Theorem 5.10). The following is due to Bronk ; see also \[7, §5\]. ###### Theorem 5.1. Let $`c1`$ and let $`Y`$ be an $`n\times n`$ random matrix whose density with respect to Lebesgue measure on $`M_n(𝐂)`$ is $$K_{c,n}^{(1)}\|detY\|^{2(c1)n}\mathrm{exp}\left(n\mathrm{Tr}(Y^{}Y)\right),$$ where $`K_{c,n}^{(1)}`$ is a constant. Then the symmetrized joint distribution of the eigenvalues of $`Y^{}Y`$ has density $$K_{c,n}^{(2)}\left(\underset{i=1}{\overset{n}{}}\lambda _i\right)^{2(c1)n}\left(\underset{1i<jn}{}(\lambda _i\lambda _j)^2\right)\mathrm{exp}\left(n\underset{i=1}{\overset{n}{}}\lambda _i\right)$$ (49) with respect to Lebesgue measure on $`(𝐑_+)^n`$, where $`K_{c,n}^{(2)}`$ is a constant. The next theorem is a corollary of a result of Hewitt and Savage . ###### Theorem 5.2. Let $`(\mathrm{\Omega },𝐄)`$ be a standard Borel space. Let $`\sigma `$ be a Borel probability measure on the product set $`\mathrm{\Omega }^𝐍=_{n=1}^{\mathrm{}}\mathrm{\Omega }`$ endowed with the product topology. Let $`\sigma _1`$ and $`\sigma _2`$ be the probability measures on $`\mathrm{\Omega }`$ and $`\mathrm{\Omega }\times \mathrm{\Omega }`$, respectively, determined by $`\sigma _1(A)`$ $`=\sigma (A\times \mathrm{\Omega }\times \mathrm{\Omega }\times \mathrm{}),A𝐄`$ $`\sigma _2(A_1\times A_2)`$ $`=\sigma (A_1\times A_2\times \mathrm{\Omega }\times \mathrm{\Omega }\times \mathrm{}),A_1,A_2𝐄.`$ Suppose that 1. $`\sigma `$ is invariant under all finite permutations of coordinates in $`\mathrm{\Omega }^𝐍`$, (i.e. those permutations leaving all but finitely many coordinates fixed); 2. $`\sigma _2=\sigma _1\times \sigma _1`$. Then $`\sigma `$ is equal to the product measure $`\underset{n=1}{\overset{\mathrm{}}{\times }}\sigma _1`$. ###### Proof. Since any noncountable standard Borel space is Borel isomorphic to the unit interval, and since $`(𝐍,2^𝐍)`$ is Borel isomorphic to the one–point compactification of $`𝐍`$, it is no loss of generality to assume that $`\mathrm{\Omega }`$ is a separable compact Hausdorff space and $`𝐄`$ is the Borel $`\sigma `$–algebra associated to this topology. For any compact set $`K`$, let $`P(K)`$ denote the set of Borel probability measures on $`K`$. Consider the folowing subsets of $`P(\mathrm{\Omega }^𝐍)`$: $`\stackrel{~}{P}`$ $`=\{\underset{n=1}{\overset{\mathrm{}}{\times }}\mu \mu P(\mathrm{\Omega })\},`$ $`\stackrel{~}{S}`$ $`=\{\nu P(\mathrm{\Omega }^𝐍)\nu \text{ is invariant under all finite permutations of the coordinates of }\mathrm{\Omega }^𝐍\}.`$ Clearly $`\stackrel{~}{P}\stackrel{~}{S}`$. By \[8, Theorem 7.2\], every $`\nu \stackrel{~}{S}`$ has a representation $$\nu =_{P(\mathrm{\Omega })}\left(\underset{n=1}{\overset{\mathrm{}}{\times }}\mu \right)\text{d}\rho (\mu ),$$ for a unique $`\rho P(P(\mathrm{\Omega }))`$. In fact, (see \[13, Theorem 3.1\]), $`\stackrel{~}{P}`$ is the set of extreme points of the compact simplex $`\stackrel{~}{S}`$. Now let $`\sigma `$ be as in the formulation of the theorem. Using hypothesis (i) we have $$\sigma =_{P(\mathrm{\Omega })}\left(\underset{n=1}{\overset{\mathrm{}}{\times }}\mu \right)\text{d}\rho (\mu )$$ for a unique $`\rho P(P(\mathrm{\Omega }))`$. In particular $$\sigma _1=_{P(\mathrm{\Omega })}\mu \text{d}\rho (\mu )\text{and}\sigma _2=_{P(\mathrm{\Omega })}(\mu \times \mu )\text{d}\rho (\mu ).$$ By the assumption on $`\mathrm{\Omega }`$, the space $`C(\mathrm{\Omega })`$ of complex valued continuous functions on $`\mathrm{\Omega }`$ is a separable Banach space (in the uniform norm), so we may let $`F`$ be a countable dense subset of $`C(\mathrm{\Omega })`$. Given $`fC(\mathrm{\Omega })`$ and $`\lambda P(\mathrm{\Omega })`$ let us write $$\lambda (f)=_\mathrm{\Omega }f\text{d}\lambda .$$ With this notation, we have for all $`fF`$, $`{\displaystyle _{P(\mathrm{\Omega })}}\|\mu (f)\sigma _1(f)\|^2\text{d}\rho (\mu )`$ $`={\displaystyle _{P(\mathrm{\Omega })}}(\mu \times \mu )(f\overline{f})\text{d}\rho (\mu )2\mathrm{R}\mathrm{e}\left(\overline{\sigma _1(f)}{\displaystyle _{P(\mathrm{\Omega })}}\mu (f)\text{d}\rho (\mu )\right)+\|\sigma _1(f)\|^2`$ $`=\sigma _2(f\overline{f})2\mathrm{R}\mathrm{e}\left(\overline{\sigma _1(f)}\sigma _1(f)\right)+\|\sigma _1(f)\|^2`$ $`=\sigma _2(f\overline{f})\|\sigma _1(f)\|^2.`$ But hypothesis (ii) shows that the above quantity is zero. Hence $`\mu (f)=\sigma _1(f)`$ for all $`fF`$, for $`\rho `$–almost all $`\mu P(\mathrm{\Omega })`$. Hence $`\mu =\sigma _1`$ for $`\rho `$–almost all $`\mu P(\mathrm{\Omega })`$, which implies $`\rho =\delta _{\sigma _1}`$, the Dirac measure at the point $`\sigma _1`$. Therefore $`\sigma =\underset{n=1}{\overset{\mathrm{}}{\times }}\sigma _1`$. ∎ ###### Lemma 5.3. Let $`c1`$. Given $`n𝐍`$ let $`\mu _n`$ be the measure having density (49) with respect to Lebesgue measure on $`𝐑_+^n`$. Given $`p\{1,2,\mathrm{},n\}`$ let $`\mu _n^{(p)}`$ be the marginal distribution of $`\mu _n`$ corresponding to the variables $`\lambda _1,\mathrm{},\lambda _p`$. Fix $`p𝐍`$. Then the distribution $`\mu _n^{(p)}`$ converges in the weak topology as $`n\mathrm{}`$ to the product measure $`\underset{1}{\overset{p}{\times }}\tau `$, where $`\tau `$ has density with respect to Lebesgue measure $$\frac{\text{d}\tau }{\text{d}\lambda }=\frac{\sqrt{(\lambda a)(b\lambda )}}{2\pi \lambda }1_{[a,b]}(\lambda ),$$ (50) with $`a=(1\sqrt{c})^2`$ and $`b=(1+\sqrt{c})^2`$. Moreover, if $`f`$ is a continuous function on $`[0,\mathrm{})^p`$ with polynomial growth, in the sense that $`f(t_1,\mathrm{},t_p)K^{(3)}(1+t_1^{k_1}t_2^{k_2}\mathrm{}t_p^{k_p})`$ for some constant $`K^{(3)}>0`$ and positive integers $`k_1,\mathrm{},k_p`$, then $$\underset{n\mathrm{}}{lim}_{𝐑_+^p}f\text{d}\mu _n^{(p)}=_{𝐑_+^p}f\text{d}\left(\underset{1}{\overset{p}{\times }}\tau \right).$$ (51) ###### Proof. It will be more convenient to consider the measure $`\sigma _n`$ whose density with respect to Lebesgue measure on $`𝐑_+^n`$ is $$K_{c,n}^{(4)}\left(\underset{i=1}{\overset{n}{}}\lambda _i\right)^{2(c1)n}\left(\underset{1i<jn}{}(\lambda _i\lambda _j)^2\right)\mathrm{exp}\left(\underset{i=1}{\overset{n}{}}\lambda _i\right),$$ for some constant $`K_{c,n}^{(4)}`$; thus $`\sigma _n`$ is the push forward measure of $`\mu _n`$ under the transformation $`(\lambda _1,\lambda _2,\mathrm{})(n\lambda _1,n\lambda _2,\mathrm{})`$. We will find the limit as $`n\mathrm{}`$ of the marginal distributions, $`\sigma _n^{(p)}`$, of $`\sigma _n`$ corresponding to the variables $`\lambda _1,\mathrm{},\lambda _p`$. Let $`\alpha =2(c1)n`$ and let $`\varphi _0^{(\alpha )},\varphi _1^{(\alpha )},\varphi _2^{(\alpha )},\mathrm{}`$ be the polynomials obtained via Gram–Schmidt orthonormalizaton of $`1,\lambda ,\lambda ^2,\mathrm{}`$ in $`L^2([0,\mathrm{}),\lambda ^\alpha e^\lambda \text{d}\lambda )`$. Thus $`\varphi _k^{(\alpha )}(\lambda )=\sqrt{\frac{n!}{\mathrm{\Gamma }(\alpha +n+1)}}L_k^\alpha (\lambda )`$, where $`L_k^\alpha `$ are the (generalized) Laguerre polynomials. Using the Vandermonde determinant we have $`{\displaystyle \underset{1i<jn}{}}(\lambda _j\lambda _i)`$ $`=det\left(\begin{array}{cccc}1& 1& \mathrm{}& 1\\ \lambda _1& \lambda _2& \mathrm{}& \lambda _n\\ \mathrm{}& \mathrm{}& & \mathrm{}\\ \lambda _1^{n1}& \lambda _2^{n1}& \mathrm{}& \lambda _n^{n1}\end{array}\right)`$ $`=K_{\alpha ,n}^{(5)}det\left(\begin{array}{cccc}\varphi _0^{(\alpha )}(\lambda _1)& \varphi _0^{(\alpha )}(\lambda _2)& \mathrm{}& \varphi _0^{(\alpha )}(\lambda _n)\\ \varphi _1^{(\alpha )}(\lambda _1)& \varphi _1^{(\alpha )}(\lambda _2)& \mathrm{}& \varphi _1^{(\alpha )}(\lambda _n)\\ \mathrm{}& \mathrm{}& & \mathrm{}\\ \varphi _{n1}^{(\alpha )}(\lambda _1)& \varphi _{n1}^{(\alpha )}(\lambda _2)& \mathrm{}& \varphi _{n1}^{(\alpha )}(\lambda _n)\end{array}\right)`$ for some constant $`K_{\alpha ,n}^{(5)}`$. Therefore, the density of $`\sigma _n`$ with respect to Lebesgue measure on $`𝐑_+^n`$ is $$D_n(\lambda _1,\mathrm{},\lambda _n)=K_{c,n}^{(6)}\left(\underset{i=1}{\overset{n}{}}\lambda _i\right)^\alpha \left(\underset{\pi S_n}{}\text{sign}(\pi )\underset{i=1}{\overset{n}{}}\varphi _{\pi (i)1}^{(\alpha )}(\lambda _i)\right)^2\mathrm{exp}\left(\underset{i=1}{\overset{n}{}}\lambda _i\right),$$ for a constant $`K_{c,n}^{(6)}`$. Writing $$\underset{i=1}{\overset{n}{}}\varphi _{\pi (i)1}^{(\alpha )}(\lambda _i)=\varphi _{\pi (1)1}^{(\alpha )}(\lambda _1)\varphi _{\pi (2)1}^{(\alpha )}(\lambda _2)\mathrm{}\varphi _{\pi (n)1}^{(\alpha )}(\lambda _n)$$ and noting that as $`\pi `$ ranges over the permutation group $`S_n`$ these form an orthonormal family with respect to the measure $`\left(_{i=1}^n\lambda _i\right)^\alpha \mathrm{exp}\left(_{i=1}^n\lambda _i\right)\text{d}\lambda _1\mathrm{}\text{d}\lambda _n`$ on $`𝐑_+^n`$, we find $`K_{c,n}^{(6)}=(n!)^1`$. Moreover, the density with respect to Lebesgue measure on $`𝐑_+`$ of the marginal distribution $`\sigma _n^{(1)}`$ is $$D_{n,1}(\lambda _1)\stackrel{\text{def}}{=}_{𝐑_+^{n1}}D_n(\lambda _1,\mathrm{},\lambda _n)\text{d}\lambda _2\mathrm{}\text{d}\lambda _n=\frac{1}{n}\left(\underset{k=0}{\overset{n1}{}}\varphi _k^{(\alpha )}(\lambda _1)^2\right)\lambda _1^\alpha e^{\lambda _1}.$$ But then the treatment in §6 of shows that $`\mu _n^{(1)}`$ converges in the weak topology and in moments as $`n\mathrm{}`$ to $`\tau `$. The density with respect to Lebesgue measure on $`𝐑_+`$ of the marginal distribution $`\sigma _n^{(2)}`$ is $`D_{n,2}(\lambda _1)`$ $`\stackrel{\text{def}}{=}{\displaystyle _{𝐑_+^{n2}}}D_n(\lambda _1,\mathrm{},\lambda _n)\text{d}\lambda _3\mathrm{}\text{d}\lambda _n`$ $`=\begin{array}{cc}\hfill {\displaystyle \frac{1}{n(n1)}}& (\lambda _1\lambda _2)^\alpha e^{(\lambda _1+\lambda _2)}\hfill \\ \hfill & {\displaystyle \underset{\begin{array}{c}0k,\mathrm{}n1\\ k\mathrm{}\end{array}}{}}\varphi _k^{(\alpha )}(\lambda _1)\varphi _{\mathrm{}}^{(\alpha )}(\lambda _2)\left(\varphi _k^{(\alpha )}(\lambda _1)\varphi _{\mathrm{}}^{(\alpha )}(\lambda _2)\varphi _{\mathrm{}}^{(\alpha )}(\lambda _1)\varphi _k^{(\alpha )}(\lambda _2)\right)\hfill \end{array}`$ $`={\displaystyle \frac{n}{n1}}D_{n,1}(\lambda _1)D_{n,1}(\lambda _2){\displaystyle \frac{1}{n(n1)}}(\lambda _1\lambda _2)^\alpha e^{(\lambda _1+\lambda _2)}\left({\displaystyle \underset{j=0}{\overset{n1}{}}}\varphi _j^{(\alpha )}(\lambda _1)\varphi _j^{(\alpha )}(\lambda _2)\right)_.^2`$ As elements of $`C_0(𝐑_+^2)^{}`$, we thus have $`\sigma _n^{(2)}`$ $`\sigma _n^{(1)}\sigma _n^{(1)}`$ (52) $`{\displaystyle \frac{1}{n1}}{\displaystyle _0^{\mathrm{}}}{\displaystyle _0^{\mathrm{}}}D_{n,1}(\lambda _1)D_{n,1}(\lambda _2)\text{d}\lambda _1\text{d}\lambda _2+`$ $`+{\displaystyle \frac{1}{n(n1)}}{\displaystyle \underset{j=0}{\overset{n1}{}}}{\displaystyle _0^{\mathrm{}}}{\displaystyle _0^{\mathrm{}}}\varphi _j^{(\alpha )}(\lambda _1)^2\varphi _j^{(\alpha )}(\lambda _2)^2(\lambda _1\lambda _2)^\alpha e^{(\lambda _1+\lambda _2)}\text{d}\lambda _1\text{d}\lambda _2`$ $`={\displaystyle \frac{2}{n1}}.`$ Since we know that $`\mu _n^{(1)}`$ converges in the weak topology as $`n\mathrm{}`$ to $`\tau `$, it follows from (52) that $`\mu _n^{(2)}`$ converges in weak topology as $`n\mathrm{}`$ to $`\tau \times \tau `$. Consider the measures $`\stackrel{~}{\mu }_n=\mu _n\times \delta _0\times \delta _0\times \mathrm{}`$ on $`[0,\mathrm{})^𝐍`$ and let $`\nu `$ be a w cluster point in $`C_0([0,\mathrm{})^𝐍)^{}`$ of these. Let $`\nu ^{(p)}`$ be the marginal distribution of $`\nu `$ corresponding to the first $`p`$ coordinates of $`[0,\mathrm{})^𝐍`$. Then from what we have proved above we have 1. $`\nu `$ is invariant under finite permutations of the coordinates in $`[0,\mathrm{})^𝐍`$; 2. $`\nu ^{(1)}=\tau `$; 3. $`\nu ^{(2)}=\tau \times \tau `$. Hence, by Theorem 5.2, $`\nu =\underset{p=1}{\overset{\mathrm{}}{\times }}\nu ^{(1)}`$. Since $`\nu `$ was an arbitrary cluster point of $`(\stackrel{~}{\mu }_n)_{n=1}^{\mathrm{}}`$ it follows that $`\stackrel{~}{\mu }_n`$ converges in weak topology to $`\underset{p=1}{\overset{\mathrm{}}{\times }}\tau `$ as $`n\mathrm{}`$. Therefore, for all $`p𝐍`$, the marginal distribution $`\nu _n^{(p)}`$ converges in weak topology to the measure $`\underset{1}{\overset{p}{\times }}\tau `$ as $`n\mathrm{}`$. It remains to show that (51) holds whenever $`f`$ is of polynomial growth. ###### Claim 5.3.1. Let $`p𝐍`$ and let $`h`$ be a positive continuous function on $`[0,\mathrm{})^p`$. Then $`lim\; inf_n\mathrm{}h\text{d}\mu _n^{(p)}h\text{d}\nu ^{(p)}`$. ###### Proof. Choose $`h_jC_0\left([0,\mathrm{})^p\right)`$, $`h_j0`$, so that $`h_j`$ increases pointwise to $`h`$ as $`j\mathrm{}`$. Then for all $`j1`$, $$\underset{n\mathrm{}}{lim\; inf}h\text{d}\mu _n^{(p)}\underset{n\mathrm{}}{lim\; inf}h_j\text{d}\mu _n^{(p)}=h_j\text{d}\nu ^{(p)}.$$ But $`\nu ^{(p)}=\underset{1}{\overset{p}{\times }}\nu ^{(1)}`$ is supported on $`[a,b]^p`$; therefore $`sup_jh_j\text{d}\nu ^{(p)}=h\text{d}\nu ^{(p)}`$, and the claim is proved. ∎ ###### Claim 5.3.2. Let $`p𝐍`$ and suppose $`f`$ and $`g`$ are continuous functions on $`[0,\mathrm{})^p`$ satisfying $`g0`$ and $`gfg`$, and suppose that $`lim_n\mathrm{}g\text{d}\mu _n^{(p)}=g\text{d}\nu ^{(p)}`$. Then $`lim_n\mathrm{}f\text{d}\mu _n^{(p)}=f\text{d}\nu ^{(p)}`$. ###### Proof. Applying Claim 5.3.1 to $`gf`$ gives $`lim\; sup_n\mathrm{}f\text{d}\mu _n^{(p)}f\text{d}\nu ^{(p)}`$, while applying Claim 5.3.1 to $`g+f`$ yields $`lim\; inf_n\mathrm{}f\text{d}\mu _n^{(p)}f\text{d}\nu ^{(p)}`$; the claim is proved. ∎ In order to finish the proof of the lemma, it will suffice to show $$\underset{n\mathrm{}}{lim}\lambda _1^{k_1}\lambda _2^{k_2}\mathrm{}\lambda _p^{k_p}\text{d}\mu _n^{(p)}=\lambda _1^{k_1}\lambda _2^{k_2}\mathrm{}\lambda _p^{k_p}\text{d}\nu ^{(p)},$$ (53) for every $`p𝐍`$ and all integers $`k_1,k_2,\mathrm{},k_p0`$. Letting $`\mathrm{}=k_1+k_1+\mathrm{}+k_p`$, we have $`0\lambda _1^{k_1}\lambda _2^{k_2}\mathrm{}\lambda _p^{k_p}1+\lambda _1^{\mathrm{}}+\lambda _2^{\mathrm{}}+\mathrm{}+\lambda _p^{\mathrm{}}`$. Moreover, because $`\mu _n^{(1)}`$ converges in moments to $`\nu ^{(1)}`$, we have $`\underset{n\mathrm{}}{lim}{\displaystyle (1+\lambda _1^{\mathrm{}}+\lambda _2^{\mathrm{}}+\mathrm{}+\lambda _p^{\mathrm{}})\text{d}\mu _n^{(p)}}`$ $`=1+p\underset{n\mathrm{}}{lim}{\displaystyle \lambda _1^{\mathrm{}}\text{d}\mu _n^{(1)}}=1+p{\displaystyle \lambda _1^{\mathrm{}}\text{d}\nu ^{(1)}}=`$ $`={\displaystyle (1+\lambda _1^{\mathrm{}}+\lambda _2^{\mathrm{}}+\mathrm{}+\lambda _p^{\mathrm{}})\text{d}\nu ^{(p)}}.`$ Now (53) follows from Claim 5.3.2, and the lemma is proved. ∎ ###### Theorem 5.4. Let $`c1`$ and let $`Y(n)`$ be an $`n\times n`$ random matrix whose density with respect to Lebesgue measure on $`M_n(𝐂)`$ is $$K_{c,n}^{(1)}\|detY\|^{2(c1)n}\mathrm{exp}\left(n\mathrm{Tr}(Y^{}Y)\right).$$ Then $`Y(n)`$ converges in $``$–moments as $`n\mathrm{}`$ to a circular free Poisson element of parameter $`c`$. ###### Proof. Clearly for every non–random $`n\times n`$ unitary matrix $`U`$, the distribution of $`UY(n)`$ is equal to the distribution of $`Y(n)`$. Let $`\sigma _n`$ be the symmetrized joint distribution of the eigenvalues of $`\left(Y(n)^{}Y(n)\right)^{1/2}`$ and let $`\mu _n`$ be the symmetrized joint distribution of the eigenvalues of $`Y(n)^{}Y(n)`$. For $`p\{1,\mathrm{},n\}`$ let $`\sigma _n^{(p)}`$, respectively $`\mu _n^{(p)}`$, be the marginal distribution of $`\sigma _n`$, respectively $`\mu _n`$, corresponding to the first $`p`$ variables. Given $`k_1,\mathrm{},k_p𝐍\{0\}`$, $$\lambda _1^{k_1}\lambda _2^{k_2}\mathrm{}\lambda _p^{k_p}\text{d}\sigma _n^{(p)}(\lambda _1,\mathrm{},\lambda _p)=\lambda _1^{k_1/2}\lambda _2^{k_2/2}\mathrm{}\lambda _p^{k_p/2}\text{d}\mu _n^{(p)}(\lambda _1,\mathrm{},\lambda _p).$$ By Theorem 5.1 and Lemma 5.3, it follows that $$\underset{n\mathrm{}}{lim}\lambda _1^{k_1}\lambda _2^{k_2}\mathrm{}\lambda _p^{k_p}\text{d}\sigma _n^{(p)}(\lambda _1,\mathrm{},\lambda _p)=\underset{i=1}{\overset{p}{}}\lambda _i^{k_i/2}\text{d}\nu _c(\lambda _i),$$ where $`\nu _c`$ is the free Poisson distribution of parameter $`c`$. Therefore $`\sigma _n^{(p)}`$ converges in moments to $`\underset{1}{\overset{p}{\times }}\rho `$, where $`\rho `$ has density $$\frac{\text{d}\rho }{\text{d}t}=\frac{\sqrt{(d_1^2t^2)(t^2d_0^2)}}{\pi t}1_{[d_0,d_1]}(t),$$ with $`d_0=1\sqrt{c}`$ and $`d_1=1+\sqrt{c}`$. Now Theorem 4.6 applies and finishes the proof. ∎ Every complex $`n\times n`$ matrix $`A`$ is unitarily conjugate to an upper triangular matrix: $`A=USU^{}`$ where $`U`$ is unitary and the $`(i,j)`$th entry of $`S`$ is zero if $`i>j`$. If $`A`$ has $`n`$ distinct eigenvalues then the pair $`(U,S)`$ is unique up to replacement by $`(UD,D^{}SD)`$, where $`D`$ is a diagonal unitary. Given a random matrix $`X_n`$, one may ask for a corresponding random upper triangular matrix $`S`$ and random unitary matrix $`U`$ so that the distribution of $`USU^{}`$ is equal to the distribution of $`X`$. Then $`X`$ and $`S`$ will have the same $``$–moments with respect to the functional $`\tau _n`$. For specificity, we may insist that the joint distribution of the pair $`(U,S)`$ be the same as the joint distribution of $`(UD,D^{}SD)`$ for every non–random diagonal unitary $`D`$, (i.e. that the joint distribution of $`(U,S)`$ be invairant under this action of the $`n`$–torus $`𝐓^n`$). If the distribution of $`X`$ is invariant under conjugation by non–random unitaries and if $`(U,S)`$ is the pair of random matrices as described above, then it is clear that the random unitary $`U`$ is distributed according to Haar measure on the $`n\times n`$ unitaries and that $`U`$ and $`S`$ are independent. In this case, the relavant question is only the distribution of $`S`$. F. Dyson answered this question when $`X\text{GRM}(n,\frac{1}{n})`$. We state his result, and then make a slight modification to give, in conjunction with Theorem 5.4, an upper triangular matrix model for a circular free Poisson element. ###### Theorem 5.5 (Dyson, see \[9, A.35\]). Let $`T(n)\text{UTGRM}(n,\frac{1}{n})`$ and let $`D(n)_n`$ be a diagonal random matrix, whose diagonal entries have joint density $$K_n^{(7)}\mathrm{exp}\left(n\underset{i=1}{\overset{n}{}}\|z_i\|^2\right)\underset{1i<jn}{}\|z_iz_j\|^2$$ (54) with respect to Lebesgue measure on $`𝐂^n`$, for some constant $`K_n^{(7)}`$. Let $`U(n)\mathrm{HURM}(n)`$, suppose that $`(D(n),T(n),U(n))`$ is an independent family of matrix–valued random variables and let $$X(n)=U(n)\left(D(n)+T(n)\right)U(n)^{}.$$ Then $`X(n)\text{GRM}(n,\frac{1}{n})`$. Consequently, $`D(n)+T(n)`$ converges in $``$–moments as $`n\mathrm{}`$ to a circular element. ###### Corollary 5.6. Let $`c1`$, let $`T(n)\text{UTGRM}(n,\frac{1}{n})`$ and let $`D_c(n)_n`$ be a diagonal random matrix, whose diagonal entries have joint density $$K_{c,n}^{(8)}\mathrm{exp}\left(n\underset{i=1}{\overset{n}{}}\|z_i\|^2\right)\left(\underset{i=1}{\overset{n}{}}\|z_i\|\right)^{2(c1)n}\underset{1i<jn}{}\|z_iz_j\|^2$$ with respect to Lebesgue measure on $`𝐂^n`$, for some constant $`K_{c,n}^{(8)}`$. Let $`U(n)\mathrm{HURM}(n)`$, suppose that $`(D_c(n),T(n),U(n))`$ is an independent family of matrix–valued random variables and let $$Y(n)=U(n)\left(D_c(n)+T(n)\right)U(n)^{}.$$ (55) Then $`Y(n)`$ has density with respect to Lebesgue measure on $`M_n(𝐂)`$ equal to $$K_{c,n}^{(1)}\|detY\|^{2(c1)n}\mathrm{exp}\left(n\mathrm{Tr}(Y^{}Y)\right).$$ (56) Consequently, $`D_c(n)+T(n)`$ converges in $``$–moments as $`n\mathrm{}`$ to a circular free Poisson element of parameter $`c`$. ###### Proof. Let $`Mc_n`$ be the manifold of matrices in $`M_n(𝐂)`$ having $`n`$ distinct eigenvalues. Then $`Mc_n`$ has full Lebesgue measure in $`M_n(𝐂)`$. Let $`𝒰_n`$ be the Lie group of $`n\times n`$ unitary matrices, and let $`𝒯_n`$ be the manifold of all upper triangular $`n\times n`$ complex matrices, no two of whose diagonal elements are the same. Let $`\pi :𝒰_n\times 𝒯_nMc_n`$ be given by $`\pi (U,S)=USU^{}`$. Dyson proved his result by evaluating the Jacobian of $`\pi `$ (after throwing away the directions in $`\mathrm{ker}\text{d}\pi `$) and thereby finding the measure $`\sigma _n`$ on $`𝒯_n`$ such that letting $`\mu _n`$ be Haar measure on $`𝒰_n`$, the push–forward measure $`\pi _{}(\mu _n\times \sigma _n)`$ on $`Mc_n`$ has density $`K_{1,n}^{(1)}\mathrm{exp}\left(n\mathrm{Tr}(Y^{}Y)\right)`$ with respect to Lebesgue measure on $`Mc_n`$, i.e. the density of a random matrix $`X(n)\text{GRM}(n,\frac{1}{n})`$. This measure $`\sigma _n`$ was found to have density $$K_n^{(9)}\underset{1i<jn}{}\|S_{ii}S_{jj}\|^2\mathrm{exp}\left(n\mathrm{Tr}(S^{}S)\right)$$ (57) with respect to Lebesgue measure on $`𝒯_n`$, where for a matrix $`S𝒯_n`$, $`S_{ii}`$ is the $`i`$th diagonal entry of $`S`$; this density (57) is that of the matrix $`D(n)+T(n)`$ in Theorem 5.5. The matrix $`D_c(n)+T(n)`$ in the corollary has density $$K_{c,n}^{(10)}\underset{1i<jn}{}\|S_{ii}S_{jj}\|^2\mathrm{exp}\left(n\mathrm{Tr}(S^{}S)\right)\|det(S)\|^{2(c1)n}$$ with respect to Lebesgue measure on $`𝒯_n`$; since $`det(USU^{})=det(S)`$, and building on Dyson’s calculation, it follows that the random matrix $`Y(n)`$ of (55) has density (56) with respect to Lebesgue measure on $`Mc_n`$, as required. An application of Theorem 5.4 shows that $`Y(n)`$, and hence also $`D_c(n)+T(n)`$, converges in $``$–moments as $`n\mathrm{}`$ to a circular free Poisson element. ∎ The following lemma shows that the diagonal entries of $`D_c(n)`$ are in a specific sense asymptotically independent. This will allow their eventual decoupling; (see Remark 5.9). ###### Lemma 5.7. For $`c1`$ and $`n𝐍`$ let $`\mu _n`$ be the probability measure on $`𝐂^n`$ whose density with respect to Lebesgue measure is $$D_n(z_1,\mathrm{},z_n)=K_{c,n}^{(8)}\mathrm{exp}\left(n\underset{i=1}{\overset{n}{}}\|z_i\|^2\right)\left(\underset{i=1}{\overset{n}{}}\|z_i\|\right)^{2(c1)n}\underset{1i<jn}{}\|z_iz_j\|^2.$$ Given $`p\{1,2,\mathrm{},n\}`$ let $`\mu _n^{(p)}`$ be the marginal distribution of $`\mu _n`$ corresponding to the first $`p`$ variables $`z_1,\mathrm{},z_p`$. Then for every $`p𝐍`$, $`\mu _n^{(p)}`$ converges in weak topology and in $``$–moments as $`n\mathrm{}`$ to the product measure $`\underset{1}{\overset{p}{\times }}\rho `$, where $`\rho `$ is uniform distribution on the annulus $`\{z𝐂\sqrt{c1}<\|z\|<\sqrt{c}\}`$. ###### Proof. This is quite similar to the proof of Lemma 5.3. Let $`\alpha =(c1)n`$. Consider first the case $`p=1`$. Let $`\psi _0,\psi _1,\psi _2,\mathrm{}`$ be the polynomials obtained via Gram–Schmidt orthonormalization of the sequence $`1,z,z^2,\mathrm{}`$ in $`L^2(𝐂,\|z\|^{2\alpha }e^{n\|z\|^2}\text{d}(\mathrm{Re}z)\text{d}(\mathrm{Im}z))`$. Then $$\psi _k(z)=\sqrt{\frac{n^{k+\alpha +1}}{\pi \mathrm{\Gamma }(k+\alpha +1)}}z^k.$$ Using the Vandermonde determinant we have $$\underset{1i<jn}{}(z_jz_i)=K_{c,n}^{(11)}det\left(\begin{array}{cccc}\psi _0(z_1)& \psi _0(z_2)& \mathrm{}& \psi _0(z_n)\\ \psi _1(z_1)& \psi _1(z_2)& \mathrm{}& \psi _1(z_n)\\ \mathrm{}& \mathrm{}& & \mathrm{}\\ \psi _{n1}(z_1)& \psi _{n1}(z_2)& \mathrm{}& \psi _{n1}(z_n)\end{array}\right)$$ for some constant $`K_{c,n}^{(11)}`$. Therefore $$D_n(z_1,\mathrm{},z_n)=K_{c,n}^{(12)}\left(\underset{i=1}{\overset{n}{}}\|z_i\|\right)^{2\alpha }\left\|\underset{\pi S_n}{}\text{sign}(\pi )\underset{i=1}{\overset{n}{}}\psi _{\pi (i)1}(z_i)\right\|^2\mathrm{exp}\left(n\underset{i=1}{\overset{n}{}}\|z_i\|^2\right).$$ Writing $$\underset{i=1}{\overset{n}{}}\psi _{\pi (i)1}(z_i)=\psi _{\pi (1)1}(z_1)\psi _{\pi (2)1}(z_2)\mathrm{}\psi _{\pi (n)1}(z_n)$$ and noting that as $`\pi `$ ranges over the permutation group $`S_n`$ these form an orthonormal family with respect to the measure $`\left(_{i=1}^n\|z_i\|\right)^{2\alpha }\mathrm{exp}\left(n_{i=1}^n\|z_i\|^2\right)`$ on $`𝐂^n`$, we find $`K_{c,n}^{(12)}=(n!)^1`$. Moreover, the density of $`\mu _n^{(1)}`$ with respect to Lebesgue measure on $`𝐂`$ is $`D_{n,1}(z)`$ $`={\displaystyle _{𝐂^{n1}}}D_n(z,z_2,\mathrm{},z_n)\text{d}(\mathrm{Re}z_1)\text{d}(\mathrm{Im}z_1)\mathrm{}\text{d}(\mathrm{Re}z_n)\text{d}(\mathrm{Im}z_n)`$ $`={\displaystyle \frac{1}{n}}\left({\displaystyle \underset{k=0}{\overset{n1}{}}}\|\psi _k(z)\|^2\right)\|z\|^{2\alpha }e^{n\|z\|^2}={\displaystyle \frac{n^\alpha }{\pi }}{\displaystyle \underset{k=0}{\overset{n1}{}}}{\displaystyle \frac{n^k\|z\|^{2k+2\alpha }}{\mathrm{\Gamma }(k+\alpha +1)}}e^{n\|z\|^2}.`$ We shall show that $`\mu _n^{(1)}`$ converges in $``$–moments to $`\rho `$. Clearly if $`a,b𝐍\{0\}`$ and if $`ab`$ then $$_𝐂z^a\overline{z}^bD_{n,1}(z)\text{d}(\mathrm{Re}z)\text{d}(\mathrm{Im}z)=0=_𝐂z^a\overline{z}^b\text{d}\rho (z).$$ Hence we need only show $$\underset{n\mathrm{}}{lim}_𝐂\|z\|^{2b}\text{d}\mu _n^{(1)}(z)=_𝐂\|z\|^{2b}\text{d}\rho (z)$$ for all $`b𝐍\{0\}`$. We have $`{\displaystyle _𝐂}\|z\|^{2b}\text{d}\mu _n^{(1)}(z)`$ $`={\displaystyle \frac{n^{(c1)n}}{\pi }}{\displaystyle \underset{k=0}{\overset{n1}{}}}{\displaystyle \frac{n^k}{\mathrm{\Gamma }(k+(c1)n+1)}}{\displaystyle _𝐂}\|z\|^{2(b+k+(c1)n)}e^{n\|z\|^2}\text{d}(\mathrm{Re}z)\text{d}(\mathrm{Im}z)`$ $`={\displaystyle \underset{k=0}{\overset{n1}{}}}{\displaystyle \frac{n^{k+(c1)n}}{\mathrm{\Gamma }(k+(c1)n+1)}}{\displaystyle _0^{\mathrm{}}}t^{b+k+(c1)n}e^{nt}\text{d}t.`$ Writing $$f_n(t)=\underset{k=0}{\overset{n1}{}}\frac{n^{k+(c1)n}}{\mathrm{\Gamma }(k+(c1)n+1)}t^{b+k+(c1)n}e^{nt},$$ we have $`{\displaystyle _𝐂}\|z\|^{2b}\text{d}\mu _n^{(1)}(z)`$ $`={\displaystyle _0^{\mathrm{}}}t^bf_n(t)\text{d}t={\displaystyle \frac{1}{b+1}}{\displaystyle _0^{\mathrm{}}}t^{b+1}f_n^{}(t)\text{d}t=`$ $`={\displaystyle \frac{1}{b+1}}\left({\displaystyle \frac{\mathrm{\Gamma }(cn+b+1)}{n^{b+1}\mathrm{\Gamma }(cn)}}{\displaystyle \frac{\mathrm{\Gamma }((c1)n+b+1)}{n^{b+1}\mathrm{\Gamma }((c1)n)}}\right)=`$ $`={\displaystyle \frac{1}{b+1}}\left(c(c+\frac{1}{n})\mathrm{}(c+\frac{b}{n})(c1)\left((c1)+\frac{1}{n}\right)\mathrm{}\left((c1)+\frac{b}{n}\right)\right)`$ $`\stackrel{n\mathrm{}}{}{\displaystyle \frac{1}{b+1}}\left(c^{b+1}(c1)^{b+1}\right)={\displaystyle _𝐂}\|z\|^{2b}\text{d}\rho (z).`$ Hence $`\mu _n^{(1)}`$ converges in $``$–moments to $`\rho `$ as $`n\mathrm{}`$; since $`\rho `$ is compactly supported it follows that $`\mu _n^{(1)}`$ converges in the weak topology to $`\rho `$. The density of $`\mu _n^{(2)}`$ with respect to Lebesgue measure is $`D_{n,2}(z_1,z_2)`$ $`={\displaystyle _{𝐂^{n2}}}D_n(z_1,z_2,\mathrm{},z_n)\text{d}(\mathrm{Re}z_3)\text{d}(\mathrm{Im}z_3)\mathrm{}\text{d}(\mathrm{Re}z_n)\text{d}(\mathrm{Im}z_n)`$ $`=\begin{array}{cc}\hfill {\displaystyle \frac{1}{n!}}{\displaystyle \underset{\pi ,\sigma S_n}{}}\text{sign}(\pi )\text{sign}(\sigma ){\displaystyle _{𝐂^{n2}}}{\displaystyle \underset{i=1}{\overset{n}{}}}(& \psi _{\pi (i)1}(z_i)\overline{\psi _{\sigma (i)1}(z_i)}\|z_i\|^{2\alpha }e^{n\|z_i\|^2})\hfill \\ & \text{d}(\mathrm{Re}z_3)\text{d}(\mathrm{Im}z_3)\mathrm{}\text{d}(\mathrm{Re}z_n)\text{d}(\mathrm{Im}z_n)\hfill \end{array}`$ $`={\displaystyle \frac{1}{n(n1)}}{\displaystyle \underset{\begin{array}{c}0k,\mathrm{}n1\\ k\mathrm{}\end{array}}{}}\begin{array}{cc}& (\|\psi _k(z_1)\|^2\|\psi _{\mathrm{}}(z_2)\|^2\psi _k(z_1)\psi _{\mathrm{}}(z_2)\overline{\psi _{\mathrm{}}(z_1)}\overline{\psi _k(z_2)})\hfill \\ & \|z_1\|^{2\alpha }\|z_2\|^{2\alpha }e^{n(\|z_1\|^2+\|z_2\|^2)}\hfill \end{array}`$ $`={\displaystyle \frac{n}{n1}}D_{n,1}(z_1)D_{n,1}(z_2){\displaystyle \frac{1}{n(n1)}}\left\|{\displaystyle \underset{k=0}{\overset{n1}{}}}\psi _k(z_1)\overline{\psi _k(z_2)}\right\|^2\|z_1\|^{2\alpha }\|z_2\|^{2\alpha }e^{n(\|z_1\|^2+\|z_2\|^2)}.`$ Hence as a linear functional on $`C_0(𝐂^2)`$, the norm of $`\mu _n^{(2)}\mu _n^{(1)}\mu _n^{(1)}`$ is bounded above by $`2/(n1)`$. Therefore $`\mu _n^{(2)}`$ converges in the weak topology as $`n\mathrm{}`$ to $`\rho \times \rho `$. Arguing as in the proof of Lemma 5.3 and using Theorem 5.2, we conclude that for every $`p1`$, $`\mu _n^{(p)}`$ converges in weak topology to $`\underset{1}{\overset{p}{\times }}\rho `$, which we will denote by $`\nu ^{(p)}`$. It remains to show that $`\mu _n^{(p)}`$ converges to $`\nu ^{(p)}`$ in $``$–moments, namely that $$\underset{n\mathrm{}}{lim}z_1^{k_1}\overline{z_1}^\mathrm{}_1\mathrm{}z_p^{k_p}\overline{z_p}^\mathrm{}_p\text{d}\mu _n^{(p)}(z_1,\mathrm{},z_p)=z_1^{k_1}\overline{z_1}^\mathrm{}_1\mathrm{}z_p^{k_p}\overline{z_p}^\mathrm{}_p\text{d}\nu ^{(p)}(z_1,\mathrm{},z_p)$$ (58) for every $`k_1,\mathrm{},k_p,\mathrm{}_1,\mathrm{},\mathrm{}_p𝐍\{0\}`$. Exactly as in the proof of Claim 5.3.1, one shows that if $`h`$ is a positive continuous function on $`𝐂^p`$ then $$\underset{n\mathrm{}}{lim\; inf}h\text{d}\mu _n^{(p)}h\text{d}\nu ^{(p)}.$$ Then, considering the real and imaginary parts separately and arguing as in the proof of Claim 5.3.2, one shows that if $`f`$ and $`g`$ are continuous functions on $`𝐂^p`$, if $`g0`$, if $`\|f\|g`$ and if $`lim_n\mathrm{}g\text{d}\mu _n^{(p)}=g\text{d}\nu ^{(p)}`$ then $`lim_n\mathrm{}f\text{d}\mu _n^{(p)}=f\text{d}\nu ^{(p)}`$. But letting $`m=k_1+\mathrm{}+k_p+\mathrm{}_1+\mathrm{}+\mathrm{}_p`$, we have $$\|z_1^{k_1}\overline{z_1}^\mathrm{}_1\mathrm{}z_p^{k_p}\overline{z_p}^\mathrm{}_p\|1+\|z_1\|^{2m}+\mathrm{}+\|z_p\|^{2m}.$$ Moreover, because $`\mu _n^{(1)}`$ converges in $``$–moments to $`\nu ^{(1)}`$, we have $`\underset{n\mathrm{}}{lim}{\displaystyle (1+\|z_1\|^{2m}+\mathrm{}+\|z_p\|^{2m})\text{d}\mu _n^{(p)}}`$ $`=1+p\underset{n\mathrm{}}{lim}{\displaystyle \|z_1\|^{2m}\text{d}\mu _n^{(1)}}=1+p{\displaystyle \|z_1\|^{2m}\text{d}\nu ^{(1)}}=`$ $`={\displaystyle (1+\|z_1\|^{2m}+\mathrm{}+\|z_p\|^{2m})\text{d}\nu ^{(p)}}.`$ Hence we have (58) and the lemma is proved. ∎ ###### Lemma 5.8. Let $`c1`$ and for every $`n𝐍`$ let $`\mu _n`$ and $`\mu _n^{}`$ be the probability measures on $`𝐂^n`$ whose densities with respect to Lebesgue measure are, respectively, $`D_n(z_1,\mathrm{},z_n)`$ $`=K_{c,n}^{(8)}\left({\displaystyle \underset{i=1}{\overset{n}{}}}\|z_i\|\right)^{2(c1)n}\left({\displaystyle \underset{1i<jn}{}}\|z_iz_j\|^2\right)\mathrm{exp}\left(n{\displaystyle \underset{i=1}{\overset{n}{}}}\|z_i\|^2\right)`$ (59) $`D_n^{}(z_1,\mathrm{},z_n)`$ $`=K_{c,n}^{(13)}\left({\displaystyle \underset{i=1}{\overset{n}{}}}\|z_i\|\right)^{2(c1)n}\left({\displaystyle \underset{\pi S_n}{}}{\displaystyle \underset{i=1}{\overset{n}{}}}\|z_i\|^{2(\pi (i)1)}\right)\mathrm{exp}\left(n{\displaystyle \underset{i=1}{\overset{n}{}}}\|z_i\|^2\right).`$ (60) For $`p\{1,\mathrm{},n\}`$ let $`\mu _n^{(p)}`$ and $`(\mu _n^{})^{(p)}`$ denote the marginal distributions of $`\mu _n`$ and, respectively, $`\mu _n^{}`$ corresponding to the variables $`z_1,\mathrm{},z_p`$. Then for every $`p`$, $`(\mu _n^{})^{(p)}`$ is obtained from $`\mu _n^{(p)}`$ by averaging over the action of the torus $`𝐓^p`$ on $`𝐂^p`$ given by coordinate–wise multiplication: $$𝐓^p\times 𝐂^p((w_1,\mathrm{},w_p),(z_1,\mathrm{},z_p))(w_1z_1,\mathrm{},w_pz_p).$$ Consequently, $`(\mu _n^{})^{(p)}`$ converges in $``$–moments and in weak topology as $`n\mathrm{}`$ to the measure $`\underset{1}{\overset{p}{\times }}\rho `$, where $`\rho `$ is the uniform distribution on the annulus $`\{z𝐂\sqrt{c1}<\|z\|<\sqrt{c}\}`$. ###### Proof. Using the Vandermonde determinant we find $$\underset{1i<jn}{}\|z_iz_j\|^2=\underset{\pi ,\sigma S_n}{}\text{sign}(\pi )\text{sign}(\sigma )\underset{i=1}{\overset{n}{}}z_i^{\pi (i)1}\overline{z_i}^{\sigma (i)1}.$$ (61) Averaging (61) over the action of $`𝐓^n`$ gives $$\underset{\pi S_n}{}\underset{i=1}{\overset{n}{}}\|z_i\|^{2(\pi (i)1)}.$$ From this one easily sees that$`K_{c,n}^{(8)}=K_{c,n}^{(13)}`$ and that $`\mu _n^{}`$ is obtained from $`\mu _n`$ by averaging. In order to show that $`(\mu _n^{})^{(p)}`$ is obtained from $`\mu _n^{(p)}`$ by averaging, it suffices to note that for any measure $`\tau `$ on $`𝐂^n`$, the average over the action of $`𝐓^p`$ on the marginal distribution, $`\tau ^{(p)}`$, corresponding to the first $`p`$ variables, is equal to the marginal distribution of the average over the action of $`𝐓^n`$ on $`\tau `$. Since, by Lemma 5.7, $`\mu _n^{(p)}`$ converges in $``$–moments and in weak topology as $`n\mathrm{}`$ to $`\underset{1}{\overset{p}{\times }}\rho `$, which is invariant under the action of $`𝐓^p`$, it follows that $`(\mu _n^{})^{(p)}`$ converges in $``$–moments and in weak topology to $`\underset{1}{\overset{p}{\times }}\rho `$. ∎ ###### Remark 5.9. Our main purpose in proving the immediately preceding two lemmas was to be able to conclude that $$\underset{n\mathrm{}}{lim}\left(z_1^{k_1}\overline{z_1}^\mathrm{}_1\mathrm{}z_p^{k_p}\overline{z_p}^\mathrm{}_p\text{d}\mu _n^{(p)}z_1^{k_1}\overline{z_1}^\mathrm{}_1\mathrm{}z_p^{k_p}\overline{z_p}^\mathrm{}_p\text{d}(\mu _n^{})^{(p)}\right)=0$$ (62) for every $`p𝐍`$ and every $`k_1,\mathrm{},k_p,\mathrm{}_1,\mathrm{},\mathrm{}_p𝐍\{0\}`$. It is possible to prove (62) directly using the Vandermonde determinant and combinatorial arguments, though this sort of proof is not as satisfying as the one above involving Lemma 5.7, where the limit measure is found. ###### Theorem 5.10. Let $`c1`$ and $`N𝐍`$, and let $`(A,\varphi )`$ be a W–noncommutative probability space with random variables $`a_1,\mathrm{},a_NA`$ and $`b_{ij}A`$ ($`1i<jN`$), where $`a_j`$ is a circular free Poisson element of parameter $`(c1)N+j`$, where each $`b_{ij}`$ is a circular element with $`\varphi (b_{ij}^{}b_{ij})=1`$, and where the family $$((\{a_j^{},a_j\})_{1jN},(\{b_{ij}^{},b_{ij}\})_{1i<jN})$$ is free. Consider the W–noncommutative probability space $`(M_N(A),\varphi _N)`$, where $$\varphi _N\left((x_{ij})_{1i,jN}\right)=N^1\underset{j=1}{\overset{N}{}}\varphi (x_{jj}),$$ and consider the random variable $$x=\frac{1}{\sqrt{N}}\left(\begin{array}{cccccc}a_1& b_{12}& b_{13}& \mathrm{}& b_{1,N1}& b_{1N}\\ 0& a_2& b_{23}& \mathrm{}& b_{2,N1}& b_{2N}\\ 0& 0& a_3& \mathrm{}& \mathrm{}& b_{3N}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& \mathrm{}& 0& a_{N1}& b_{N1,N}\\ 0& 0& \mathrm{}& 0& 0& a_N\end{array}\right)M_N(A).$$ Then $`x`$ is a circular free Poisson element of parameter $`c`$. ###### Proof. For every $`n𝐍`$ let $`Y(nN)`$ be an $`nN\times nN`$ random matrix whose distribution has density with respect to Lebesgue measure $$K_{c,nN}^{(1)}\|det(Y)\|^{2(c1)nN}\mathrm{exp}\left(nN\mathrm{Tr}(Y^{}Y)\right).$$ Then by Theorem 5.4, $`Y(nN)`$ converges in $``$–moments as $`n\mathrm{}`$ to a circular free Poisson element of parameter $`c`$. By Corollary 5.6, each $`Y(nN)`$ has the same $``$–moments as $`S^{(1)}(nN)\stackrel{\text{def}}{=}D^{(1)}(nN)+T(nN)`$, where $`T(nN)\text{UTGRM}(nN,\frac{1}{nN})`$, $`D^{(1)}(nN)`$ is a diagonal $`nN\times nN`$ random matrix, the distribution of whose diagonal entries has density $$K_{c,nN}^{(8)}\left(\underset{i=1}{\overset{nN}{}}\|z_i\|\right)^{2(c1)nN}\left(\underset{1i<jnN}{}\|z_iz_j\|^2\right)\mathrm{exp}\left(nN\underset{i=1}{\overset{n}{}}\|z_i\|^2\right)$$ (63) with respect to Lebesgue measure on $`𝐂^n`$, and where $`D^{(1)}(nN)`$ and $`T(nN)`$ are independent. We will use previous results to show that also each $`S^{(k)}(nN)\stackrel{\text{def}}{=}D^{(k)}(nN)+T(nN)`$, ($`k\{2,3,4,5,6\}`$), converges in $``$–moments as $`n\mathrm{}`$ to a circular free Poisson element of parameter $`c`$, where $`D^{(k)}(nN)`$ is a diagonal random matrix such that $`D^{(k)}(nN)`$ and $`T(nN)`$ are independent and where the joint distributions of the diagonal entries of $`D^{(k)}(nN)`$ have the following densities with respect to Lebesgue measure on $`𝐂^n`$: $`\text{for }k=2:`$ $`K_{c,nN}^{(13)}\left({\displaystyle \underset{i=1}{\overset{nN}{}}}\|z_i\|\right)^{2(c1)nN}\left({\displaystyle \underset{\pi S_{nN}}{}}{\displaystyle \underset{i=1}{\overset{nN}{}}}\|z_i\|^{2(\pi (i)1)}\right)\mathrm{exp}\left(nN{\displaystyle \underset{i=1}{\overset{nN}{}}}\|z_i\|^2\right)`$ $`\text{for }k=3:`$ $`K_{c,nN}^{(14)}\left({\displaystyle \underset{i=1}{\overset{nN}{}}}\|z_i\|\right)^{2(c1)nN}\left({\displaystyle \underset{i=1}{\overset{nN}{}}}\|z_i\|^{2(i1)}\right)\mathrm{exp}\left(nN{\displaystyle \underset{i=1}{\overset{nN}{}}}\|z_i\|^2\right)`$ $`\text{for }k=4:`$ $`\begin{array}{cc}\hfill K_{c,nN}^{(15)}{\displaystyle \underset{j=1}{\overset{N}{}}}(({\displaystyle \underset{i=1}{\overset{n}{}}}\|z_{(j1)n+i}\|)^{2((c1)N+j1)n}& ({\displaystyle \underset{i=1}{\overset{n}{}}}\|z_{(j1)n+i}\|^{2(i1)})\hfill \\ & \mathrm{exp}(nN{\displaystyle \underset{i=1}{\overset{n}{}}}\|z_{(j1)n+i}\|^2))\hfill \end{array}`$ $`\text{for }k=5:`$ $`\begin{array}{cc}\hfill K_{c,nN}^{(16)}{\displaystyle \underset{j=1}{\overset{N}{}}}(({\displaystyle \underset{i=1}{\overset{n}{}}}\|z_{(j1)n+i}\|)^{2((c1)N+j1)n}& ({\displaystyle \underset{\pi S_n}{}}{\displaystyle \underset{i=1}{\overset{n}{}}}\|z_{(j1)n+i}\|^{2(\pi (i)1)})\hfill \\ & \mathrm{exp}(nN{\displaystyle \underset{i=1}{\overset{n}{}}}\|z_{(j1)n+i}\|^2))\hfill \end{array}`$ $`\text{for }k=6:`$ $`\begin{array}{cc}\hfill K_{c,nN}^{(17)}{\displaystyle \underset{j=1}{\overset{N}{}}}(({\displaystyle \underset{i=1}{\overset{n}{}}}\|z_{(j1)n+i}\|)^{2((c1)N+j1)n}& ({\displaystyle \underset{1i<i^{}n}{}}\|z_{(j1)n+i}z_{(j1)n+i^{}}\|^2)\hfill \\ & \mathrm{exp}(nN{\displaystyle \underset{i=1}{\overset{n}{}}}\|z_{(j1)n+i}\|^2)).\hfill \end{array}`$ The proof that $`S^{(k)}(nN)`$ converges in $``$–distribution to a circular free Poisson element relies for $`k=2`$ on Lemmas 5.7 and 5.8, (see Remark 5.9), and Theorem 3.2; for $`k=3`$ we use Theorem 3.6; the density for $`k=4`$ is just a rewriting of that for $`k=3`$; for $`k=5`$ we use again Theorem 3.6; for $`k=6`$ we use again Lemmas 5.7 and 5.8, and Theorem 3.2. We may characterize the above successive transformations as follows: from (63) to $`k=2`$ is decoupling; from $`k=2`$ to $`k=3`$ is desymmetrization; from $`k=3`$ to $`k=4`$ is regrouping; from $`k=4`$ to $`k=5`$ is partial resymmetrization; from $`k=5`$ to $`k=6`$ is partial recoupling. Taking blocks of consequetive rows and columns to write $`D^{(6)}(nN)+T(nN)`$ as an $`N\times N`$ matrix of $`n\times n`$ random matrices, we have $$S^{(6)}(nN)=\frac{1}{\sqrt{N}}\left(\begin{array}{cccccc}A_1^{(6)}& B_{12}^{(6)}& B_{13}^{(6)}& \mathrm{}& B_{1,N1}^{(6)}& B_{1N}^{(6)}\\ 0& A_2^{(6)}& B_{23}^{(6)}& \mathrm{}& B_{2,N1}^{(6)}& B_{2N}^{(6)}\\ 0& 0& A_3^{(6)}& \mathrm{}& \mathrm{}& B_{3N}^{(6)}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& \mathrm{}& 0& A_{N1}^{(6)}& B_{N1,N}^{(6)}\\ 0& 0& \mathrm{}& 0& 0& A_N^{(6)}\end{array}\right),$$ where $$(\left(A_j^{(6)}(n)\right)_{1jN},\left(B_{ij}^{(6)}(n)\right)_{1i<jN})$$ is an independent family of matrix–valued random variables, where $`B_{ij}^{(6)}(n)\text{GRM}(n,\frac{1}{n})`$ for every $`1i<jN`$ and where $`A_j^{(6)}(n)=D_j^{(6)}(n)+T_j(n)`$ with $`T_j(n)\text{UTGRM}(n,\frac{1}{n})`$, with $`D_j^{(6)}(n)`$ a diagonal random matrix, the joint distribution of whose diagonal entries has density with respect to Lebesgue measure $$K_{(c1)N+j,n}^{(8)}\left(\underset{i=1}{\overset{n}{}}\|z_i\|\right)^{2((c1)N+j1)n}\left(\underset{1i<i^{}n}{\overset{n}{}}\|z_iz_i^{}\|^2\right)\mathrm{exp}\left(n\underset{i=1}{\overset{n}{}}\|z_i\|^2\right)$$ and with $`D_j^{(6)}(n)`$ and $`T_j(n)`$ independent. Let $`U_j(n)\mathrm{HURM}(n)`$, ($`1jN`$), be such that $$(\left(A_j^{(6)}(n)\right)_{1jN},\left(B_{ij}^{(6)}(n)\right)_{1i<jN},\left(U_j(n)\right)_{1jN})$$ is an independent family of matrix–valued random variables. By conjugating the matrix $`S^{(6)}(nN)`$ with $`\text{diag}(U_1(n),U_2(n),\mathrm{},U_N(n))`$ and by using Corollary 5.6 and the fact that the class $`\text{GRM}(n,1/n)`$ is invariant under left and right multiplication by independent unitaries, it follows that $`S^{(6)}(nN)`$ has the same $``$–moments, as $$S^{(7)}(nN)=\frac{1}{\sqrt{N}}\left(\begin{array}{cccccc}A_1^{(7)}& B_{12}^{(7)}& B_{13}^{(7)}& \mathrm{}& B_{1,N1}^{(7)}& B_{1N}^{(7)}\\ 0& A_2^{(7)}& B_{23}^{(7)}& \mathrm{}& B_{2,N1}^{(7)}& B_{2N}^{(7)}\\ 0& 0& A_3^{(7)}& \mathrm{}& \mathrm{}& B_{3N}^{(7)}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& \mathrm{}& 0& A_{N1}^{(7)}& B_{N1,N}^{(7)}\\ 0& 0& \mathrm{}& 0& 0& A_N^{(7)}\end{array}\right),$$ where $$(\left(A_j^{(7)}(n)\right)_{1jN},\left(B_{ij}^{(7)}(n)\right)_{1i<jN})$$ (64) is an independent family of matrix–valued random variables, where $`B_{ij}^{(7)}(n)\text{GRM}(n,\frac{1}{n})`$ for every $`1i<jN`$ and where the distribution of $`A_j^{(7)}(n)`$ has density $$K_{(c1)N+j,n}^{(1)}\|det(A)\|^{2((c1)N+j1)n}\mathrm{exp}(n\mathrm{Tr}(A^{}A))$$ with respect to Lebesgue measure on $`M_n(𝐂)`$. If $`(V_j^{(1)})_{1jN}`$, $`(V_j^{(2)})_{1jN}`$, $`(U_{ij}^{(1)})_{1i<jN}`$, $`(U_{ij}^{(2)})_{1i<jN}`$, are non–random $`n\times n`$ unitary matrices, then $$(\left(V_j^{(1)}A_j^{(7)}(n)V_j^{(2)}\right)_{1jN},\left(U_{ij}^{(1)}B_{ij}^{(7)}(n)U_{ij}^{(2)}\right)_{1i<jN})$$ (65) continues to be an independent family of matrix–valued random variables, $`V_j^{(1)}A_j^{(7)}(n)V_j^{(2)}`$ has the same distribution as $`A_j^{(7)}`$ and $`U_{ij}^{(1)}B_{ij}^{(7)}(n)U_{ij}^{(2)}`$ has the same distribution as $`B_{ij}^{(7)}(n)`$. Therefore, the family (65) has the same joint $``$–moments as the family (64). Taking into account also Theorem 5.1 and Lemma 5.3 (as in the proof of Theorem 5.4), we see that the conditions of Theorem 4.6 are fulfilled, allowing us to conclude that the family (64) is asymptotically $``$–free as $`n\mathrm{}`$. Moreover, (by Theorem 5.4), each $`A_j^{(7)}(n)`$ converges in $``$–moments to a circular free Poisson element of parameter $`(c1)N+j`$, while $`B_{ij}^{(7)}(n)`$ converges in $``$–moments to a circular element. Therefore, the entries of the matrix $`S^{(7)}(n)`$ model as $`n\mathrm{}`$ the entries of the matrix $`x`$ in the statement of the theorem. As $`S^{(7)}(n)`$ converges in $``$–moments to a circular free Poisson element of parameter $`c`$, the theorem is proved. ∎ ## 6. Invariant subspaces for a circular free Poisson element In this section, we will apply Theorem 5.10 and the general results of §2 to exhibit invariant subspaces for a circular free Poisson element. We will rely on the result of Haagerup and Larsen \[6, Example 5.2\] that the spectrum of a circular free Poisson element of parameter $`c`$ is $`\{z𝐂\sqrt{c1}\|z\|\sqrt{c}\}`$. ###### Theorem 6.1. Let $`(Mc,\psi )`$ be a W–noncommutative probability space with $`\psi `$ faithful, let $`c1`$ and let $`yMc`$ be a circular free Poisson element of parameter $`c`$. Given $`r0`$, let $`p_r(y)Mc`$ be the projection onto the invariant subspace of $`y`$ as in Definition 2.4. Then $$\psi (p_r(y))=\{\begin{array}{cc}0\hfill & \text{if }r\sqrt{c1}\hfill \\ r^2(c1)\hfill & \text{if }\sqrt{c1}r\sqrt{c}\hfill \\ 1\hfill & \text{if }r\sqrt{c}.\hfill \end{array}$$ (66) ###### Proof. We may without loss of generality assume that $`Mc=\{y\}^{\prime \prime }`$, which implies $`McL(F_2)`$ and $`\psi `$ is a trace. Let $`N𝐍`$ and let $$x=\frac{1}{\sqrt{N}}\left(\begin{array}{cccccc}a_1& b_{12}& b_{13}& \mathrm{}& b_{1,N1}& b_{1N}\\ 0& a_2& b_{23}& \mathrm{}& b_{2,N1}& b_{2N}\\ 0& 0& a_3& \mathrm{}& \mathrm{}& b_{3N}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& \mathrm{}& 0& a_{N1}& b_{N1,N}\\ 0& 0& \mathrm{}& 0& 0& a_N\end{array}\right)M_N(A).$$ be the circular free Poisson element of parameter $`c`$ as in Theorem 5.10, where we take $`(A,\varphi )`$ to be a W–noncommutative probability space. Thus $`a_j`$ is a circular free Poisson element of parameter $`(c1)N+j`$, each $`b_{ij}`$ is a circular element and the collection of all $`a_j`$ and $`b_{ij}`$ is $``$–free. For $`k\{1,\mathrm{},N\}`$, let $$e_k=\text{diag}(\underset{k\text{ times}}{\underset{}{1,\mathrm{},1}},0,\mathrm{},0)M_n(𝐂1)M_n(A).$$ Another application of Theorem 5.10 shows that in the W–noncommutative probability space $$(e_kM_N(A)e_k,\sqrt{\frac{N}{k}}\varphi _N_{e_kM_N(A)e_k})(M_k(A),\varphi _k),$$ the element $`\sqrt{\frac{N}{k}}e_kxe_k`$ is a circular free Poisson element of parameter $`\frac{N}{k}(c1)+1`$. Hence by \[6, Example 5.2\], $`e_kxe_k`$ has spectrum $$\left\{z𝐂\right\|\sqrt{c1}\|z\|\sqrt{(c1)+\frac{k}{N}}\}.$$ Similarly, if $`k<N`$ then denoting by $`1_N`$ the identity element of $`M_N(A)`$, we find that in the W–noncommutative probability space $$((1_Ne_k)M_N(A)(1_Ne_k),\sqrt{\frac{N}{Nk}}\mathrm{\Phi }_N_{(1_Ne_k)M_N(A)(1_Ne_k)})(M_{Nk}(A),\varphi _{Nk})$$ the element $`\sqrt{\frac{N}{Nk}}(1_Ne_k)x(1_Ne_k)`$ is a circular free Poisson element of parameter $`\frac{N}{Nk}c`$. Hence $`(1_Ne_k)x(1_Ne_k)`$ has spectrum $$\left\{z𝐂\right\|\sqrt{(c1)+\frac{k}{N}}\|z\|\sqrt{c}\}.$$ Therefore, by Proposition 2.2, if $`kN2`$ and if $$\sqrt{(c1)+\frac{k}{N}}r<\sqrt{(c1)+\frac{k+1}{N}}$$ then $`e_kp_r(x)e_{k+1}`$ and consequently $`\frac{k}{N}\psi (p_r(y))\frac{k+1}{N}`$. Letting $`N`$ grow without bound and choosing $`k`$ appropriately implies (66). ∎ Some further facts concerning these projections $`p_r(y)`$ are collected below in Theorem 6.3, for the proof of which we will use the following lemma. ###### Lemma 6.2. For every $`c1`$ let $`(A_c,\varphi _c)`$ be a W–noncommutative probability space and let $`y_cA`$ be a circular free Poisson element of parameter $`c`$. Given $`c_01`$ and a sequence $`(c_n)_1^{\mathrm{}}`$ in $`[1,\mathrm{})`$ converging to $`c_0`$, we have that $`y_{c_n}`$ converges in $``$–moments to $`y_{c_0}`$ as $`n\mathrm{}`$. ###### Proof. The positive part $`h_c`$ of $`y_c`$ has the same moments as the measure $`\rho _c`$ on $`𝐑_+`$ whose density with respect to Lebesgue measure is $$\frac{\text{d}\rho _c}{\text{d}t}=\frac{\sqrt{(d_1^2t^2)(t^2d_0^2)}}{\pi t}1_{[d_0,d_1]}(t),$$ with $`d_0=1\sqrt{c}`$ and $`d_1=1+\sqrt{c}`$. Since $`y_c`$ has the polar decomposition $`y_c=u_ch_c`$ where $`u_c`$ is a Haar unitary and where $`u_c`$ and $`h_c`$ are $``$–free, the $``$–moments of $`y_c`$ can be expressed as certain polynomials in the moments of $`\rho _c`$. Clearly, the $`k`$th moment of $`\rho _{c_n}`$ converges to the $`k`$th moment of $`\rho _{c_0}`$ as $`n\mathrm{}`$. ∎ ###### Theorem 6.3. Let $`y`$ be a circular free Poisson element of parameter $`c`$ in some W–noncommutative probability space $`(Mc,\psi )`$, with $`\psi `$ faithful. Then 1. $`p_s(y)`$ converges in the strong topology to $`p_r(y)`$ as $`sr`$. 2. If $`\sqrt{c1}<r<\sqrt{c}`$ then in the W–noncommutative probability space $$(p_r(y)Mcp_r(y),\frac{1}{\psi (p_r(y))}\psi _{p_r(y)Mcp_r(y)}),$$ (67) $`\psi (p_r(y))^{1/2}yp_r(y)`$ is a circular free Poisson element of parameter $`1+(c1)/\psi (p_r(y))`$. Hence the spectrum of $`yp_r(y)`$ relative to $`p_r(y)Mcp_r(y)`$ is $$\sigma (yp_r(y))=\{z𝐂\sqrt{c1}\|z\|\sqrt{r}\}.$$ 3. If $`\sqrt{c1}<r<\sqrt{c}`$ then in the W–noncommutative probability space $$((1p_r(y))Mc(1p_r(y)),\frac{1}{1\psi (p_r(y))}\psi _{(1p_r(y))Mc(1p_r(y))}),$$ $`\left(1\psi (p_r(y))\right)^{1/2}(1p_r(y))y`$ is a circular free Poisson element of parameter $`c/\left(1\psi (p_r(y))\right)`$. Hence the spectrum of $`(1p_r(y))y`$ relative to $`(1p_r(y))Mc(1p_r(y))`$ is $$\sigma ((1p_r(y))y)=\{z𝐂\sqrt{r}\|z\|\sqrt{c}\}.$$ 4. If $`x`$ is the upper triangular $`N\times N`$ matrix given in Theorem 5.10 that is a circular free Poisson element of parameter $`c`$, then for every $`k\{0,1,,\mathrm{},N\}`$ and letting $`r=\sqrt{(c1)+k/N}`$, we have $`p_r(x)=\text{diag}(\underset{k\text{ times}}{\underset{}{1,\mathrm{},1}},0,\mathrm{},0)`$. ###### Proof. We know from general principles that $`p_r^{}(y)p_r(y)`$ if $`r^{}<r`$, and from Theorem 6.1 we have that $`lim_{sr}\psi (p_s(y))=\psi (p_r(y))`$; as $`\psi `$ is faithful we conclude (i). Let us now prove (iv). Arguing as in the proof of Theorem 6.1, we have $`e_kp_r(x)`$ whenever $`r>\sqrt{(c1)+k/N}`$. We may take the W-noncommutative probability space $`(A,\varphi )`$ so that $`\varphi `$ is a faithful trace, in which case, since $`inf\{\psi (p_r(x))r>\sqrt{(c1)+k/N}\}=k/N`$, it follows that $$e_k=\left\{p_r(x)\right\|r>\sqrt{(c1)+k/N}\}.$$ Thus $`e_k`$ is the limit in strong topology of $`p_r(x)`$ as $`r`$ tends to $`\sqrt{(c1)+k/N}`$ from above. Using (i), it follows that $`e_k=p_{\sqrt{(c1)+k/N}}(x)`$. For (ii), let us show that $`\psi (p_r(y))^{1/2}yp_r(y)`$ is circular free Poisson of the desired parameter, first in the case when $`\psi (p_r(y))=k/N`$ is rational. We may take $`(Mc,\psi )`$ to be $`(M_N(A),\varphi _N)`$ and $`y`$ to be equal to the $`N\times N`$ matrix $`x`$ as in Theorem 5.10. By (iv), the noncommutative probability space (67) is $`(M_k(A),\varphi _k)`$ and $$\frac{1}{\sqrt{\psi (p_r(y))}}yp_r(y)=\frac{1}{\sqrt{k}}\left(\begin{array}{cccccc}a_1& b_{12}& b_{13}& \mathrm{}& b_{1,k1}& b_{1k}\\ 0& a_2& b_{23}& \mathrm{}& b_{2,k1}& b_{2k}\\ 0& 0& a_3& \mathrm{}& \mathrm{}& b_{3k}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& \mathrm{}& 0& a_{k1}& b_{k1,k}\\ 0& 0& \mathrm{}& 0& 0& a_k\end{array}\right),$$ where $`a_j`$ is circular free Poisson of parameter $`(c1)N+j`$. Applying again Theorem 5.10, we obtain that $`\psi (p_r(y))^{1/2}yp_r(y)`$ is circular free Poisson of parameter $`1+(c1)/\psi (p_r(y))`$. When $`r`$ is such that $`\psi (p_r(y))`$ is irrational, then using (i) we have that $`yp_r(y)`$ is the strong limit of $`yp_s(y)`$ as $`s`$ tends to $`r`$ through rational numbers. Hence by Lemma 6.2 and the continuity in $`r`$ of of $`\psi (p_r(y))`$ implied by Theorem 5.10, it follows that $`\psi (p_r(y))^{1/2}yp_r(y)`$ is circular free Poisson of parameter $`1+(c1)/\psi (p_r(y))`$. The statement about the spectrum follows from the result of Haagerup and Larsen that we’ve been using repeatedly. Part (iii) is proved similarly. When $`\psi (p_r(y))=k/N`$ is rational then we get $$\frac{1}{\sqrt{1\psi (p_r(y))}}(1p_r(y))y=\frac{1}{\sqrt{Nk}}\left(\begin{array}{cccccc}a_{k+1}& b_{k+1,k+2}& b_{k+1,k+3}& \mathrm{}& b_{k+1,N1}& b_{k+1,N}\\ 0& a_{k+2}& b_{k+2,k+3}& \mathrm{}& b_{k+2,N1}& b_{k+2,N}\\ 0& 0& a_{k+3}& \mathrm{}& \mathrm{}& b_{k+3,N}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& \mathrm{}& 0& a_{N1}& b_{N1,N}\\ 0& 0& \mathrm{}& 0& 0& a_N\end{array}\right),$$ where $`a_{k+j}`$ is circular free Poisson of parameter $`(c1)N+k+j`$. The remaining part of the argument is like for (ii) above. ∎ The following proposition shows that $`p_r(y)`$ is characterized by the spectral conditions in (ii) and (iii) of Theorem 6.3. ###### Proposition 6.4. Let $`Y`$ be a circular free Poisson element of parameter $`c1`$ in a W–probability space $`(Mc,\psi )`$, with $`\psi `$ faithful, and let $`\sqrt{c1}<r<\sqrt{c}`$. Suppose $`pMc`$ is a projection such that 1. $`yp=pyp`$ 2. $`\sigma _{pMcp}(yp)\{z𝐂\sqrt{c1}\|z\|r\}`$ 3. $`\sigma _{(1p)Mc(1p)}((1p)y)\{z𝐂r\|z\|\sqrt{c}\}`$. Then $`p=p_r(y)`$. ###### Proof. Note that (i) implies $`(1p)y=(1p)y(1p)`$. Let $`Mc`$ be normally and faithfully represented on a Hilbert space $``$. For $`\xi p`$ we have $$\underset{n\mathrm{}}{lim\; sup}y^n\xi ^{1/n}\underset{n\mathrm{}}{lim\; sup}(pyp)^n^{1/n}r,$$ where the last inequality is because the spectral radius of $`pyp`$ is $`r`$. Hence $`pp_r(y)`$. In order to prove the reverse inequality, it will suffice to show $`pp_s(y)`$ for all $`0s<r`$, because $`sp_s(y)`$ is strong–continuous by Theorem 6.3(i). Let $`0s<r`$, $`\xi (1p)`$ and let $`\eta E_s(y)`$, i.e. $$\underset{n\mathrm{}}{lim\; sup}y^n\eta ^{1/n}s.$$ (68) Set $`\xi _n=\left((1p)y^{}(1p)\right)^n\xi `$. Then $`\xi _n(1p)`$ and, because the spectral radius of $`\left((1p)y^{}(1p)\right)^1`$ is $`1/r`$, we have $`lim\; sup_n\mathrm{}\xi _n^{1/n}1/r`$. Since $`(1p)`$ is an invariant subspace for $`y^{}`$, we have $`\xi =\left((1p)y^{}(1p)\right)^n\xi _n=(y^{})^n\xi _n`$. Therefore $`\xi ,\eta =(y^{})^n\xi _n,\eta =\xi _n,y^n\eta `$, so using (68) and Schwarz’s inequality we have $`lim\; sup_n\mathrm{}\|\xi ,\eta \|^{1/n}s/r<1`$, which shows that $`\xi ,\eta =0`$. Hence $`(1p)\overline{E_s(y)}=p_s(y)`$ and therefore $`p_s(y)p`$. ∎ ###### Remark 6.5. The proof above shows that the subspace $$E_r(y)=\{\xi \underset{k\mathrm{}}{lim\; sup}y^k\xi ^{1/k}r\}$$ is closed; thus we have $`p_r(y)=E_r(y)`$, without taking the closure. The next example, however, shows that the sort of spectral decomposition found in Theorem 6.3 and closedness of the subspace $`E_r(y)`$ do not always hold. ###### Example 6.6. Let $`=_{k=2}^{\mathrm{}}_k`$, where $`_k`$ is $`k`$–dimensional Hilbert space with orthonormal basis $`e_1^{(k)},\mathrm{},e_k^{(k)}`$ and let $`T=_{k=2}^{\mathrm{}}T_k`$, where $`T_kB(_k)`$ is the nilpotent operator $$T_ke_j^{(k)}=\{\begin{array}{cc}0\hfill & j=1\hfill \\ e_{j1}^{(k)}\hfill & 2jk.\hfill \end{array}$$ It is well-known that the spectrum of $`T`$ is the closed unit disk $`\overline{𝐃}`$ — see for example Brown \[2, Example 4.10\]. However, if $`r>0`$ then $`E_r(T)`$ is dense in $``$, so $`p_r(T)=1`$ and the spectrum of $`Tp_r(T)`$ is $`\overline{𝐃}`$. Moreover, if $`0<r<1`$ then the vector $`_{k=2}^{\mathrm{}}\frac{1}{k}e_k^{(k)}`$ is not an element of $`E_r(T)`$; this shows that $`E_r(T)`$ is not closed. Note that $`T_{k=2}^{\mathrm{}}B(_k)`$, which is a finite von Neumann algebra. Ken Dykema Department of Mathematics, Texas A&M University, College Station TX 77843–3368, USA E-mail: Ken.Dykema@math.tamu.edu Internet URL: http://www.math.tamu.edu/$`\stackrel{~}{}`$Ken.Dykema/ Uffe Haagerup\* Department of Mathematics and Computer Science, University of Southern Denmark — Odense University, Campusvej 55, 5230 Odense M, Denmark E-mail: haagerup@imada.sdu.dk Internet URL: http://www.imada.sdu.dk/$`\stackrel{~}{}`$haagerup/ \*MaPhySto, Centre for Mathematical Physics and Stochastics, funded by a grant from The Danish National Research Foundation.
warning/0004/astro-ph0004137.html	ar5iv	text	# Constraining the Location of Microlensing Objects by using the Finite Source Effect in EAGLE events ## 1 Introduction Several groups have been engaged in gravitational microlensing observations toward the Large Magellanic Cloud (LMC) and the Small Magellanic Cloud (SMC) in order to investigate the nature of MAssive Compact Halo Objects (MACHOs) in the Galaxy’s halo. Until now, 8 microlensing event candidates have been found toward the LMC (Alcock et al. (1997)), and 2 candidates, including one possible binary lens event, have been found toward the SMC (Afonso et al. (1998); Afonso et al. (1999); Albrow et al. (1999); Alcock et al. 1999a ). The microlensing light curve informs us only of the event duration $`\widehat{t}`$, in which the mass, velocity and distance of the lensing object are degenerate, and thus it is difficult to determine the mass or the distance of the lensing object for each microlensing event. Therefore, the nature of MACHOs in the Galactic halo still remains unclear, and one cannot rule out the possibility that the microlensing events are not caused by MACHOs in the Galaxy’s halo but by unknown populations lying between the Galaxy and the Clouds (e.g., Zaritsky & Lin 1997; Zhao 1998), or by normal stars in the Clouds themselves (e.g.,Sahu 1994). In order to understand the nature of lensing objects, some additional information that partially or fully break the three-fold degeneracy is required. Over the past years several investigations have been made to extract additional information from special types of microlensing events. Examples of which are binary events and parallax events, for which the lens proper motion and lens distance can be constrained (e.g., Schneider & Weiss 1984; Alcock et al. (1995); Miyamoto & Yoshii 1995). A transit event, in which the lens transits the extended surface of the source star, is another example of such types of microlensing events (Gould 1992, 1994; Nemiroff & Wickramasinghe (1994); Witt & Mao (1994); Peng (1997)). When the lens passes over the surface of the source star, the point source approximation fails and the peak of light curve deviates from that for a point source. From such a light curve, one can derive the radius of the source star scaled by the Einstein ring radius projected onto the source plane. Since the radius of the source star can be estimated from its spectrum, one can obtain the proper motion of the lens for transit events. The proper motion of the lens will be of great use in investigating the lens location, as the expected proper motion distribution of objects in the Galaxy’s halo differs significantly from that in the Clouds. Unfortunately however, the probability for transit events is expected to be extremely small ( $`0.1\%`$) for normal microlensing events (Gould (1994)). The probability for transit events is relatively high for EAGLE events, in which an invisible star is amplified above the detection limit, because EAGLE events are guaranteed to have a small impact parameter (Wang & Turner (1996); Gould (1997); Nakamura & Nishi (1998)). Since there are many more faint stars below the detection limit than visible stars that are usually used for microlensing search, the EAGLE event rate is expected to be fairly high (Nakamura & Nishi (1998)). However EAGLE events have rarely been seen in the current observation programs in which the brightness of visible stars are measured based on a DoPHOT-type analysis. Recently however, a new CCD photometry method called ‘image subtraction’ has been developed (Alard (1998), 2000; Alcock et al. 1999b ; I.A. Bond 1999, private communication). This method, which compares an exposure frame directly with a reference frame, is much more powerful for detecting EAGLE events than the current monitoring of visible stars, and one can expect that plenty of EAGLE events will be detected with this method in the near future. If many EAGLE events are detected, one may be able to find several transit events because EAGLE events have small impact parameters. Therefore, a search for EAGLE events will possibly provide a new opportunity to measure the lens proper motion in transit events, and thereby to constrain strongly the lens location. Furthermore, statistics of transit EAGLE events alone may be useful for studying the nature of lensing objects. The fraction of transit events depends on the ratio of the physical radius of the source star $`R_{}`$ to the Einstein radius $`r_E`$ in source plane, i.e. $`R_{}/r_E`$. In the case where the lenses are the stars in the LMC itself (self-lensing), the Einstein radius $`r_E`$ in source plane is much smaller than in the case where the lenses are MACHOs in the Galactic halo. Thus the fraction of transit EAGLE events out of all EAGLE events for self-lensing in the LMC will be much larger than that for lensing due to MACHOs in the halo. Therefore the fraction of transit EAGLE events may be used to constrain the location and the nature of lens. EAGLE events will thus become one of the most powerful tools for investigating the nature of lensing objects. For these reasons, in this paper we investigate the properties of EAGLE events, and present a detailed calculation of the probability for transit EAGLE events. In section 2 we briefly summarize the basic equations, and in section 3 we estimate the event rate of EAGLEs. In section 4 we calculate the fraction of transit events, and in section 5 we derive the event duration of EAGLEs. Discussion on the observation strategy for EAGLE search is made in section 6. ## 2 Introductory description of EAGLE events ### 2.1 Basic equations for EAGLE events The amplification A of the microlensing event for point source approximation is given by $$A(u)=\frac{u^2+2}{u\sqrt{u^2+4}},$$ (1) where $$u=\frac{d}{R_E}=\frac{D_s}{D_d}\frac{d}{r_E},R_E=\sqrt{\frac{4GM}{c^2}\frac{(D_sD_d)D_d}{D_s}}.$$ (2) Here $`D_s`$ and $`D_d`$ are the distances from the observer to the source star, and to the lens object respectively. $`d`$ is the distance from the lens to the source projected onto the lens plane. $`R_E`$ and $`r_E=(D_s/D_d)R_E`$ are the Einstein radius and the projected Einstein radius in the source plane, respectively. The parameter $`u`$ denotes the impact parameter in the lens plane scaled by $`R_E`$. In an EAGLE event, an invisible star with magnitude $`V`$ is amplified beyond the observational detection limit $`V_{obs}`$. More precisely, to identify an EAGLE event, the source must become brighter than the EAGLE detection threshold $`V_{\mathrm{th}}`$, which is slightly brighter than the observational detection limit $`V_{obs}`$. Thus, for an EAGLE event, the amplification $`A`$ must be larger than the threshold amplitude $`A_T`$, which is written as: $$A_T=10^{0.4(VV_{\mathrm{th}})}$$ (3) We define the threshold impact parameter $`u_T`$ as the largest impact parameter for an EAGLE event with a $`V_{th}`$ magnitude source star. The value of $`u_T`$ is determined so that $`A(u_T)=A_T(V,V_{\mathrm{th}})`$. In the case of a point source and a small impact parameter, we can make the approximation that $`A1/u`$. The threshold impact parameter can be then be written as (Nakamura & Nishi (1998)): $$u_T=10^{0.4(VV_{\mathrm{th}})}$$ (4) However, in most of the cases we are interested in, the effect of a finite source is not negligible, and this approximation is not applicable. We thus calculate the threshold impact parameter $`u_T`$ by taking the finite-source effect into account. When the source size is finite, the amplification $`A`$ cannot be given by equation (1), but approximated by the following expression (Gould (1994)), $$A_{\mathrm{fin}}(u,u_{})=A(u)B(z),z\frac{u}{u_{}},u_{}\frac{D_d}{D_s}\frac{R_{}}{R_E}=\frac{R_{}}{r_E}$$ (5) Here $`u_{}`$ is the source radius scaled by $`r_E`$ ($`R_{}`$ is the physical radius of the source), and $`z`$ is the impact parameter in the source plane scaled by the radius of the source star $`R_{}`$. The function B(z) describes the ratio of the amplification of a finite source to the amplification of a point source, which is written as $$B(z)=\frac{z^2}{\pi }_0^{2\pi }𝑑\theta _0^{1/z}𝑑qq(1+q^2+2q\mathrm{cos}\theta )^{1/2}.$$ (6) We show the relation of the amplification and other parameters in figure Constraining the Location of Microlensing Objects by using the Finite Source Effect in EAGLE events . To compute the amplification using equations (5) and (6), one has to know the physical size of the source star that has a magnitude of $`V`$. In this paper, we approximate the relation between $`R_{}`$ and $`V`$ for a main sequence star (Allen (1973); Blaauw (1963)) as follows, $`R_{}(V)`$ $`=`$ $`10^{0.053V+1.22},(21V31)`$ (7) $`=`$ $`10^{0.125V+2.78},(15V21).`$ (8) Using equations (5), (6), (7) and (8), one can derive the threshold impact parameter $`u_T`$ from the condition that $`A_{\mathrm{th}}=A_{\mathrm{fin}}(u_T,u_{})`$. ### 2.2 Mass and luminosity function of source stars The EAGLE event rate depends heavily on the mass function and the luminosity function of source stars in the LMC and the SMC. Following Nakamura & Nishi (1998), we assume that a stellar initial mass function is in the power law form with index $`\alpha `$, that the star formation rates in the LMC have been constant for 5 Gyr and was zero before 5 Gyr ago. Assuming that the lifetime of the Sun is 10 Gyr, we obtain the present day mass function as $`\varphi (M)`$ $`=`$ $`CM^\alpha ,(M_l<M<1.4M_{})`$ (9) $`=`$ $`2({\displaystyle \frac{M}{M_{}}})^2CM^\alpha ,(1.4M_{}<M<M_u)`$ (10) where $`M_l`$ and $`M_u`$ are the lower and the upper mass limit of the stars, respectively. The initial mass function of the stars in the LMC are not so different from the Salpeter IMF $`\alpha =2.35`$ (Holtzman (1996)). However for field stars the slope may be as steep as $`\alpha 5.0`$ (Westerlund (1997)). In the following analysis we use the two extreme values of $`\alpha `$, 2.35 and 5.0. We note that the results in the following sections do not depend strongly on the shape of the mass function beyond $`M=1.4M_{}`$ because the contribution of low mass stars is much more significant than that of massive stars. The mass-luminosity relation of the main sequence stars is expressed by (Kippenhahn et al. (1990)) $$L(M)=L_{}(\frac{M}{M_{}})^3$$ (11) The Bolometric Correction ($`B.C.`$) can be approximated as (Schild et al (1971); Davis & Webb (1970); Johnson (1964)) $`B.C.`$ $`=`$ $`0.6V11.6,(13V19)`$ (12) $`=`$ $`0,(19V24)`$ (13) $`=`$ $`0.37V+8.9,(24V31)`$ (14) Then we can drive the $`V`$ band luminosity function of stars in the LMC using equations (9) to (14). $`\varphi _L(V)`$ $`=`$ $`C^{}10^{(\alpha +2)(V+B.C.)/(7.5)},(15<V<21)`$ (15) $`=`$ $`C^{\prime \prime }10^{\alpha (V+B.C.)/(7.5)},(21<V<31)`$ (16) In this paper we will use these mass and luminosity functions. ## 3 EAGLE event rate In this section we estimate the EAGLE event rate. Instead of the absolute event rate, we estimate the relative event rate normalized with the rate for normal microlensing events. When a lens object passes the line of sight with a normalized impact parameter $`u=d/R_E`$, the source star is amplified according to equation (1), where $`d`$ is the distance from the lens to the source in the lens plane. A microlensing event is defined as a phenomenon where the source star is amplified more than the given threshold amplification $`A_{th}`$. This occurs whenever a lens enters the microlensing “tube” whose radius is $`u_{th}R_E`$, where $`u_{th}`$ is the detection threshold impact parameter corresponding to $`A(u=u_{th})=A_{th}`$ (Griest 1991). This means that all detected microlensing events have to satisfy $`du_{th}R_E`$. The probability that a lens passes this region is proportional to the length of this radius $`u_{th}R_E`$. For a normal microlensing event $`u_{th}`$ is set to be unity corresponding to $`A_{th}=1.34`$, and for an EAGLE event we use $`u_T`$ defined in section 2 as $`u_{th}`$. Hence the event rate is proportional to $`u_{th}R_E=R_E`$ for a normal microlensing event and $`u_{th}R_E=u_TR_E`$ for an EAGLE event, i.e., these event rates are proportional to the radius of each appropriate circle in figure Constraining the Location of Microlensing Objects by using the Finite Source Effect in EAGLE events and Einstein radius $`R_E`$. The event rate is proportional to the density distribution of the lens $`n_s(M,D_d)`$ (Griest 1991). The event rate is also proportional to the number of source stars which is expressed as the luminosity function of source stars $`\varphi _L(V)`$ in the LMC. Moreover, the event rate should be averaged over the probability distribution of the mass of the lenses using the mass function $`\varphi (M)`$. Thus, the differential form of the event rate averaged over the source stars for normal microlensing events is written as $$d\mathrm{\Gamma }_NR_E\varphi _L(V)\varphi (M)n_s(M,D_d)dD_ddD_sdVdM.$$ (17) Note that the integration with respect to $`V`$ means averaging over the source stars. Similarly, the differential form of the rate for EAGLE events can be expressed as $$d\mathrm{\Gamma }_Eu_TR_E\varphi _L(V)\varphi (M)n_s(M,D_d)dD_ddD_sdVdM.$$ (18) Note that the threshold impact parameter $`u_T`$ appears in case of EAGLE events, since the event rate for EAGLE events is proportional to $`u_TR_E`$ instead of $`R_E`$. To calculate the normal and the EAGLE event rates, we must integrate them over the possible ranges of the parameters. The normal event rate $`\mathrm{\Gamma }_N`$ can be estimated as $$\mathrm{\Gamma }_N=C_V_M_{D_s}_{D_d}𝑑\mathrm{\Gamma }_N,$$ (19) where the integration is performed over the following ranges. $`D_{d,min}`$ $`D_d`$ $`D_s,`$ (20) $`D_{s,min}`$ $`D_s`$ $`D_{s,max},`$ (21) $`M_l`$ $`M`$ $`M_u,`$ (22) $`V_l`$ $`V`$ $`V_{obs}.`$ (23) Here $`M_l`$ and $`M_u`$ denote the lower and the upper mass limit for the mass function, and $`V_l`$ denotes the luminous end of the luminosity function $`\mathrm{\Phi }_L(V)`$. Similarly, the EAGLE event rate $`\mathrm{\Gamma }_E`$ can be estimated as $$\mathrm{\Gamma }_E=C_V_M_{D_s}_{D_d}𝑑\mathrm{\Gamma }_E.$$ (24) Note that the integral ranges for $`D_d`$, $`D_s`$ and $`M`$ are the same as for $`\mathrm{\Gamma }_N`$, but $`V`$ is in the range: $$V_{obs}VV_u.$$ (25) Note that $`V_u`$ denotes the faint end of the luminosity function. Since the constant C is equal in both equations, one can easily obtain the ratio $`\mathrm{\Gamma }_E/\mathrm{\Gamma }_N`$ by integrating above equations numerically. As for the integral ranges for the luminosity in the LMC stars, we assume that $`V_l=15`$ and $`V_u=31`$. We note that the results are not affected strongly by stars outside this magnitude range. We also assume $`V_{obs}=21`$, which is a typical observational limit for current microlensing programs. We consider two different values for the threshold magnitude; $`V_{th}`$ = 19 and 20 . We calculate these integrals in the case where the lenses are stars in the LMC itself (self-lensing) and in the case where the lenses are MACHOs in the halo. For these two cases, there are the following differences which arise in the integral ranges of ($`M`$, $`D_s`$ and $`D_d`$) and the functions of ($`n_s(M,D_d)`$ and $`\varphi (M)`$). * In the case where the lenses are known stars in the LMC itself, for simplicity, we assume a constant spatial density distribution of the lens $`n_s(M,D_d)`$ with depth $`d_{max}=D_{s,max}D_{s,min}`$. For the depth of the LMC $`d_{max}`$, one finds that the bar can be as thin as the disk itself (Binney & Tremaine (1987)), and the thickness of the LMC disk is about 300pc (Westerlund (1991)). Since there is not any more accurate information, we adopt two values, 300pc and 1kpc. We also assume that the distance to the LMC is 50 kpc, and also assume $`M_l=0.1M_{}`$, $`M_u=50M_{}`$ for the lens mass distribution expressed by equations (9) and (10) in the case of self-lensing. * In the case where the lenses are MACHOs in the halo, we fix $`D_s=50`$ kpc and integrate with $`D_d`$ from $`D_{s,min}=0`$ to $`D_s=50`$ kpc. We also adopt the standard halo model for the spatial distribution of the lens $`n_s(M,D_d)`$ as follows (Griest 1991): $$n_s(M,D_d)=\frac{\rho }{M}=\frac{\rho _0}{M}\frac{a^2+r_0^2}{a^2+r^2}$$ (26) where $`\rho _0=0.0079`$ is the local mass density, $`r`$ is the Galactocentric radius, $`r_0=8.5`$ kpc is the Galactocentric distance of the Sun, and $`a=5`$ kpc is the core radius. We also assume a delta function mass distribution for MACHOs, and we take two values of $`M`$; $`M=0.1M_{}`$ and $`1.0M_{}`$. The results for $`V_{\mathrm{th}}=`$ 19 and 20 are listed in Table 1, and the probability distribution of EAGLE events as a function of the source magnitude $`V`$ are shown in figure Constraining the Location of Microlensing Objects by using the Finite Source Effect in EAGLE events . From these results we conclude the following: * From the results in table 1 we can expect the EAGLE event rate is of the same order ($`7`$ events/yr) as the rate for normal events when $`\alpha =2.35`$, and much larger ($`70170`$ events/yr) when $`\alpha =5.0`$. Here we assume the normal event rate $`\mathrm{\Gamma }_N=4`$ events/yr (Alcock et al. (1997)). Thus, the event rate for EAGLEs is high enough to allow the events to be detected. * The event rate for EAGLEs for $`\alpha =5.0`$ is $`825`$ times higher than that for $`\alpha =2.35`$. This strong dependence of the EAGLE event rate on $`\alpha `$ indicates that we can obtain some information on the power index $`\alpha `$ of the mass function of stars in the LMC from the EAGLE event rate. * When $`\alpha =5.0`$, the EAGLE event rate for halo MACHOs is 3 times higher than that for the self-lensing case. So if $`\alpha =5.0`$, it is possible to determine whether the lens masses are MACHOs in the halo or known stars in the LMC itself from the event rate of EAGLEs. * In figure Constraining the Location of Microlensing Objects by using the Finite Source Effect in EAGLE events , we can see the EAGLE event probability decrease for large $`V`$, though the number of stars increases with increasing $`V`$. The reasons are that $`u_T`$ is small for large $`V`$ and the source can not be amplified enough to be observed due to the finite-source effect. Also, if $`\alpha =5.0`$, the distribution of the source star magnitudes may be useful for determining whether the lens masses are MACHOs in the halo or are stars in the LMC, because this distribution is significantly different in each case (see figure Constraining the Location of Microlensing Objects by using the Finite Source Effect in EAGLE events ). ## 4 Fraction of finite-source transit events In this section, we investigate the fraction of transit EAGLE events in all EAGLE events. The finite-source effect appears when $`z`$ becomes smaller than about 2 (Gould (1994)), where $`zu/u_{}`$ is the impact parameter (in the source plane) scaled by the radius of the source star $`R_{}`$ (see section 2) and $`u_{}`$ is defined in equation (5) as the source star radius normalized by the Einstein radius in the source plane. However, when $`1<z<2`$, the effect is not strong and it is quite difficult to determine whether the event is a transit event or a non-transit event (Peng (1997)). Therefore, here we define a transit event as an event with $`z<1`$. The event rate of EAGLEs is proportional to $`u_TR_E`$ (see section 3), and similarly, the probability that the lens transits the surface of the source star is proportional to $`u_{}R_E`$ . In short, these rates are proportional to the radius of the corresponding circle in figure Constraining the Location of Microlensing Objects by using the Finite Source Effect in EAGLE events . Thus, the ratio of the transit EAGLE event rate to the total EAGLE event rate is proportional to $`u_{}/u_T`$. In order to estimate the mean value of the fraction of transit EAGLE events, one must integrate this ratio over all possible parameter ranges as above, in section 3. The mean value of this fraction can be estimated by evaluating the following: $$\frac{u_{}}{u_T}=\frac{1}{\mathrm{\Gamma }_E}C_M_V_{D_s}_{D_d}\frac{u_{}}{u_T}𝑑\mathrm{\Gamma }_E$$ (27) The results for $`V_{\mathrm{th}}=19`$ and $`20`$ are shown in Table 2. From these results we conclude the following: * The fraction of transit events in all EAGLE events is much higher than that for normal events ($`0.1\%`$) (Gould 1994). So we can expect a few transit EAGLE events per year using a 1-m class telescope depending on the lens location. If we obtain an accurate light curve and a good estimate of the radius of the source star by follow-up observations, we can get the proper motion of the lens for each event (Gould (1994)). Since the lens proper motion is useful for constraining the lensing location (and so the lens population), we can determine the location of the lens for each event. * Moreover, the fraction of transit events strongly depends on whether the lens is in the halo or in the LMC. The transit fraction for the self-lensing case for $`\alpha =2.35`$ is $`415`$ times larger than that for halo MACHO lensing case (for $`\alpha =5.0`$, $`25`$ times larger). We can constrain the power index $`\alpha `$ using the event rate (see Table 1) and the distribution of the source star magnitude (see figure Constraining the Location of Microlensing Objects by using the Finite Source Effect in EAGLE events ). Hence, if we can measure the fraction of transit events in all EAGLE events, we will be able to constrain the lens location (and so the lens population) statistically. So the EAGLE event search will be of great importance in investigating the location and the nature of the lensing objects. We discuss some quantitative examples about these in section 6. ## 5 Event duration In this section, we estimate the duration, $`t_E`$ of an EAGLE event. We define the duration of an EAGLE event $`t_E`$ as the time when the source star becomes brighter than the observation apparent magnitude threshold $`V_{obs}`$. The critical amplification $`A_{obs}`$ is written as $`A_{obs}=10^{0.4(VVobs)}`$. If $`A_{obs}>2.5`$, we use equation (5) to derive $`u_{obs}`$. On the other hand, when $`A_{obs}2.5`$, we use the standard equation (1), because equation (5) is invalid and the finite-source effect is negligible for this case. The mean value of $`t_E`$ for the event with source magnitude $`V`$ can be written as $`t_E(V)`$ $`=`$ $`2\widehat{t}{\displaystyle \frac{1}{u_T}}{\displaystyle _0^{u_T}}\sqrt{u_{obs}^2u_{min}^2}𝑑u_{min},\widehat{t}={\displaystyle \frac{R_E}{v_t}}`$ (28) $``$ $`1.94\widehat{t}u_{obs}(V),(V_{\mathrm{th}}=20)`$ (29) $``$ $`1.98\widehat{t}u_{obs}(V),(V_{\mathrm{th}}=19)`$ (30) where $`\widehat{t}`$ is the true duration of the microlensing event and $`v_t`$ is the transverse velocity of the lens object relative to the observer-source line of sight. We assume $`v_t=220`$ km/s for the MACHOs in the halo and $`v_t=30`$ km/s for the stars in the LMC. The mean value of $`t_E(V)`$ averaged over all possible values of distance, lens mass and source magnitude is: $$t_E=\frac{1}{\mathrm{\Gamma }_E}C_M_V_{D_s}_{D_d}t_E(V)𝑑\mathrm{\Gamma }_E$$ (31) The results for $`V_{\mathrm{th}}=19`$ and $`20`$ are shown in Table 3, and the probability distribution of $`t_E`$ for EAGLE events is shown in figure Constraining the Location of Microlensing Objects by using the Finite Source Effect in EAGLE events . In figure Constraining the Location of Microlensing Objects by using the Finite Source Effect in EAGLE events , the distributions are shifted to shorter durations for $`\alpha =5.0`$ with respect to those for $`\alpha =2.35`$ because the number of dim sources for EAGLEs increases. We can also see a cut-off in the probability curves at smaller $`t_E`$ for self-lensing. This is because the event rate of EAGLEs whose impact parameter is smaller than some threshold value decreases according to the finite-source effect. From figure Constraining the Location of Microlensing Objects by using the Finite Source Effect in EAGLE events , we can see that the duration of the EAGLE is usually short (1 day $``$ 40 days), especially when $`\alpha =5.0`$. Hence, the nightly observational mode currently undertaken is not adequate for detecting EAGLE events, but an hourly observational mode is more suitable for detecting EAGLE events. If observations are made several times per night, there will be sufficient time resolution to detect the EAGLE events and to issue an alert for an occurring EAGLE event. ## 6 Discussion and Conclusion We have seen that the EAGLE event rate is as high as that for normal events. Since the duration of an EAGLE event is usually short (1 day $``$ 40 days), the nightly observation program currently undertaken by most groups is not adequate. Hourly observation is necessary for finding EAGLE events. Moreover, monitoring with a 1-m class telescope is sufficient to search for EAGLEs and to make follow up observations near the peak, but not enough to follow up the wing of the light curve. This is because the source is visible only at the peak of the light curve but invisible during most of the remainder of the event. To get overall structure of the light curve, one has to observe the event with a larger telescope for a longer period (but we do not need to observe so frequently). Current microlensing experiments are being carried out toward very dense stellar fields. So, the determination of the absolute flux of the source has been an issue in microlensing because of source star blending. In the case of EAGLE events, the source stars are fainter than the detection limit and thus the blending of brighter stars is significant, making the problem worse. Attempts have been made to correct the blending effects for individual events by introducing an additional parameter, the blended flux, but this method suffers from very large uncertainties in the derived lensing parameters due to degeneracies among the parameters (Alard (1997); Han (1997)). One solution to correct for blending problem is to use the Hubble Space Telescope (HST). The high resolving power of the HST and the color information from ground-based observations enable us to uniquely separate the lensed source star from blended stars (Han (1997)). The image subtraction method (Alard (1998), 2000; Alcock et al. 1999b ; I.A. Bond 1999, private communication) should be used to locate the centroid of the lensed flux in the ground-based event image accurately to find the lensed star in the HST image. Other high resolution telescopes like the VLT are also good candidates for performing follow-up observations. Using a large telescope is also required for spectroscopic observations of the source star, from which one can estimate the source radius. We can do this even after the event is finished. Using the image subtraction method has another advantage in microlensing experiments. The color change induced by the limb-darkened extended source helps the measurement of the proper motion of the lens and increases the possible number of proper motion measurement events (Witt (1995); Gould (1996)). Han (Han (2000)) found the color change measurement by using the image subtraction method enables one to obtain the same information about the lens and the source star, but with significantly reduced uncertainties due to the absence of the blending effect. So measuring the color change with differential photometry helps to derive the proper motion of the lens along with the other microlensing parameters, and at least, inform us whether a lens transited the surface of the source or not without requiring the absolute baseline flux. This information is useful in understanding the location of the lens object which we discuss below. In order to perform follow-up monitoring and spectroscopic observations, one has to detect an EAGLE event in real-time and issue an alert. For the real-time detection of EAGLE events, the image subtraction method is more suitable than the DoPHOT analysis or the pixel analysis, since the image subtraction method can detect the luminosity variation at any position of the fields, even where no star was identified previously. Thus, the most reasonable and practical observation strategy is to observe hourly (with a 1-m telescope) and to perform the real-time analysis with the subtraction method. An alert can then be issued, to observers around the world enabling frequent observations around the event peak. After that, high resolution observations with larger telescope should be carried out to determine the baseline flux of the source. As seen in section 4, a large fraction of EAGLE events is likely to be transit EAGLE events. If we can issue an alert for such an EAGLE event, observations with the global network will allow us to obtain a precise light curve near the maximum amplification. Because the transit time is $`13`$ hours for self-lensing and $`2`$ hours for halo MACHOs, it is possible to measure precisely the light curve in the transit region. Since the overall structure of the light curve can be determined by follow-up observations with larger telescopes, one can obtain the source radius scaled by Einstein radius $`u_{}`$ from the full light curve and from the color change measurements due to a limb-darkened extended source (Han (2000)). If one can obtain $`u_{}`$ and also the radius of the source star $`R_{}`$ from the color of the source star, one can obtain the Einstein radius $`r_E`$ projected in source plane. Since the event duration $`\widehat{t}`$ can be also determined from the full light curve, one can measure the proper motion of the lens object for a transit EAGLE event (Gould (1994)). The proper motion of the lens object is of great use in determining the lens location. Since the proper motion of objects in the halo is $`3040`$ times larger than that of stars in the LMC, one can know whether the lensing object exists in the halo or not from the proper motion of transit EAGLE events. Even if we cannot derive an accurate proper motion of lens, the fraction of transit events can be used to constrain the lens location, because that ratio for the self-lensing case is $`215`$ times larger than that for the halo MACHOs case. To demonstrate how we can constrain the nature of microlensing objects, we show in figure Constraining the Location of Microlensing Objects by using the Finite Source Effect in EAGLE events the probability distribution of the transit events for both the halo lensing case and the LMC self-lensing case. For the typical parameters $`\alpha =2.35`$, $`V_{th}=20`$ and a detection efficiency of 50%, one can expect to find $`21`$ EAGLE events after a 3 year observation period of 11 square degrees of the LMC central region (as the MACHO collaboration does). In these 21 EAGLE events, we can expect 1.1 (0.4) and 6.1 (4.6) transit events for halo MACHOs whose mass is $`0.1(1.0)M_{}`$ and for self-lensing with $`d_{max}=300`$ (1k) pc, respectively. So if we do not detect any transit events in 3 years, we can reject the possibility of self-lensing at more than $`99\%`$ confidence level according to Poisson statistics. In table 4, we show some other results on how strongly we can discriminate between the two possible locations of the lensing objects based on Poisson statistics. It is possible that we may be able to constrain strongly the lens location based on the 3-year statistics of transit EAGLE events. In conclusion, the study of EAGLE events is very effective for constraining the nature of lens objects. If an EAGLE search is made for a few years, we can form strong constraints on lens objects in microlensing events. We are grateful to Prof. Yasushi Muraki for his supervision, and also to Yukitoshi Kan-ya, Yutaka Matsubara and Philip Yock for their helpful comments. We would like to acknowledge the careful reading our manuscript with Nicholas Rattenbury.
warning/0004/math0004177.html	ar5iv	text	# Shock capturing by anisotropic diffusion oscillation reduction ## I Introduction Shock wave is a common phenomenon in nature, such as in aerodynamics and hydrodynamics. Mathematically, nonlinear hyperbolic conservation equations provide a good description to shock waves. The construction of numerical schemes that are capable of shock capturing for hyperbolic and inviscid hydrodynamic equations is a major objective of computational methodology. However, it is noted that the concept of discontinuity does not apply in digital computations. Therefore, a shock in computational sense may refer to systematic, rapid variation of function values over a few grid points. The difficulty is that hyperbolic equation may have weak solutions that are discontinuous at the so-called shock front. Such a discontinuity will cause Gibbs’ oscillations in a high (spatial) order numerical scheme. The numerically induced oscillations are usually amplified in (time) iterations. A variety of numerical schemes have been proposed for shock capturing. As early as 50 years ago, a solution to this problem was constructed by von Neumann and Richtmyer. The essence of their approach is to introduce small artificial viscosity so that a smooth solution can alway be attained in a finite difference approximation. A variety of modifications to von Neumann and Richtmyer’s method are made in the last 5 decades to address problems of possible failure in spatial scaling and errors due to additional momentum flux and production, as well as unbalanced heat flux and its over production in the simulation of hydrodynamic conservation laws. A different approach was proposed by Godunov to construct a full solution by using low order piecewise discontinuous approximations. Such a piecewise solution is a good approximation at the smooth regions, and is capable of representing the shock front over a small region of grid. The knowledge of wave propagation and wave interaction is built in the numerical scheme in the form of a Riemann solver. The Godunov’s approach has been extended to higher-order schemes. In doing so, a higher-order approximation is constructed in the smooth region, while near the shock the solution is still of first order accuracy. The Godunov method is very stable and thus, easy to design and use. However the major disadvantage of this method is the complexity introduced into a numerical scheme through a Riemann solver. Such a drawback reduces its computational efficiency. Another general approach is the hybrid scheme, which utilizes a high-order scheme for a smooth region while using a low-order scheme near a discontinuity. A linear combination of these two types of schemes is then used at each interface using weight factors which may be nonlinear functions of the local flow field. DeBar constructed a linear hybridization of first-order and second-order difference schemes as early as 1968. Harten and Zwas devised another early linear hybridization scheme. Boris and Book reported a blending algorithm which yields sharp discontinuities without oscillations. A total variation diminishing (TVD) scheme was proposed by Harten to control the spurious oscillations in the numerical solution by using the total variation as a measure. The TVD scheme typically degenerates to first-order accuracy at locations with smooth extrema and was later generalized to an essentially non-oscillatory (ENO) scheme. The major idea of the ENO scheme is to suppress spurious oscillations near the shock or discontinuities, while maintaining a higher-order accuracy at smooth regions. This line of thinking was further polished recently in a weighted essentially non-oscillatory (WENO) scheme. The WENO approach takes a linear combination of a number of high-order schemes of both central difference and up-wind type. The central difference type schemes have a larger weight factor at the smooth region while the up-wind schemes play a major role at the shock or discontinuity. In general, these approaches are quite expensive since checks are performed before making a decision at each grid point. In terms of accuracy, all existing methods are at best of first order near the shock or discontinuity. Perhaps modified artificial viscosity methods are the most popular approaches in practical computations. Unfortunately, the artificial viscosity smears shocks over three or more grid zones, which can lead to serious errors in the physical interpretation of the numerical results. Special care is required to ensure that the smeared numerical shock is consistent with the true thickness of the shock in a practical problem under study. This difficulty has led to enormous and continuous effort at developing efficient and robust approaches. One approach is to locally refine the computational grid. An alternative approach is the use of an adaptive mesh. Both methods are aimed at matching between the physical shock and the numerically smoothed one with respect to the spatial extension. However, a major drawback is their restriction for extremely small time step sizes as required by the Courant constraint. In many cases, they have to be formulated in an implicit scheme, which imposes an extra complexity in practical implementation and an extra requirement in computer memory. Another problem is the unbalanced momentum flux and production, and additional heat flux and heat production in hydrodynamic equations. The addition of a viscous term to the inviscid hydrodynamic equations or the hyperbolic equation has a physical justification, that the true physical flow has no discontinuities. Therefore, mathematical model can be modified so as to reflect the true physics. In their original work, von Neumann and Richtmyer have introduced an artificial viscosity $`q_{NR}`$ of the form $$q_{NR}\left(𝐯\right)^2$$ (1) for the equation of motion. Their term is of second order in gradient of the velocity field $`𝐯`$ (The velocity gradient is a second rank tensor and has the divergence of velocity as its component) and will not be very sensitive to small gradients. Landshoff proposed an additional term that vanishes less rapidly for small gradients $$q_L𝐯.$$ (2) A generic form that contains both $`q_{NR}`$ and $`q_L`$ (i.e., $`\alpha 𝐯+\beta \left(𝐯\right)^2`$) has been carefully studied by many researchers. In particular, Caramana et al discussed a number of intuitive criteria for this form and the parameterization of artificial viscosity. Special considerations are given to dissipativity, Galilean invariance, self-similar motion invariance, wave front invariance and viscous force continuity. This generic form, as it was proposed for hydrodynamic conservation laws, has its advantages over other approaches for certain fluid mechanical computations. However, it is very unstable for certain hyperbolic equations due to the Courant-Friedrich-Lewy (CFL) stability condition. For example, it does not work well in resolving a sharp shock front for inviscid Burgers’ equation with a smooth initial condition. We believe that this failure is due to the high-order in the gradient of the Neumann and Richtmyer form, which is too sensitive to large gradients at a sharp shock front. Despite the fact that a number of functional forms of the artificial viscosity have been proposed, the procedure is still ad hoc. For a given form, parameter selection is quite tricky and often varies from problem to problem. The existence of so many different approaches for shock capturing indicates both the importance and the difficulty of the problem. The objective of the present paper is to introduce an anisotropic diffusion oscillation reduction (ADOR) approach for shock wave computations. The method of anisotropic diffusion in association with a partial differential equation was proposed by Perona and Malik for digital image processing in 1990. Since then, much research has been stimulated in image processing and applied mathematics communities. The method uses the heat equation, which has a solution of the Gaussian type, as a low pass filter to eliminate noise in an image, while it detects and preserves the edges of the image. It has been shown that the Perona-Malik equation provides a computational approach to image segmentation, noise removal, edge detection, and image enhancement. In a recent work, a generalized Perona-Malik equation was proposed for image restoration and edge enhancement. In fact, digital edge detection has much in common with numerical shock capturing. In this work, we introduce a set of generalized anisotropic diffusion operators and edge enhancing functionals for shock wave computations. This paper is organized as follows: Section II is devoted to kinetic theory analysis of the equation of change and artificial viscosity. The physical origin and mathematical properties of artificial viscosity are discussed from the kinetic theory point of view. A set of hydrodynamic equations describing fluid flow consisting of microscopic particles is derived from the quantum Boltzmann equation, i.e., the Waldmann-Snider equation. The latter can be regarded as a consequence of a reduction of the von Neumann equation, or the quantum Liouville equation to the level of a single particle density operator. The mathematical form of the pressure tensor is analyzed by using the group theory consideration. A set of generalized artificial viscosities, artificial heat and artificial virial correction are proposed as the results of the kinetic theory analysis. Generalized Perona-Malik equation is discussed in Section III for shock wave computations. This approach follows a very different line of thinking as it was originally proposed for image analysis and processing. An edge-controlled image enhancing functional is introduced for shock representation. Anisotropic diffusion and diffusion super operators are introduced for fast and effective shock capturing. Numerical experiments are presented in Section IV. A few standard test problems, including Burgers’ equation and inviscid Burgers’ equation in one- and two-dimensions, the incompressible Euler and Navier-Stokes equations are used for exploring usefulness, testing efficiency and illustrating the validity of the ADOR approach. A recently developed discrete singular convolution (DSC) algorithm is utilized for the numerical integration. This paper ends with a conclusion. ## II Kinetic theory analysis of artificial viscosity Artificial viscosity was original proposed for shock capturing in association with hydrodynamic equations in the fluid mechanics. The choice of its form and parameter should be consistent with and motivated by the physical origin of hydrodynamic equations. Since quantum theory is the foundation for the modern fluid mechanics, a microscopic analysis is presented for appropriate understanding from the kinetic theory. The latter, such as the Boltzmann equation, serves as a theoretical basis for hydrodynamic conservation laws. ### A Microscopic Analysis Consider the flow of a quantum gas system consisting of total $`N`$ particles in a volume $`S`$. Its behavior is governed by the Schrödinger equation $$i\mathrm{}\frac{\mathrm{\Psi }}{t}=H^{(N)}\mathrm{\Psi },$$ (3) where $`H^{(N)}`$ is the self-adjoint Hamiltonian of the system and $`\mathrm{\Psi }`$ is a vector in the Hilbert space associated with the system. For the description of physical observable, we adopt the density operator $`\rho ^{(N)}`$ whose time evolution is governed by the quantum Liouville equation $$i\frac{\rho ^{(N)}}{t}=^{(N)}\rho ^{(N)}=\frac{1}{\mathrm{}}[H^{(N)},\rho ^{(N)}]_{}=\frac{1}{\mathrm{}}(H^{(N)}\rho ^{(N)}\rho ^{(N)}H^{(N)}).$$ (4) A physical observable $`O^{(N)}`$ is a Hermitian operator of the Hilbert space and has the expectation value given by $$<O^{(N)}>=\mathrm{Tr}_{1,\mathrm{},N}O^{(N)}\rho ^{(N)},$$ (5) where $`\mathrm{Tr}_{1,\mathrm{},N}`$ is the trace over all the states of the $`N`$ particle system and in particular, $$<1^{(N)}>=\mathrm{Tr}_{1,\mathrm{},N}1^{(N)}\rho ^{(N)}=1$$ (6) gives the normalization of the total probability for finding all of the $`N`$ particles in the volume $`S`$. Here $`1^{(N)}`$ is the identity operator of the system. A standard form for the Hamiltonian is given by $$H^{(N)}=K^{(N)}+V^{(N)}=\underset{i}{}K_i+\underset{i<j}{}V_{ij},$$ (7) where $`K_i`$ is the kinetic energy (including internal state energy) operator of particle $`i`$ and $`V_{ij}`$ is the potential operator of particles $`i`$ and $`j`$. The physical behavior of the $`N`$-particle system is far too complicated to compute by any means as a macroscopic gas flow may consist of $`10^{23}`$ particles or more. Fortunately, for ordinary gases, it is sufficient to consider the state of a typical particle, say particle 1 $$\rho _1^{(1)}=N\mathrm{Tr}_{2,\mathrm{},N}\rho ^{(N)}.$$ (8) In general, the state of $`n`$ particles is defined $$\rho _{1,\mathrm{},n}^{(n)}=N(N1)\mathrm{}(Nn+1)\mathrm{Tr}_{n+1,\mathrm{},N}\rho ^{(N)}.$$ (9) The time evolution of particle 1 is governed by $$i\frac{\rho _1^{(1)}}{t}=^{(1)}\rho _1^{(1)}+\mathrm{Tr}_2,𝒱_{12}^{(2)}\rho _{12}^{(2)}=\frac{1}{\mathrm{}}[H^{(1)},\rho _1^{(1)}]_{}+\frac{1}{\mathrm{}}[V_{12},\rho _{12}^{(2)}]_{}$$ (10) where $`H^{(1)}=K_1`$. This is the first member of the BBGKY hierarchy in the quantum form. The general form of the BBGKY hierarchy is given by $$i\frac{\rho _{1,\mathrm{},n}^{(n)}}{t}=^{(n)}\rho _{1,\mathrm{},n}^{(n)}+\mathrm{Tr}_2,𝒱_{1,\mathrm{},n+1}^{(n+1)}\rho _{1,\mathrm{},n+1}^{(n+1)},$$ (11) where $`𝒱_{1,\mathrm{},n+1}^{(n+1)}`$ is the potential superoperator between particles $`1,\mathrm{},n`$ and particle $`n+1`$. This set of equations is formal and exact, since no approximation has been made. However, to determine the time evolution of $`\rho _1^{(1)}`$, it is necessary to know the behavior of $`\rho _{12}^{(2)}`$, which is, in turn, determined by the second order BBGKY equation and the latter involves three-particle density operator $`\rho _{123}^{(3)}`$. ### B Mesoscopic analysis Kinetic theory attempts to explain macroscopic properties in terms of microscopic properties of the atoms and/or molecules based on the classical or quantum mechanics. Typical macroscopic observations deal with $`N(10^{23})`$ particles over a volume much larger than the size of the individual molecules, and over a time period much longer compared to the time scale of the individual molecular dynamics. An ab-initio description of the time dependence of such macroscopic observations from the quantum Liouville equation is practically impossible not only because of the conceptual difficulty with the meaning of measurement, but also due to the large number of degrees of freedom involved. The kinetic theory approach simplifies the $`N`$-particle problem dramatically by looking at the behavior of one typical particle under the influence of all other particles. The most successful kinetic theory has been based on the Boltzmann equation. Specifically this has been related to the hydrodynamic equations and has given rise to molecular expressions for the transport coefficients. Various attempts have been made to understand the Boltzmann equation from the $`N`$-body Liouville equation, or utilizing the BBGKY hierarchy described in the last subsection. A first principle “derivation” of the Boltzmann equation from the Liouville equation has been achieved by Bogoliubov, and independently by Green. The Wang Chang-Uhlenbeck-de Boer equation was the first attempt to generalize the quantum Boltzmann equation, and to account for the internal degrees of freedom of a polyatomic gas. It was not generally realized until 1957 that polyatomic gases can only be treated properly by a quantum mechanical method due to the fact that the internal energy levels of such molecules are, in general, degenerate. Waldmann and, independently, Snider introduced a quantum kinetic equation, the Waldmann-Snider equation, $$i\frac{\rho _1^{(1)}}{t}=_1^{(1)}\rho _1^{(1)}+\mathrm{Tr}_2𝒱_{12}^{(2)}\mathrm{\Omega }_{L_{12}}\rho _1^{(1)}\rho _2^{(1)}$$ (12) where $`\mathrm{\Omega }_{L_{12}}`$ is the Møller superoperator for quantum mechanical scattering $$\mathrm{\Omega }_{L_{12}}=\underset{t\mathrm{}}{lim}e^{i_{12}^{(2)}t}e^{i𝒦_{12}^{(2)}t}.$$ (13) The Waldmann-Snider equation is still the basis for the kinetic theory of quantum gas flow. Note that a dramatically simplified version of the Boltzmann equation provides the (lattice) Boltzmann gas (LBG) approach for the computational fluid dynamics in the recent years. ### C Macroscopic analysis The most important properties of a fluid flow are physical observables, or physical measurements. For conservative observables, the equations of change of physical observables give rise to conservation laws. The equations of change are important not only because they govern the time dependencies of the macroscopic quantities but also because they are needed for solving kinetic equations using the Chapman-Enskog method. The equations of change for physical observables are derived in this subsection. A single-particle observable for particle 1 is denoted by $`\psi _1`$. The expectation value of the physical observable $`\psi _1`$ is determined by the singlet density operator $`\rho _1^{(1)}`$ according to $$\psi _1=\mathrm{Tr}_1\psi _1\rho _1^{(1)}.$$ (14) The most important physical observables are the number density $`n_1`$, stream velocity $`𝐯_1`$ and kinetic energy $`\epsilon _1^K`$. Firstly, the physical observable for the number density $`n_1`$ is the Dirac delta distribution $`\psi _1=\delta _1\delta (𝐫𝐫_1)`$ so that $$n_1=\delta _1.$$ (15) Secondly, the stream velocity $`𝐯_1`$ is given by $$𝐯_1\frac{1}{n_1}\mathrm{Tr}_1\frac{1}{2m_1}(𝐩_1\delta _1+\delta _1𝐩_1)\rho _1,$$ (16) where $`m_1`$ is the mass of particle 1. Finally, the kinetic energy $`\epsilon _1^K`$ per particle is given by $$\epsilon _1^K\frac{1}{n_1}\mathrm{Tr}_1\left(\frac{1}{8m_1}(𝐩_1^2\delta _1+2𝐩_1\delta _1𝐩_1+\delta _1𝐩_1^2)\frac{1}{2}m_1𝐯_1^2\delta _1\right)\rho _1.$$ (17) These expressions are appropriately symmetrized to account for the fact that position and momentum do not commute, while a physical observable is a Hermitian operator. Although this operator approach is used here for for the number density $`n_1`$, stream velocity $`𝐯_1`$ and kinetic energy $`\epsilon _1^K`$, phase space expressions can be easily obtained for these fluid dynamic quantities by using the standard Wigner representation. The equations of change are of course independent of the detailed method used for the expression of various quantities. Equations of change for a particle observable is obtained from the quantum Boltzmann equation (12) according to $`{\displaystyle \frac{\psi _1_1}{t}}{\displaystyle \frac{\psi _1}{t}}_1={\displaystyle \frac{1}{i\mathrm{}}}[\psi _1,H^{(1)}]_{}_1+{\displaystyle \frac{1}{i}}\mathrm{Tr}_{1,2}𝒱_{12}^{(2)}\mathrm{\Omega }_{L_{12}}\rho _1^{(1)}\rho _2^{(1)}.`$ (18) In case that a physical observable $`\psi _1`$ is time independent, the second term on the left hand side of Eq. (18) vanishes. Conservation laws are the equation of continuity $$\frac{n_1}{t}=(n_1𝐯_1);$$ (19) the equation of motion $$n_1m_1\frac{𝐯_1}{t}+n_1m_1𝐯_1𝐯_1=\text{P}_1,$$ (20) and the kinetic energy equation $$\frac{n_1\epsilon _1^K}{t}=(n_1𝐯_1\epsilon _1^K+𝐪_1^K+𝐪_{\mathrm{coll}})\text{P}_1^t:𝐯_1+\sigma ^K,$$ (21) where the superscript $`t`$ denotes the transpose and the kinetic heat flux is $$𝐪_1^K=\mathrm{Tr}_1\left\{\left(\frac{𝐩_1}{m_1}𝐯_1\right)\frac{(𝐩_1m_1𝐯_1)^2}{2m_1}\delta _1\right\}_s\rho _1,$$ (22) where $`\{\}_s`$ designates appropriately operator-symmetrized quantities. Here the collisional heat flux is $`𝐪_{\mathrm{coll}}`$ $``$ $`{\displaystyle \frac{1}{8m_1}}\mathrm{Tr}_{12}{\displaystyle _1^1}d\nu \{\delta [𝐫{\displaystyle \frac{𝐫_1+𝐫_2}{2}}{\displaystyle \frac{\nu }{2}}(𝐫_1𝐫_2)]`$ (23) $`\times `$ $`(𝐩_1+𝐩_22m_1𝐯_1){\displaystyle \frac{V_{12}}{𝐫_1}}\}_s\mathrm{\Omega }_{L_{12}}\rho ^{(1)}_1\rho ^{(1)}_2`$ (24) and the kinetic energy production results from the collisional transfer of energy $`\sigma ^K`$ $`=`$ $`{\displaystyle \frac{1}{8m_1}}\mathrm{Tr}_{12}\{(\delta _1+\delta _2)(𝐩_1𝐩_2){\displaystyle \frac{V_{12}}{𝐫_1}}`$ (26) $`+{\displaystyle \frac{V_{12}}{𝐫_1}}(𝐩_1𝐩_2)\}_s\mathrm{\Omega }_{L_{12}}\rho ^{(1)}_1\rho ^{(1)}_2.`$ The pressure tensor $`\text{P}_1^K`$ has kinetic and collisional contributions $$\text{P}_1=\text{P}_1^K+\text{P}_1^{\mathrm{coll}}$$ (27) where the kinetic pressure tensor is $$\text{P}_1^K\frac{1}{4m_1}\delta _1𝐩_1𝐩_1+𝐩_1\delta _1𝐩_1+(𝐩_1\delta _1𝐩_1)^t+𝐩_1𝐩_1\delta _1_1n_1m_1𝐯_1𝐯_1$$ (28) and the collisional pressure tensor is given by $`\text{P}_1^{\mathrm{coll}}`$ $`=`$ $`{\displaystyle \frac{1}{4}}\mathrm{Tr}_{12}(𝐫_1𝐫_2)_1V_{12}`$ (29) $`\times `$ $`{\displaystyle _1^1}\delta \left({\displaystyle \frac{𝐫_1+𝐫_2}{2}}+{\displaystyle \frac{\nu }{2}}(𝐫_1𝐫_2)𝐫\right)𝑑\nu \mathrm{\Omega }_{L_{12}}\rho _1^{(1)}\rho _2^{(1)}.`$ (30) The transpose $`\text{P}_1^t`$ of the pressure tensor $`\text{P}_1`$ enters to couple the convective energy to the internal kinetic energy per particle $`\epsilon _1^K`$, heat flux contributions arise from kinetic $`𝐪_1^K`$ and collisional $`𝐪_{\mathrm{coll}}`$ motion and finally there is a production term $`\sigma ^K`$ which involves the transfer between kinetic and potential energies. Although the equations of continuity and motion are consistent with the conservation of the number of particles and total momentum in the absence of bound pairs, the presence of the production term in the kinetic energy implies that the total kinetic energy is not conserved because of the possible conversion to potential energy due to the non-locality of the collisions (in the macroscopic sense). Thus it is necessary to look at the equation of change for the potential energy per particle $`\epsilon _1^V`$ $$n_1\epsilon _1^V=\frac{1}{4}\mathrm{Tr}_{12}(\delta _1+\delta _2)V_{12}\rho _{12}^{(2)}.$$ (31) This is obtained in a manner consistent with pair particle interactions and shown to be of the form $$\frac{n_1\epsilon _1^V}{t}=(n_1𝐯_1\epsilon _1^V+𝐪_1^V)\sigma ^K,$$ (32) where $$𝐪_1^V=\frac{1}{4m_1}\mathrm{Tr}_{12}\left[(𝐩_1m_1𝐯_1)\delta _1+\delta _1(𝐩_1m_1𝐯_1)\right]V_{12}\mathrm{\Omega }_{L_{12}}\rho _1^{(1)}\rho _2^{(1)}.$$ (33) The production of potential and kinetic energy exactly cancels and there is an added heat flux contribution $`𝐪_1^V`$ associated with the conductive potential energy flow. ### D Pressure tensor and artificial viscosity The estimation of transport coefficients from the Boltzmann equation is quite complicated and involves many aspects. The most important transport quantities for physical and engineering applications are the mass, momentum and energy. Sometimes angular momentum transport may also be important for certain physical phenomena and thus an equation of change for angular momentum may be required. However, angular momentum conservation is rarely discussed in the context of hydrodynamic conservation laws. This is because the angular momentum is not measured like other physical observables in experiments. It is noted that all hydrodynamic conservation laws have the structure of convection, conduction and production. The equation of momentum conservation leads to shock wave as the divergence of the pressure tensor vanishes. The existence of shock wave devastates the numerical simulation as noted by non Neumann and Richtmyer. They introduced artificial viscosity of the form of $`(𝐯)^2`$ to overcome the numerical difficulty. In fact, a much better choice can be selected by analyzing the form of the pressure tensor. The Pressure tensor is a tensor of second rank and has two indices. It has contributions from the kinetic transport $`\text{P}^K`$ and collision transport $`\text{P}^{\mathrm{coll}}`$ (Here, the subscript 1 is omitted for convenience). For example, the kinetic transport $`\text{P}_{xy}^K`$ is the rate of transport (kg m/s)/(m<sup>2</sup>s) in the x-direction, of the momentum with respect to the y-direction. The hydrodynamic pressure $$\text{P}_{\mathrm{eq}}=nk_BT\text{U},$$ (34) given by the equation of state, is the local equilibrium contribution to the pressure tensor. Here, U is the identity tensor of second rank and $`k_B`$ is the Boltzmann constant. In general, the gas is not at local equilibrium and the structure of the nonequilibrium part of the pressure tensor, $$𝚷=\text{P}\text{P}_{\mathrm{eq}}$$ (35) can be analyzed according to the irreducible representation of the improper three-dimensional rotational group $`O(3)`$. As such, $`𝚷`$ can be expressed as a linear combination the irreducible representations $$𝚷=\mathrm{\Pi }\text{U}+\epsilon \text{P}^a+𝚷^{(2)},$$ (36) where scalar, antisymmetric and symmetric traceless parts of the second rank tensor are given by $`\mathrm{\Pi }`$ $`=`$ $`{\displaystyle \frac{1}{3}}\text{U}:𝚷`$ (37) $`\text{P}^a`$ $`=`$ $`{\displaystyle \frac{1}{2}}\epsilon 𝚷`$ (38) $`𝚷^{(2)}`$ $`=`$ $`{\displaystyle \frac{1}{2}}[𝚷+𝚷^t]{\displaystyle \frac{1}{3}}\text{U}[\text{U}:𝚷].`$ (39) Here $`\epsilon `$ is the Levi-Civita tensor, $`\mathrm{\Pi }`$ is a scalar (one dimension), $`\text{P}^a`$ is a vector (three dimensions) and $`𝚷^{(2)}`$ is a tensor (five dimensions). In the theory of the Boltzmann equation and irreversible thermodynamics, the pressure tensor is treated as proportional to velocity gradient as in the case of a classical Newtonian flow. The velocity gradient, $`𝐯`$, is also a second rank tensor and can also be decomposed into three different components of order zero, one and two, similar to Eq. (37), $`𝐯`$ $`=`$ $`{\displaystyle \frac{1}{3}}\text{U}:𝐯`$ (40) $`\times 𝐯`$ $`=`$ $`{\displaystyle \frac{1}{2}}\epsilon 𝐯`$ (41) $`\left[𝐯\right]^{(2)}`$ $`=`$ $`{\displaystyle \frac{1}{2}}[𝐯+(𝐯)^t]{\displaystyle \frac{1}{3}}\text{U}𝐯.`$ (42) Here, $`𝐯`$ is also a nine-dimensional quantity. In principle, there are 81 possible coefficients connecting the pressure tensor and the velocity gradient. Symmetry analysis indicates that there are three phenomenological equations relating the corresponding components of $`𝚷`$ and $`𝐯`$ $`\mathrm{\Pi }`$ $`=`$ $`\eta _V^0𝐯`$ (43) $`\text{P}^a`$ $`=`$ $`\eta _r^0\times 𝐯`$ (44) $`𝚷^{(2)}`$ $`=`$ $`2\eta ^0[𝐯]^{(2)},`$ (45) where $`\eta _V^0,\eta _r^0`$ and $`\eta ^0`$ are the viscosities of bulk, rotational and shear. The pressure tensor given by Eq (43) provides an excellent approximation for a gas system near the local equilibrium, i.e. the Boltzmann distribution $$\rho _1^{(1)}=\frac{1}{\mathrm{Tr}_1e^{K_1/k_BT}}e^{K_1/k_BT}.$$ (46) For the purpose of numerical stability and shock wave simulation, we can impose artificial pressure tensor as $`\mathrm{\Pi }_{\mathrm{art}}`$ $`=`$ $`\zeta _V^0𝐯`$ (47) $`\text{P}_{\mathrm{art}}^a`$ $`=`$ $`\zeta _r^0\times 𝐯`$ (48) $`𝚷_{\mathrm{art}}^{(2)}`$ $`=`$ $`2\zeta ^0[𝐯]^{(2)},`$ (49) where $`\zeta _V^0,\zeta _r^0`$ and $`\zeta ^0`$ are artificial bulk viscosity, artificial rotational viscosity and artificial shear viscosity, respectively. These artificial viscosities modify the non equilibrium part of the pressure tensor. In fact, from the point of view of computations, the equilibrium part of the pressure tensor can also be modified, i.e. a term which is proportional to $`nk_BT`$ $$\text{P}_{\mathrm{nT}}=ϵnk_BT\text{U}.$$ (50) $`\text{P}_{\mathrm{nT}}`$ can be interpreted either as artificial heat or artificial virial correction. Hence, we have the artificial pressure tensor of the form $`\text{P}_{\mathrm{art}}`$ $`=`$ $`\text{P}_{\mathrm{nT}}+\mathrm{\Pi }_{\mathrm{art}}\text{U}+\epsilon \text{P}_{\mathrm{art}}^a+𝚷_{\mathrm{art}}^{(2)}`$ (51) $`=`$ $`ϵnk_BT\text{U}\zeta _V^0𝐯\text{U}\zeta _r^0\epsilon \times 𝐯2\zeta ^0[𝐯]^{(2)}.`$ (52) It should be noted that there is an advantage in keeping only a part of the artificial pressure tensor that corresponds to the non-equilibrium pressure tensor. The latter makes no additional contribution to conservative quantities, such as, number density, momentum and energy, as required by the Fredholm alternative for the existence of a solution for the kinetic equation in the Chapman-Enskog method. However, for a gaseous system far from the local equilibrium, higher-order terms in the velocity gradient become important. Large velocity gradient can build up from special boundary condition, such as the boundary layer phenomena, or from the lack of relaxation in a very dilute gas flow. Therefore artificial viscosities and artificial heat of Eqs. (47) can be modified to include higher-order velocity gradient contributions. Recognizing that the pressure tensor is a second rank tensor, thus all velocity gradients contribute in one of the forms of $`\mathrm{\Pi },\text{P}^a`$ and $`𝚷^{(2)}`$. In this regard, the artificial viscosity form proposed by von Neumann and Richtmyer, $`(𝐯)^2`$, contributes to $`\mathrm{\Pi }`$. Certainly, artificial viscosities of the forms of $`𝐯`$, $`\times 𝐯`$ and $`[𝐯]^{(2)}`$ can be used for numerical computations. All of the aforementioned forms are at least of the first order in gradient, which may lead a strong smoothing at the edge of a shock front and thus lead to large errors in numerical applications. Furthermore, these forms can result in instability when the shock front is very sharp, such as in inviscid Burgers’ equation. For these reasons, it is appropriate to propose a general expression for the artificial viscosity and artificial heat so as to allow all coefficients to be functions of $`𝐯`$, $`\times 𝐯`$ and $`[𝐯]^{(2)}`$ $`ϵ`$ $`=`$ $`ϵ(𝐯,\times 𝐯,[𝐯]^{(2)})`$ (53) $`\zeta _V^0`$ $`=`$ $`\zeta _V^0(𝐯,\times 𝐯,[𝐯]^{(2)})`$ (54) $`\zeta _r^0`$ $`=`$ $`\zeta _r^0(𝐯,\times 𝐯,[𝐯]^{(2)})`$ (55) $`\zeta ^0`$ $`=`$ $`\zeta ^0(𝐯,\times 𝐯,[𝐯]^{(2)})`$ (56) where $``$ denotes the magnitude and is computed as $$𝐀=\sqrt{\underset{i}{}A_i^2}$$ (57) for a vector $`𝐀`$, and $$\text{B}=\sqrt{\underset{ij}{}\text{B}_{ij}^2}$$ (58) for a tensor B. Obviously, artificial viscosity of von Neumann and Richtmyer, $`(𝐯)^2`$, becomes a special case of the present treatment. The von Neumann and Richtmyer form can be classified either as an artificial heat term, $`\text{P}_{\mathrm{nT}}`$, or as an artificial bulk viscosity term, $`\mathrm{\Pi }_{\mathrm{art}}\text{U}`$. The inclusion of an artificial rotational viscosity can be an efficient way for the treatment of fast rotational flow. It is noted that the rotational viscosity is in the original Navier-Stokes equation derived from the kinetic theory. However, for flows with natural boundary conditions, the angular momentum is also conserved and it does not couple the linear momentum except through internal spinor relaxations. However, this is no longer the case in a flow with an irregular geometry. For such a case, the internal-conversion between angular momentum and linear momentum is substantial and therefore the rotational viscosity, $`\eta _r^0`$, should be retained in both theoretical modeling and numerical simulation. A detailed discussion and numerical simulation of this aspect will be accounted elsewhere. In the next section, we analyze the shock capturing algorithm further from the point of view of image processing, particularly, image edge detection. ## III Oscillation reduction by anisotropic diffusion Although anisotropic diffusion was original proposed for image processing, there is much in common between digital image processing and computational fluid dynamics. An image function $`I(𝐫)`$ is a two-dimensional projection of certain physical quantities, such as matter, velocity, energy, electromagnetic field, etc., under appropriate illumination conditions. The edges in an image usually refer to rapid changes in some physical properties, such as geometry, illumination, and reflectance. Mathematically, a discontinuity may be involved in the function representing such physical properties. Therefore image edges are very similar to shocks in fluid dynamics. Numerical shock capturing can be formulated on the lines of iterative digital edge detection. Edge detection is a key issue in image processing, computer vision, and pattern recognition. A variety of algorithms, such as the Sobel operator, the Prewitt operator, the Canny operator and the DSC algorithm are proposed for image edge detection and representation. Anisotropic diffusion is a promising new mathematical algorithm for image edge detection and image processing. This basic idea can be adopted and modified for numerical shock capturing in association with the hyperbolic conservation laws. The Perona-Malik algorithm is reviewed before the method of anisotropic diffusion oscillation reduction (ADOR) is discussed. ### A The Perona-Malik equation The basic idea behind the Perona-Malik algorithm is to evolve an original image, $`I(𝐫)`$, under an edge-controlled diffusion operator $`{\displaystyle \frac{u(𝐫,t)}{t}}`$ $`=`$ $`\left[d(u(𝐫,t))u(𝐫,t)\right]`$ (59) $`u(𝐫,0)`$ $`=`$ $`I(𝐫).`$ (60) Here, $`d(u)`$ is a generalized diffusion coefficient which is so designed that its values are very small at the edges of an image. Many edge stopping functions $`d(\|u\|)`$ are appropriate for anisotropic diffusion. For example, the Gaussian $$d(u)=e^{\|u\|^2/2\sigma ^2}$$ (61) and the Lorentz $$d(u)=\frac{1}{1+\|u\|^2/\sigma ^2}$$ (62) are both suitable for edge representation. They provide perceptually similar results in practical applications. Numerically, the diffusion coefficient becomes very small near an image edge due to the effect of edge stopping functions $`d(u)`$. As a result, the image edge is preserved in the diffusion process. The pixel values at a non-edge part will be smoothed and reduced due to the substantial diffusion coefficient prescribed by $`d(u)`$. Perona and Malik argued that the solution of their anisotropic diffusion equation has no additional maxima (minima) which does not belong to the initial image data. However, this point has been challenged recently. It is well-known that this anisotropic diffusion algorithm may break down when the gradient generated by noise is comparable to image edges and features. Numerically, this can be alleviated by using a regularization procedure. ### B Shock capturing by anisotropic diffusion Hyperbolic conservation laws of the type $$\frac{u(𝐫,t)}{t}+𝐅(u)=0$$ (63) describe the rate of change of a physical quantity $`u`$ given by the generalized convection $`𝐅(u)`$. Without the balance of conduction and/or production, Eq. (63) may have a discontinuous solution. The task is to construct a stable computational scheme which is capable of resolving the “shock”. We first note that computationally, if a shock is defined as a discontinuity, there is no shock to capture. This is because, the original notion of discontinuity is undefined on a discrete mesh. Therefore, a numerical shock is characterized by rapid variation of function values over a small grid zone. Numerical shock capturing, at best, is globally a first order approximation scheme for resolving the large gradient feature of the true solution. However, locally, it is preferred to compute the solution as accurate as possible, except at the discontinuity of the true solution. To this end, we introduce an anisotropic diffusion term to Eq. (63) $`{\displaystyle \frac{u(𝐫,t)}{t}}+𝐅(u)`$ $`=`$ $`\left[d_1(u(𝐫,t))u(𝐫,t)\right]`$ (64) where the diffusivity, $`d_1(u(𝐫,t)`$, is chosen such that it is essentially zero except at a numerical shock position. Obviously, $`u(𝐫,t)`$ is important for shock detection. Apparently, Eq. (64) reduces to the Perona-Malik equation (63) without the convection term. However, the selections of the anisotropic diffusivity in these two equations are entirely different and serve opposite purposes. For example, one may choose $`d_1(u(𝐫,t)`$ as $$d_1=d_1^1\mathrm{ln}[(u)^2+1],$$ (65) where $`d_1^1`$ is a constant. Equation (65) differs much from the Gaussian and the Lorentz. The derivation of Eq. (64) has its roots in the conservation of a physical quantity involving a phenomenological flux and its divergence. However, as a computational algorithm, an anisotropic diffusion of the form $`{\displaystyle \frac{u(𝐫,t)}{t}}+𝐅(u)`$ $`=`$ $`\mathrm{\Gamma }_1(u(𝐫,t))^2u(𝐫,t)`$ (66) can be used. In fact, expressions in both Eqs. (64) and (66) are efficient for numerical shock capturing. Here, $`\mathrm{\Gamma }_1`$ is chosen to smooth the oscillations near a shock and is essentially zero in other regions. ### C Edge enhancing functional and super diffusion From the hydrodynamic point of view, the governing equation of a conservative quantity has a general structure, i.e., the rate of change is balanced by convection, conduction and production. Therefore, the nonlinear advective motion can be counterbalanced by an appropriate production. Therefore, we propose a real-valued, bounded shock (edge) enhancing functional $$e(u(𝐫,t)).$$ (67) Here $`e(u(𝐫,t))`$ is appropriately chosen so that it is edge sensitive and is essentially zero away from a numerical shock or “discontinuity”. This leads to another shock capturing equation $`{\displaystyle \frac{u(𝐫,t)}{t}}+𝐅(u)`$ $`=`$ $`\left[d_1(u(𝐫,t))u(𝐫,t)\right]`$ (68) $`+`$ $`e(u(𝐫,t)).`$ (69) Appropriate choice of the shock enhancing functional will result in a stable numerical algorithm. Obviously, the von Neumann and Richtmyer artificial viscosity term, Eq. (1), is a special case of the present formulation. The diffusion equation can be derived from Fick’s law for mass flux, $$𝐣_\mathrm{𝟏}(𝐫,t)=D_1u(𝐫,t)$$ (70) with $`D_1`$ being a constant. From the point of view of kinetic theory, this is an approximation to a quasi homogeneous system which is near equilibrium. A better approximation can be expressed as a super flux $$𝐣_𝐪(𝐫,t)=\underset{q}{}D_q^{2q}u(𝐫,t),(q=1,2,\mathrm{}),$$ (71) where $`D_q`$ are constants and higher order terms ($`q>1`$) describe corrections to mass flux by the influence of inhomogeneity in density distribution and of flux-flux correlations. The mass conservation leads to $`{\displaystyle \frac{u(𝐫,t)}{t}}`$ $`=`$ $`𝐣_𝐪(𝐫,t)+s(𝐫,t)`$ (72) $`=`$ $`{\displaystyle \underset{q}{}}[D_q^{2q}u(𝐫,t)]+s(𝐫,t),(q=1,2,\mathrm{}),`$ (73) where $`s`$ is a source term which can be a nonlinear function describing chemical reactions. Equation (72) is a generalized reaction-diffusion equation which includes not only the usual diffusion and production terms, but also super diffusion terms. A truncation at the second order super flux, $`𝐣_\mathrm{𝟐}(𝐫,t)=D_1u(𝐫,t)D_2^2u(𝐫,t)`$, leads to the expression that is important in many phenomenological theories, such as the Cahn-Hilliard equation and the Kuramoto-Sivashinsky equation. The latter have been used for the description of a number of physical phenomena, such as the Taylor-Couette flow, parametric waves and pattern formation in alloys, glasses, polymers, combustion and biological systems. In a shock wave simulation process, the distribution of the physical quantity under study can be highly inhomogeneous and/or oscillatory. Hence, the present shock capturing algorithm can be made more efficient by incorporating a shock (edge) sensitive super diffusion operator $`{\displaystyle \frac{u(𝐫,t)}{t}}+𝐅(u)`$ $`=`$ $`{\displaystyle \underset{q}{}}\left[d_q(u(𝐫,t))^{2q}u(𝐫,t)\right]`$ (74) $`+`$ $`e(u(𝐫,t)),(q=1,2,\mathrm{}).`$ (75) Here $`d_q(u,\|u\|)`$ are shock (edge) sensitive diffusion functions. In most applications, a truncation at $`q=2`$ is sufficient $`{\displaystyle \frac{u(𝐫,t)}{t}}+𝐅(u)`$ $`=`$ $`[d_1(u(𝐫,t))u(𝐫,t)]`$ (76) $`+`$ $`[d_2(u(𝐫,t))^2u(𝐫,t)]`$ (77) $`+`$ $`e(u(𝐫,t)).`$ (78) This equation can be simplified further for practical implementation $`{\displaystyle \frac{u(𝐫,t)}{t}}+𝐅(u)`$ $`=`$ $`\mathrm{\Gamma }_1(u(𝐫,t))^2u(𝐫,t)`$ (79) $`+`$ $`\mathrm{\Gamma }_2(u(𝐫,t))^4u(𝐫,t)`$ (80) $`+`$ $`e(u(𝐫,t)),`$ (81) where $`\mathrm{\Gamma }_1`$ and $`\mathrm{\Gamma }_2`$ are to be appropriately chosen to capture the shock. Usually, $`\mathrm{\Gamma }_1`$ is a positive function, while $`\mathrm{\Gamma }_2`$ is negative. From the point view of image processing, both operators $`^2`$ and $`^4`$ are high pass filters. In the Fourier representation, operators $`^2`$ and $`^4`$ are proportional to $`\omega ^2`$ and $`\omega ^4`$ of the frequency $`\omega `$. Hence, $`^2`$ is more sensitive to low frequency oscillations while $`^4`$ has a large filter response for high frequency oscillations. Although both $`^2`$ and $`^4`$ become band pass filters on a discrete grid, their filter responses have similar properties as in their abstract operator forms. Note that many of these anisotropic terms presented in this section resemble those expressions derived in the last section for the artificial pressure tensor, $`\text{P}_{\mathrm{art}}`$. Although the starting points for these two approaches are entirely different, the results are obviously connected, notation needs to be simplified and possible confusions are to be avoided. For these reasons, we use the acronym ADOR for both approaches. ## IV Numerical experiments In this section, the numerical schemes developed from the analysis of kinetic theory and image processing analysis are utilized for numerical computations involving shock waves. Parameter selections for the ADOR method are specified. Obviously, there are infinitely many ways to construct anisotropic diffusion functionals. For simplicity, we test the following choices of the ADOR coefficients: $$d_1=d_1^1\mathrm{ln}[(u)^2+1],d_2=d_2^1\mathrm{ln}[(u)^2+1],e=e^1(u)^2$$ (82) for Eq. (76), $$\mathrm{\Gamma }_1=\gamma _1^1(\|u_x\|)^{\frac{1}{4}},\mathrm{\Gamma }_2=\gamma _2^1(\|u_x\|)^{\frac{1}{4}},e=0$$ (83) for Eq. (79), $$\mathrm{\Gamma }_1=\gamma _1^2\mathrm{ln}[\|u_x\|+1],\mathrm{\Gamma }_2=\gamma _2^2\mathrm{ln}[\|u_x\|+1],e=0,$$ (84) for Eq. (79), $$\mathrm{\Gamma }_1=\gamma _1^3\mathrm{ln}[u_x^2+1],\mathrm{\Gamma }_2=\gamma _2^3\mathrm{ln}[u_x^2+1],e=0,$$ (85) for Eq. (79), and $$\mathrm{\Gamma }_1=\gamma _1^4\mathrm{ln}[u+1],\mathrm{\Gamma }_2=\gamma _2^4\mathrm{ln}[u+1],e=0,$$ (86) for Eq. (79). In fact, coefficients in x- and y-direction can be chosen separately $$\mathrm{\Gamma }_1=\gamma _1^5\{\begin{array}{cc}\mathrm{ln}[\|u_x\|+1]\hfill & \text{for }u_{xx}\hfill \\ \mathrm{ln}[\|u_y\|+1]\hfill & \text{for }u_{yy}\hfill \end{array},\gamma _2^5=0,e=0.$$ (87) Note that these expressions have one feature in common, i.e., they are all of low order in the gradient of $`u`$. This feature allows a sharp change in the solution to be handled by a coarse grid. Otherwise, if there are higher-order gradient terms, such as the von Neumann and Richtmyer form, the computational grid near the shock front has to be refined to reduce the flux amplitude. A finer grid, in turn, leads to another stability problem as dictated by the CFL condition. In the rest of this section, the performance of the abovementioned prescriptions is examined for a detailed choice of ADOR parameters $`d_1^1,d_2^1,e^1`$ and $`\gamma _j^i(i=1,2,3;j=1,2)`$. A few standard problems are employed to test the validity, and to demonstrate the robustness of the present approach. These problems include Burgers’ equation in one and two space dimension, the incompressible Navier-Stokes equation and Euler equation with periodic boundary conditions. For spatial discretization, we use the discrete singular convolution (DSC) algorithm which was proposed as a potential approach for computer realization of singular convolutions. Mathematical foundation of the algorithm is the theory of distributions. Sequence of approximations to the singular kernels of Hilbert type, Abel type and delta type were constructed. Applications are discussed to analytical signal processing, Radon transform and surface interpolation. Numerical solutions to differential equations are formulated via singular kernels of delta type. By appropriately choosing the DSC kernels, the DSC approach exhibits global methods’ accuracy for integration and local methods’ flexibility for handling complex geometries and boundary conditions. Unified features of the DSC approach to differential equations were analyzed in details. In particular, we demonstrated that different implementations of the DSC algorithm, such as global, local, Galerkin, collocation, and finite difference, can be deduced from a single starting point. The DSC algorithm is validated for the numerical solution of the Fokker-Planck equation, the Schrödinger equation and the Navier-Stokes equation. It was also utilized to integrate the (nonlinear) sine-Gordon equation with the initial values close to a homoclinic orbit singularity, for which conventional local methods encounter great difficulties and numerically induced chaos was reported for such an integration. The reader is referred to for a detailed discussion of the method. The DSC parameters used in the present work are $`\sigma /\mathrm{\Delta }=3.2`$ and $`M=31`$ for the DSC kernel of regularized Shannon. For time discretization, the explicit 4th-order Runge-Kutta scheme is used for Burgers’ equation. The Navier-Stokes equation is integrated by using the implicit Euler scheme in association with a discrete singular convolution-alternating direction implicit (DSC-ADI) algorithm. Further information is given in subsections below. ### A Burgers’ equation in one dimension Burgers’ equation is an important model for the understanding of physical flows. It appears customary to test new schemes in computational fluid dynamics by applying them to Burgers’ equation. Despite much of the effort, numerical solution of Burgers’ equation is still not a trivial task, particularly at very high Reynolds numbers where the nonlinear advection leads to shock waves. In fact, many standard computational algorithms fail to predict Burgers’ inviscid shocks. Burgers’ equation is given by $$\frac{u}{t}+u\frac{u}{x}=\frac{1}{\mathrm{Re}}\frac{^2u}{x^2},$$ (88) where $`u(x,t)`$ is the dependent variable resembling the flow velocity and Re is the Reynolds number characterizing the size of the viscosity. The competition between the nonlinear advection and the viscous diffusion is controlled by the value of Re in Burgers’ equation, and thus determines the behavior of the solution. We consider Eq. (88) using the following initial and boundary conditions $`u(x,0)`$ $`=`$ $`\mathrm{sin}(\pi x),`$ (89) $`u(0,t)`$ $`=`$ $`u(1,t)=0.`$ (90) Cole has provided an exact solution for this problem in terms of a series expansion which is readily computable roughly for the parameter $`\mathrm{Re}100`$. For the parameter $`\mathrm{Re}=100`$, the present calculations use 41 grid points in the interval with a time increment of 0.01. Both $`L_1`$ and $`L_{\mathrm{}}`$ errors at 11 different times are listed in TABLE I. We next consider inviscid Burgers’ equation with same initial and boundary condition as given in Eq. (89). In this case, the solution quickly develops into a sharp shock front at $`x=1`$. The DSC algorithm does not work along because severe oscillations eventually turn into an uncontrolled error growth. This problem is treated by using the anisotropic diffusion approach discussed in the previous sections. The performance of four different prescriptions given in Eqs. (82)-(85) is examined. ADOR parameter selections are compared in FIG. 1 at 4 different times (t=0.3,0.5,0.8,2.0). These results are all obtained by using 101 grid points with a time increment of 0.002. In FIG. 1e, results for all the four different forms of coefficients are plotted for a detailed comparison. For the region away from the shock front, there is essentially no difference among four different forms. Although, results computed from different forms are slightly different at the shock front, four different choices of the diffusion coefficients have similar behavior and all are capable of correctly simulating Burgers’ shock. The super diffusion coefficient chosen in Eq. (84) does not work very well as the shock front is distorted (see FIG. 1f). However, we have found that a combination of anisotropic diffusion and super diffusion does work better than the choice of a single term (results not shown). Similar results were also obtained by using 50 grid points. It is mentioned that the tests on the use of the edge enhancing functional does not result in a stable solution with the mesh system and time increment, chosen the present study. ### B Burgers’ equation in two dimensions Let us consider Burgers’ equation of the form $`u_t+uu_x+vu_y={\displaystyle \frac{1}{\mathrm{Re}}}(u_{xx}+u_{yy})`$ (91) $`v_t+uv_x+vv_y={\displaystyle \frac{1}{\mathrm{Re}}}(v_{xx}+v_{yy})`$ (92) in a square $`[0,1]\times [0,1]`$ with the initial values $`u(x,y,0)=\mathrm{sin}(\pi x)\mathrm{sin}(\pi y)`$ (93) $`v(x,y,0)=[\mathrm{sin}(\pi x)+\mathrm{sin}(2\pi x)]+[\mathrm{sin}(\pi y)+\mathrm{sin}(2\pi y)]`$ (94) and boundary conditions $`u(0,y,t)=u(1,y,t)=u(x,0,t)=u(x,1,t)=0`$ (95) $`v(0,y,t)=v(1,y,t)=v(x,0,t)=v(x,1,t)=0.`$ (96) This problem has no analytical solution and is chosen to demonstrate that the present algorithm can be efficient, even with a very coarse mesh. This case is computed by using an ADOR parameter of $`\gamma _1^3=0.0006`$ with $`41^2`$ points. The contours of velocity field components at t=1 are plotted in FIG. 2. ### C The incompressible Navier-Stokes equations To test the present approach for shock capturing further we consider the Navier-Stokes equation $`u_t+uu_x+vu_y=p_x+{\displaystyle \frac{1}{\mathrm{Re}}}(u_{xx}+u_{yy})`$ (97) $`v_t+uv_x+vv_y=p_y+{\displaystyle \frac{1}{\mathrm{Re}}}(v_{xx}+v_{yy})`$ (98) with equation of continuity $$u_x+v_y=0,$$ (99) where $`(u,v)`$ is the velocity vector, $`p`$ is the pressure, Re (Re$`>0`$) is the Reynolds number and Re$`=\mathrm{}`$ defines the Euler equation. The domain of problem is a square $`[0,2\pi ]\times [0,2\pi ]`$ with periodic boundary conditions. With appropriate initial values, the Euler equation can be used to describe a flow field of vertically perturbed horizontal shear layers around a jet. Bell et al studied this case by a second order projection method. Recently E and Shu have employed this example to demonstrate the success of their high order ENO scheme for resolving the fine vorticity structure of the double shear layers. Case 1: Analytically solvable initial values The Navier-Stokes equation is analytically solvable for appropriate initial values $`u(x,y,0)=\mathrm{cos}(x)\mathrm{sin}(y)`$ (100) $`v(x,y,0)=\mathrm{sin}(x)\mathrm{cos}(y).`$ (101) The exact solution for this case is given by $`u(x,y,t)=\mathrm{cos}(x)\mathrm{sin}(y)e^{\frac{2t}{\mathrm{Re}}}`$ (102) $`v(x,y,t)=\mathrm{sin}(x)\mathrm{cos}(y)e^{\frac{2t}{\mathrm{Re}}}`$ (103) $`p(x,y,t)={\displaystyle \frac{1}{4}}[\mathrm{cos}(2x)+\mathrm{cos}(2y)]e^{\frac{4t}{\mathrm{Re}}}.`$ (104) This provides a benchmark test for potential numerical methods in fluid dynamics. The implicit Euler scheme is used for the time integration and the DSC algorithm is utilized for the spatial discretization. The accuracy and reliability of this combination was previously tested for this problem using a standard LU decomposition algorithm for solving linear algebraic equations. Here, DSC-ADI algorithm is used. We choose a grid of $`32^2`$ for the present calculation with a time increment of 0.001. The DSC-ADI results are summarized in TABLE II. Note that, for the inviscid case (Re=$`\mathrm{}`$) the present result is accurate to the machine precision. It is evident that the accuracy of the DSC approach is extremely high, particularly when the Reynolds numbers are very large. Case 2: The Euler equation We now test our anisotropic diffusion approach for the Euler equation (Re$`=\mathrm{}`$) with sharply varying initial values. This example is chosen to illustrate the ability of the present approach, for providing very fine resolution with a relatively coarse grid. The initial values are that of a jet in a doubly periodic geometry $`u(x,y,0)=\left\{\begin{array}{cc}\mathrm{tanh}\left(\frac{2y\pi }{2\rho }\right),\hfill & \text{if }y\pi \hfill \\ \mathrm{tanh}\left(\frac{3\pi 2y}{2\rho }\right),\hfill & \text{if }y>\pi \hfill \end{array}\right\}`$ (107) $`v(x,y,0)=\delta \mathrm{sin}(x),`$ (108) where $`\delta =0.05`$ is used for the convenience of comparison with the previous study. This initial value describes the flow field consisting of horizontal shear layers of finite thickness, perturbed by a small amplitude vertical velocity, making up the boundaries of the jet. However, this problem is not analytically solvable. A pioneer work in this study was given by Bell et al, in which they utilized a second-order Godunov scheme in association with a projection approach for divergence-free velocity fields with general boundary condition. With a periodic boundary condition, E and Shu have shown that their high-order ENO scheme performs well for this problem. We first consider the parameter $`\rho =\pi /15`$, a case studied by Bell et al using a projection method with three sets of grids (128<sup>2</sup>, 256<sup>2</sup> and 512<sup>2</sup>). E and Shu computed this case by using both spectral collocation code with 512<sup>2</sup> points and their high order ENO scheme with 64<sup>2</sup> and 128<sup>2</sup> points. The spectral collocation code produced an oscillatory solution at $`t=10`$ (see FIG. 1 of Ref. ), while the high order ENO scheme produced a defect at $`t=6`$ as the channels connecting the vorticity centers are slightly distorted (see FIG. 2 of Ref. ). In the present simulation, we choose a $`64^2`$ grid for the computational domain with a time increment of 0.002. The prescription in Eq. (85) is used with the ADOR parameter of $`\gamma _1^3=0.0006`$. The results at different times (t=4,6,8,10) are plotted in FIG 3. It is seen that our solutions are smooth (some non-smooth features in the contour plot is due to the fact that the grid is very coarse) and stable for this case. In particular, no distortion is found in vorticity contours at t=6. For early times, present results compare extremely well with those of the spectral collocation code computed with 512<sup>2</sup> points. There are no spurious numerical oscillations during the entire process. Finally, we perform simulations by the present approach using the discontinuous initial data ($`\rho 0`$) as in . The evolution under the Euler equation leads to a periodic array of large vortices, with the shear layers between the rolls being thinned by the large straining field. It is known that, for the present choice of parameters the solution quickly develops into a roll-up process with smaller and smaller scales, and the resolution is lost eventually with a fixed grid for local methods. We compute this case by using $`128^2`$ grid points as in . A time increment of 0.002 is used. The ADOR parameter is chosen as $`\gamma _1^3=0.0006`$. The present results at different times (t=4,6,8,10,12,14) are plotted in FIG 4. ## V Conclusions Connection is made between digital image processing and computational fluid dynamics. The evolution of an image surface under a partial differential operator can be viewed as a form of image processing. Computationally, numerical shock capturing can be formulated on the lines of iterative edge-detection. Hence, techniques developed in the computational fluid dynamics can be used for image processing and vice-versa. This paper introduces the method of anisotropic diffusion oscillation reduction (ADOR), an approach which has its roots in image processing, for shock wave computations. In fact, the ADOR method is much similar to the artificial viscosity algorithm. Physical origins and mathematical properties of the artificial viscosity are discussed from the kinetic theory point of view. The form of pressure tensor is derived from the first principles of quantum mechanics. Quantum kinetic theory is utilized to arrive at macroscopic transport equations from the microscopic quantum theory. Macroscopic symmetry is used to simplify the phenomenological pressure tensor expressions. The latter provides the basis for the design of artificial viscosity. The original von Neumann-Richtmyer form of artificial viscosity fits into a special case of the present generalizations. The anisotropic diffusion, which is essentially an image processing technique, is modified for shock capturing. The technique preserves image edges by introducing little diffusion, where the image gradient is large, while it provides a substantial diffusion coefficient at smooth parts of the image. In the present shock capturing algorithm, the edge-sensitive diffusion is introduced at a shock front so that large oscillations can be efficiently eliminated. An edge enhancing functional and edge-detected super diffusion operators are proposed for shock capturing. These terms are introduced from the general structure of conservation laws. In fact, kinetic theory analysis and anisotropic diffusion argument lead to a number of similar expressions for shock wave treatments. Hence, the acronym ADOR is referred for both approaches. The reliability and robustness of the ADOR method is explored in association with the discrete singular convolution (DSC) algorithm. The DSC algorithm was previously validated by handling many linear and nonlinear problems and is a potential algorithm for the computer realization of some singular convolutions. A few detailed prescriptions for the anisotropic diffusion functional are considered in this work. The coefficients in these prescriptions are chosen to be of lower order in gradient so that a coarse grid can be used. A number of standard test examples, including (inviscid) Burgers’ equation in one and two spatial dimensions, the incompressible Navier-Stokes and the Euler equations, are employed for the present test computations. For Burgers’ equation, we have examined the accuracy for spatial and temporal discretizations. A few ADOR coefficients given by Eqs. (82)-(85) are tested. The anisotropic diffusion term is found to be very robust in association with many edge-detected coefficients, while the super diffusion operator produces an over shot at Burgers’ shock front. The edge enhancing functional apparently does not perform well for this problem. The ADOR approach is found to be very stable for the 2D inviscid Burgers’ equation even with a very coarse grid. For the incompressible Navier-Stokes equation, the ADOR approach is used in association with a DSC-ADI algorithm. The accuracy of the algorithm was tested by an analytically solvable case and the machine precision is found at the inviscid limit of the Navier-Stokes equation. The algorithm was then applied to the study of doubly periodic shear layers, which was chosen to demonstrate the capability of the ADOR method for more difficult problems. These results indicate that the proposed method has the potential to capture shock waves. Work on the implementation of the present approach to more complex fluid flow systems, including general boundary conditions and for compressible flow is under progress. Acknowledgment This work was supported in part by the National University of Singapore. The author thanks Professors Tony Chan and C.-W. Shu for useful discussions about the ENO and WENO schemes. Figure Captions FIG. 1. A comparison of various ADOR coefficients for the numerical solution of inviscid Burgers’ equation. (a) $`d_1^1=0.0009,d_2^1=0,e^1=0`$; (b) $`\gamma _1^1=0.0022,\gamma _2^1=0`$; (c) $`\gamma _1^2=0.0018,\gamma _2^2=0`$; (d) $`\gamma _1^3=0.0015,\gamma _2^3=0`$; (e) Comparison of (a), (b),(c) and (d); (f) $`\gamma _1^3=0,\gamma _2^3=98\times 10^8`$. FIG. 2. The ADOR solution for the 2D inviscid Burgers’ equation with 41<sup>2</sup> points, $`\gamma _1^4=0.02,\gamma _2^4=0`$. Left: the $`u`$ field; right: the $`v`$ field. FIG. 3. The vorticity contours for the 2D Euler equation by the ADOR method with 64<sup>2</sup> points, $`\gamma _1^5=0.0006`$. Up left: t=4; up right: t=6; low left: t=8; low right: t=10. FIG. 4. The vorticity contour for the 2D Euler equation, with discontinuous initial data, by the ADOR method with 128<sup>2</sup> points, $`\gamma _1^5=0.0006`$. Up left: t=4; up right: t=6; middle left: t=8; middle right: t=10; low left: t=12; low right: t=14. FIG. 1. A comparison of various ADOR coefficients for the numerical solution of inviscid Burgers’ equation. (a) $`d_1^1=0.0009,d_2^1=0,e^1=0`$; (b) $`\gamma _1^1=0.0022,\gamma _2^1=0`$; (c) $`\gamma _1^2=0.0018,\gamma _2^2=0`$; (d) $`\gamma _1^3=0.0015,\gamma _2^3=0`$; (e) Comparison of (a), (b),(c) and (d); (f) $`\gamma _1^3=0,\gamma _2^3=98\times 10^8`$. FIG. 2. The ADOR solution for the 2D inviscid Burgers’ equation with 41<sup>2</sup> points, $`\gamma _1^4=0.02,\gamma _2^4=0`$. Left: the $`u`$ field; right: the $`v`$ field. FIG. 3. The vorticity contours for the 2D Euler equation by the ADOR method with 64<sup>2</sup> points, $`\gamma _1^5=0.0006`$. Up left: t=4; up right: t=6; low left: t=8; low right: t=10. FIG. 4. The vorticity contour for the 2D Euler equation, with discontinuous initial data, by the ADOR method with 128<sup>2</sup> points, $`\gamma _1^5=0.0006`$. Up left: t=4; up right: t=6; middle left: t=8; middle right: t=10; low left: t=12; low right: t=14.
warning/0004/quant-ph0004045.html	ar5iv	text	# Relative entropy in quantum information theory ## Abstract We review the properties of the quantum relative entropy function and discuss its application to problems of classical and quantum information transfer and to quantum data compression. We then outline further uses of relative entropy to quantify quantum entanglement and analyze its manipulation. ## 1 Quantum relative entropy In this paper we discuss several uses of the quantum relative entropy function in quantum information theory. Relative entropy methods have a number of advantages. First of all, the relative entropy functional satisfies some strong identities and inequalities, providing a basis for good theorems. Secondly, the relative entropy has a natural interpretation in terms of the statistical distinguishability of quantum states; closely related to this is the picture of relative entropy as a “distance” measure between density operators. These interpretations of the relative entropy give insight about the meaning of the mathematical constructions that use it. Finally, relative entropy has found a wide variety of applications in quantum information theory. The usefulness of relative entropy in quantum information theory should come as no surprise, since the classical relative entropy has shown its power as a unifying concept in classical information theory . Indeed, some of the results we will describe have close analogues in the classical domain. Nevertheless, the quantum relative entropy can provide insights in contexts (such as the quantification of quantum entanglement) that have no parallel in classical ideas of information. Let $`Q`$ be a quantum system described by a Hilbert space $``$. (Throughout this paper, we will restrict our attention to systems with Hilbert spaces having a finite number of dimensions.) A pure state of $`Q`$ can be described by a normalized vector $`\|\psi `$ in $``$, but a general (mixed) state requires a density operator $`\rho `$, which is a positive semi-definite operator on $``$ with unit trace. For the pure state $`\|\psi `$, the density operator $`\rho `$ is simply the projection operator $`\|\psi \psi \|`$; otherwise, $`\rho `$ is a convex combination of projections. The entropy $`S(\rho )`$ is defined to be $$S(\rho )=\text{Tr}\rho \mathrm{log}\rho .$$ (1) The entropy is non-negative and equals zero if and only if $`\rho `$ is a pure state. (By “$`\mathrm{log}`$” we will mean a logarithm with base 2.) Closely related to the entropy of a state is the relative entropy of a pair of states. Let $`\rho `$ and $`\sigma `$ be density operators, and define the quantum relative entropy $`𝒮(\rho \|\|\sigma )`$ to be $$𝒮(\rho \|\|\sigma )=\text{Tr}\rho \mathrm{log}\rho \text{Tr}\rho \mathrm{log}\sigma .$$ (2) (We read this as “the relative entropy of $`\rho `$ with respect to $`\sigma `$”.) This function has a number of useful properties: 1. $`𝒮(\rho \|\|\sigma )0`$, with equality if and only if $`\rho =\sigma `$. 2. $`𝒮(\rho \|\|\sigma )<\mathrm{}`$ if and only if $`\text{supp}\rho \text{supp}\sigma `$. (Here “$`\text{supp}\rho `$” is the subspace spanned by eigenvectors of $`\rho `$ with non-zero eigenvalues.) 3. The relative entropy is continuous where it is not infinite. 4. The relative entropy is jointly convex in its arguments . That is, if $`\rho _1`$, $`\rho _2`$, $`\sigma _1`$ and $`\sigma _2`$ are density operators, and $`p_1`$ and $`p_2`$ are non-negative numbers that sum to unity (i.e., probabilities), then $$𝒮(\rho \|\|\sigma )p_1𝒮(\rho _1\|\|\sigma _1)+p_2𝒮(\rho _2\|\|\sigma _2)$$ (3) where $`\rho =p_1\rho _1+p_2\rho _2`$ and $`\sigma =p_1\sigma _1+p_2\sigma _2`$. Joint convexity automatically implies convexity in each argument, so that (for example) $$𝒮(\rho \|\|\sigma )p_1𝒮(\rho _1\|\|\sigma )+p_2𝒮(\rho _2\|\|\sigma ).$$ (4) The properties, especially property (1), motivate us to think of the relative entropy as a kind of “distance” between density operators. The relative entropy, which is not symmetric and which lacks a triangle inequality, is not technically a metric; but it is a positive definite directed measure of the separation of two density operators. Suppose the density operator $`\rho _k`$ occurs with probability $`p_k`$, yielding an average state $`\rho ={\displaystyle \underset{k}{}}p_k\rho _k`$, and suppose $`\sigma `$ is some other density operator. Then $`{\displaystyle \underset{k}{}}p_k𝒮(\rho _k\|\|\sigma )`$ $`=`$ $`{\displaystyle \underset{k}{}}p_k\left(\text{Tr}\rho _k\mathrm{log}\rho _k\text{Tr}\rho _k\mathrm{log}\sigma \right)`$ $`=`$ $`{\displaystyle \underset{k}{}}p_k\left(\text{Tr}\rho _k\mathrm{log}\rho _k\text{Tr}\rho _k\mathrm{log}\rho +\text{Tr}\rho _k\mathrm{log}\rho \text{Tr}\rho _k\mathrm{log}\sigma \right)`$ $`=`$ $`{\displaystyle \underset{k}{}}p_k\left(\text{Tr}\rho _k\mathrm{log}\rho _k\text{Tr}\rho _k\mathrm{log}\rho \right)+\text{Tr}\rho \mathrm{log}\rho \text{Tr}\rho \mathrm{log}\sigma `$ $`{\displaystyle \underset{k}{}}p_k𝒮(\rho _k\|\|\sigma )`$ $`=`$ $`{\displaystyle \underset{k}{}}p_k𝒮(\rho _k\|\|\rho )+𝒮(\rho \|\|\sigma ).`$ (5) Equation 5 is known as Donald’s identity. The classical relative entropy of two probability distributions is related to the probability of distinguishing the two distributions after a large but finite number of independent samples. This is called Sanov’s theorem , and this result has quantum analogue . Suppose $`\rho `$ and $`\sigma `$ are two possible states of the quantum system $`Q`$, and suppose we are provided with $`N`$ identically prepared copies of $`Q`$. A measurement is made to determine whether the prepared state is $`\rho `$, and the probability $`P_N`$ that the state $`\sigma `$ passes this test—in other words, is confused with $`\rho `$—is $$P_N2^{N𝒮(\rho \|\|\sigma )}$$ (6) as $`N\mathrm{}`$. (We have assumed that the measurement made is an optimal one for the purpose, and it is possible to show that an asymptotically optimal measurement strategy can be found that depends on $`\rho `$ but not $`\sigma `$.) The quantum version of Sanov’s theorem tells us that the quantum relative entropy governs the asymptotic distinguishability of one quantum state from another by means of measurements. This further supports the view of $`𝒮(\|\|)`$ as a measure of “distance”; two states are “close” if they are difficult to distinguish, but “far apart” if the probability of confusing them is small. The remainder of this paper is organized as follows. Sections 2–5 apply relative entropy methods to the problem of sending classical information by means of a (possibly noisy) quantum channel. Sections 6–7 consider the transmission and compression of quantum information. Sections 8–9 then apply relative entropy methods to the discussion of quantum entanglement and its manipulation by local operations and classical communication. We conclude with a few remarks in Section 10. ## 2 Classical communication via quantum channels One of the oldest problems in quantum information theory is that of sending classical information via quantum channels. A sender (“Alice”) wishes to transmit classical information to a receiver (“Bob”) using a quantum system as a communication channel. Alice will represent the message $`a`$, which occurs with probability $`p_a`$, by preparing the channel in the “signal state” represented by the density operator $`\rho _a`$. The average state of the channel will thus be $`\rho ={\displaystyle \underset{a}{}}p_a\rho _a`$. Bob will attempt to recover the message by making a measurement of some “decoding observable” on the channel system. The states $`\rho _a`$ should be understood here as the “output” states of the channel, the states that Bob will attempt to distinguish in his measurement. In other words, the states $`\rho _a`$ already include the effects of the dynamical evolution of the channel (including noise) on its way from sender to receiver. The dynamics of the channel will be described by a trace-preserving, completely positive map $``$ on density operators . The effect of $``$ is simply to restrict the set of output channel states that Alice can arrange for Bob to receive. If $`𝒟`$ is the set of all density operators, then Alice’s efforts can only produce output states in the set $`𝒜=(𝒟)`$, a convex, compact set of density operators. Bob’s decoding observable is represented by a set of positive operators $`E_b`$ such that $`{\displaystyle \underset{b}{}}E_b=1`$. If Bob makes his measurement on the state $`\rho _a`$, then the conditional probability of measurement outcome $`b`$ is $$P(b\|a)=\text{Tr}\rho _aE_b.$$ (7) This yields a joint distribution over Alice’s input messages $`a`$ and Bob’s decoded messages $`b`$: $$P(a,b)=p_aP(b\|a).$$ (8) Once a joint probability distribution exists between the input and output messages (random variables $`A`$ and $`B`$, respectively), the information transfer can be analyzed by classical information theory. The information obtained by Bob is given by the mutual information $`I(A:B)`$: $$I(A:B)=H(A)+H(B)H(A,B)$$ (9) where $`H`$ is the Shannon entropy function $$H(X)=\underset{x}{}p(x)\mathrm{log}p(x).$$ (10) Shannon showed that, if the channel is used many times with suitable error-correcting codes, then any amount of information up to $`I(A:B)`$ bits (per use of the channel) can be sent from Alice to Bob with arbitrarily low probability of error . The classical capacity of the channel is $`C=\mathrm{max}I(A:B)`$, where the maximum is taken over all input probability distributions. $`C`$ is thus the maximum amount of information that may be reliably conveyed per use of the channel. In the quantum mechanical situation, for a given ensemble of signal states $`\rho _a`$, Bob has many different choices for his decoding observable. Unless the signal states happen to be orthogonal, no choice of observable will allow Bob to distinguish perfectly between them. A theorem stated by Gordon and Levitin and first proved by Holevo states that the amount of information accessible to Bob is limited by $`I(A:B)\chi `$, where $$\chi =S(\rho )\underset{a}{}p_aS(\rho _a).$$ (11) The quantity $`\chi `$ is non-negative, since the entropy $`S`$ is concave. More recently, Holevo and Schumacher and Westmoreland have shown that this upper bound on $`I(A:B)`$ is asymptotically achievable. If Alice uses the same channel many times and prepares long codewords of signal states, and Bob uses an entangled decoding observable to distinguish these codewords, then Alice can convey to Bob up to $`\chi `$ bits of information per use of the channel, with arbitrarily low probability of error. (This fact was established for pure state signals $`\rho _a=\|\psi _a\psi _a\|`$ in . In this case, $`\chi =S(\rho )`$.) The Holevo bound $`\chi `$ can be expressed in terms of the relative entropy: $`\chi `$ $`=`$ $`\text{Tr}\rho \mathrm{log}\rho +{\displaystyle \underset{a}{}}p_a\text{Tr}\rho _a\mathrm{log}\rho _a`$ $`=`$ $`{\displaystyle \underset{a}{}}p_a\left(\text{Tr}\rho _a\mathrm{log}\rho _a\text{Tr}\rho _a\mathrm{log}\rho \right)`$ $`\chi `$ $`=`$ $`{\displaystyle \underset{a}{}}p_a𝒮(\rho _a\|\|\rho ).`$ (12) In geometric terms, $`\chi `$ is the average relative entropy “directed distance” from the average state $`\rho `$ to the members of the signal ensemble. Donald’s identity (Equation 5) has a particularly simple form in terms of $`\chi `$. Given an ensemble and an additional state $`\sigma `$, $$\underset{a}{}p_a𝒮(\rho _a\|\|\sigma )=\chi +𝒮(\rho \|\|\sigma ).$$ (13) This implies, among other things, that $$\chi \underset{a}{}p_a𝒮(\rho _a\|\|\sigma )$$ (14) with equality if and only if $`\sigma =\rho `$, the ensemble average state. ## 3 Thermodynamic cost of communication In this section and the next, we focus on the transfer of classical information by means of a quantum channel. Imagine a student who attends college far from home . Naturally, the student’s family wants to know that the student is passing his classes, and so they want the student to report to them frequently over the telephone. But the student is poor and cannot affort very many long-distance telephone calls. So they make the following arrangement: each evening at the same time, the poor student will call home only if he is failing one or more of this classes. Otherwise, he will save the phone charges by not calling home. Every evening that the poor student does not call, therefore, the family is receiving a message via the telephone that his grades are good. (That the telephone is being used for this message can be seen from the fact that, if the phone lines are knocked out for some reason, the family can no longer make any inference from the absence of a phone call.) For simplicity, imagine that the student’s grades on successive days are independent and that the probability that the student will be failing on a given evening is $`p`$. Then the information conveyed each evening by the presence or absence of a phone call is $$H(p)=p\mathrm{log}p(1p)\mathrm{log}(1p).$$ (15) The cost of making a phone call is $`c`$, while not making a phone call is free. Thus, the student’s average phone charge is $`cp`$ per evening. The number of bits of information per unit cost is thus $$\frac{H(p)}{cp}=\frac{1}{c}\left(\mathrm{log}p\left(\frac{1}{p}1\right)\mathrm{log}(1p)\right).$$ (16) If the poor student is very successful in his studies, so that $`p0`$, then this ratio becomes unboundedly large, even though both $`H(p)0`$ and $`cp0`$. That is, the student is able to send an arbitrarily large number of bits per unit cost. There is no irreducible cost for sending one bit of information over the telephone. The key idea in the story of the poor student is that one possible signal—no phone call at all—has no cost to the student. The student can exploit this fact to use the channel in a cost-effective way, by using the zero-cost signal almost all of the time. Instead of a poor student using a telephone, we can consider an analogous quantum mechanical problem. Suppose that a sender can manipulate a quantum channel to produce (for the receiver) one of two possible states, $`\rho _0`$ or $`\rho _1`$. The state $`\rho _0`$ can be produced at “zero cost”, while the state $`\rho _1`$ costs a finite amount $`c_1>0`$ to produce. In the signal ensemble, the signal state $`\rho _1`$ is used with probability $`\eta `$ and $`\rho _0`$ with probability $`1\eta `$, leading to an average state $$\rho =(1\eta )\rho _0+\eta \rho _1.$$ (17) The average cost of creating a signal is thus $`c=\eta c_1`$. For this ensemble, $$\chi =(1\eta )𝒮(\rho _0\|\|\rho )+\eta 𝒮(\rho _1\|\|\rho ).$$ (18) As discussed in the previous section, $`\chi `$ is an asymptotically achievable upper bound for the information transfered by the channel. An upper bound for $`\chi `$ can be obtained from Donald’s identity. Letting $`\rho _0`$ be the “additional” state, $$\chi (1\eta )𝒮(\rho _0\|\|\rho _0)+\eta 𝒮(\rho _1\|\|\rho _0)=\eta 𝒮(\rho _1\|\|\rho _0).$$ (19) Combining this with a simple lower bound, we obtain $$\eta 𝒮(\rho _1\|\|\rho )\chi \eta 𝒮(\rho _1\|\|\rho _0).$$ (20) If we divide $`\chi `$ by the average cost, we find an asymptotically achievable upper bound for the number of bits sent through the channel per unit cost. That is, $$\frac{\chi }{c}\frac{1}{c_1}𝒮(\rho _1\|\|\rho _0).$$ (21) Furthermore, equality holds in the limit that $`\eta 0`$. Thus, $$sup\frac{\chi }{c}=\frac{1}{c_1}𝒮(\rho _1\|\|\rho _0).$$ (22) In short, the relative entropy “distance” between the signal state $`\rho _1`$ and the “zero cost” signal $`\rho _0`$ gives the largest possible number of bits per unit cost that may be sent through the channel—the “cost effectiveness” of the channel. If the state $`\rho _0`$ is a pure state, or if we can find a usable signal state $`\rho _1`$ whose support is not contained in the support of $`\rho _0`$, then $`𝒮(\rho _1\|\|\rho _0)=\mathrm{}`$ and the cost effectiveness of the channel goes to infinity as $`\eta 0`$. (This is parallel to the situation of the poor student, who can make the ratio of “bits transmitted” to “average cost” arbitrarily large.) What if there are many possible signal states $`\rho _1`$, $`\rho _2`$, etc., with positive costs $`c_1`$, $`c_2`$, and so on? If we assign the probability $`\eta q_k`$ to $`\rho _k`$ for $`k=1,2,\mathrm{}`$ (where $`{\displaystyle \underset{k}{}}q_k=1`$), and use $`\rho _0`$ with probability $`1\eta `$, then we obtain $$\eta \underset{k}{}q_k𝒮(\rho _k\|\|\rho )\chi \eta \underset{k}{}q_k𝒮(\rho _k\|\|\rho _0).$$ (23) The average cost of the channel is $`c=\eta {\displaystyle \underset{k}{}}q_kc_k`$. This means that $$\frac{\chi }{c}\frac{_kq_k𝒮(\rho _k\|\|\rho _0)}{_kq_kc_k}.$$ (24) We now note the following fact about real numbers. Suppose $`a_n,b_n>0`$ for all $`n`$. Then $$\frac{_na_n}{_nb_n}\underset{n}{\mathrm{max}}\frac{a_n}{b_n}.$$ (25) This can be proven by letting $`R=\mathrm{max}(a_n/b_n)`$ and pointing out that $`a_nRb_n`$ for all $`n`$. Then $`{\displaystyle \underset{n}{}}a_n`$ $``$ $`R{\displaystyle \underset{n}{}}b_n`$ $`{\displaystyle \frac{_na_n}{_nb_n}}`$ $``$ $`R.`$ In our context, this implies that $$\frac{_kq_k𝒮(\rho _k\|\|\rho _0)}{_kq_kc_k}\underset{k}{\mathrm{max}}\frac{q_k𝒮(\rho _k\|\|\rho _0)}{q_kc_k}$$ (26) and thus $$\frac{\chi }{c}\underset{k}{\mathrm{max}}\frac{𝒮(\rho _k\|\|\rho _0)}{c_k}.$$ (27) By using only the “most efficient state” (for which the maximum on the right-hand side is achieved) and adopting the “poor student” strategy of $`\eta 0`$, we can show that $$sup\frac{\chi }{c}=\underset{k}{\mathrm{max}}\frac{𝒮(\rho _k\|\|\rho _0)}{c_k}.$$ (28) These general considerations of an abstract “cost” of creating various signals have an especially elegant development if we consider the thermodynamic cost of using the channel. The thermodynamic entropy $`S_\theta `$ is related to the information-theoretic entropy $`S(\rho )`$ of the state $`\rho `$ of the system by $$S_\theta =k\mathrm{ln}2S(\rho ).$$ (29) The constant $`k`$ is Boltzmann’s constant. If our system has a Hamiltonian operator $`H`$, then the thermodynamic energy $`E`$ of the state is the expectation of the Hamiltonian: $$E=H=\text{Tr}\rho H.$$ (30) Let us suppose that we have access to a thermal reservoir at temperature $`T`$. Then the “zero cost” state $`\rho _0`$ is the thermal equilibrium state $$\rho _0=\frac{1}{Z}e^{\beta H},$$ (31) where $`\beta =1/kT`$ and $`Z=\text{Tr}e^{\beta H}`$. ($`Z`$ is the partition function.) The free energy of the system in the presence of a thermal reservoir at temperature $`T`$ is $`F=ETS_\theta `$. For the equilibrium state $`\rho _0`$, $`F_0`$ $`=`$ $`\text{Tr}\rho _0H+kT\mathrm{ln}2\left(\mathrm{log}Z{\displaystyle \frac{\beta }{\mathrm{ln}2}}\text{Tr}\rho _0H\right)`$ (32) $`=`$ $`kT\mathrm{ln}2\mathrm{log}Z`$ The thermodynamic cost of the state $`\rho _1`$ is just the difference $`F_1F_0`$ between the free energies of $`\rho _1`$ and the equilibrium state $`\rho _0`$. But this difference has a simple relation to the relative entropy. First, we note $$\text{Tr}\rho _1\mathrm{log}\rho _0=\mathrm{log}Z\beta \text{Tr}\rho _1H,$$ (33) from which it follows that $`F_1F_0`$ $`=`$ $`\text{Tr}\rho _1H+kT\mathrm{ln}2\text{Tr}\rho _1\mathrm{log}\rho _1+kT\mathrm{ln}2\mathrm{log}Z`$ $`=`$ $`kT\mathrm{ln}2\left(\text{Tr}\rho _1\mathrm{log}\rho _1\text{Tr}\rho _1\mathrm{log}\rho _0\right)`$ $`F_1F_0`$ $`=`$ $`kT\mathrm{ln}2𝒮(\rho _1\|\|\rho _0).`$ (34) If we use the signal state $`\rho _1`$ with probability $`\eta `$, then the average thermodynamic cost is $`f=\eta (F_1F_0)`$. The number of bits sent per unit free energy is therefore $$\frac{\chi }{f}\eta \frac{𝒮(\rho _1\|\|\rho _0)}{f}=\frac{1}{kT\mathrm{ln}2}.$$ (35) The same bound holds for all choices of the state $`\rho _1`$, and therefore for all ensembles of signal states. We can approach this upper bound if we make $`\eta `$ small, so that $$sup\frac{\chi }{f}=\frac{1}{kT\mathrm{ln}2}$$ (36) In short, for any coding and decoding scheme that makes use of the quantum channel, the maximum number of bits that can be sent per unit free energy is just $`(kT\mathrm{ln}2)^1`$. Phrased another way, the minimum free energy cost per bit is $`kT\mathrm{ln}2`$. This analysis can shed some light on Landauer’s principle , which states that the minimum thermodynamic cost of information erasure is $`kT\mathrm{ln}2`$ per bit. From this point of view, information erasure is simply information transmission into the environment, which requires the expenditure of an irreducible amount of free energy. ## 4 Optimal signal ensembles Now we consider $`\chi `$-maximizing ensembles of states from a given set $`𝒜`$ of available (output) states, without regard to the “cost” of each state. Our discussion in Section 2 tells us that the $`\chi `$-maximizing ensemble is the one to use if we wish to maximize the classical information transfer from Alice to Bob via the quantum channel. Call an ensemble that maximizes $`\chi `$ an “optimal” signal ensemble, and denote the maximum value of $`\chi `$ by $`\chi ^{}`$. (The results of this section are developed in more detail in .) The first question is, of course, whether an optimal ensemble exists. It is conceivable that, though there is a least upper bound $`\chi ^{}`$ to the possible values of $`\chi `$, no particular ensemble in $`𝒜`$ achieves it. (This would be similar to the results in the last section, in which the optimal cost effectiveness of the channel is only achieved in a limit.) However, an optimal ensemble does exist. Uhlmann has proven a result that goes most of the way. Suppose our underlying Hilbert space $``$ has dimension $`d`$ and the set $`𝒜`$ of available states is convex and compact. Then given a fixed average state $`\rho `$, there exists an ensemble of at most $`d^2`$ signal states $`\rho _a`$ that achieves the maximum value of $`\chi `$ for that particular $`\rho `$. The problem we are considering is to maximize $`\chi `$ over all choices of $`\rho `$ in $`𝒜`$. Since Uhlmann has shown that each $`\rho `$-fixed optimal ensemble need involve no more than $`d^2`$ elements, we only need to maximize $`\chi `$ over ensembles that contain $`d^2`$ or fewer members. The set of such ensembles is compact and $`\chi `$ is a continuous function on this set, so $`\chi `$ achieves its maximum value $`\chi ^{}`$ for some ensemble with at most $`d^2`$ elements. Suppose that the state $`\rho _a`$ occurs with probability $`p_a`$ in some ensemble, leading to the average state $`\rho `$ and a Holevo quantity $`\chi `$. We will now consider how $`\chi `$ changes if we modify the ensemble slightly. In the modified ensemble, a new state $`\omega `$ occurs with probability $`\eta `$ and the state $`\rho _a`$ occurs with probability $`(1\eta )p_a`$. For the modified ensemble, $`\rho ^{}`$ $`=`$ $`\eta \omega +(1\eta )\rho `$ (37) $`\chi ^{}`$ $`=`$ $`\eta 𝒮(\omega \|\|\rho ^{})+(1\eta ){\displaystyle \underset{a}{}}p_a𝒮(\rho _a\|\|\rho ^{}).`$ (38) We can apply Donald’s identity to these ensembles in two different ways. First, we can take the original optimal ensemble and treat $`\rho ^{}`$ as the other state ($`\sigma `$ in Eq. 5), obtaining: $$\underset{a}{}p_a𝒮(\rho _a\|\|\rho ^{})=\chi +𝒮(\rho \|\|\rho ^{}).$$ (39) Substituting this expression into the expression for $`\chi ^{}`$ yields: $`\chi ^{}`$ $`=`$ $`\eta 𝒮(\omega \|\|\rho ^{})+(1\eta )(\chi +𝒮(\rho \|\|\rho ^{}))`$ $`\mathrm{\Delta }\chi `$ $`=`$ $`\chi ^{}\chi `$ (40) $`=`$ $`\eta (𝒮(\omega \|\|\rho ^{})\chi )+\eta 𝒮(\rho \|\|\rho ^{})`$ Our second application of Donald’s identity is to the modified ensemble, taking the original average state $`\rho `$ to play the role of the other state: $`\eta 𝒮(\omega \|\|\rho )+(1\eta )\chi `$ $`=`$ $`\chi ^{}+𝒮(\rho ^{}\|\|\rho )`$ (41) $`\mathrm{\Delta }\chi `$ $`=`$ $`\eta (𝒮(\omega \|\|\rho )\chi )𝒮(\rho ^{}\|\|\rho ).`$ (42) Since the relative entropy is never negative, we can conclude that $$\eta (𝒮(\omega \|\|\rho ^{})\chi )\mathrm{\Delta }\chi \eta (𝒮(\omega \|\|\rho )\chi ).$$ (43) This gives upper and lower bounds for the change in $`\chi `$ if we mix in an additional state $`\omega `$ to our original ensemble. The bounds are “tight”, since as $`\eta 0`$, $`𝒮(\omega \|\|\rho ^{})𝒮(\omega \|\|\rho )`$. Very similar bounds for $`\mathrm{\Delta }\chi `$ apply if we make more elaborate modifications of our original ensemble, involving more than one additional signal state. This is described in . We say that an ensemble has the maximal distance property if and only if, for any $`\omega `$ in $`𝒜`$, $$𝒮(\omega \|\|\rho )\chi ,$$ (44) where $`\rho `$ is the average state and $`\chi `$ is the Holevo quantity for the ensemble. This property gives an interesting characterization of optimal ensembles: > Theorem: An ensemble is optimal if and only if it has the maximum distance property. We give the essential ideas of the proof here; further details can be found in . Suppose our ensemble has the maximum distance property. Then, if we add the state $`\omega `$ with probability $`\eta `$, the change $`\mathrm{\Delta }\chi `$ satisfies $$\mathrm{\Delta }\chi \eta (𝒮(\omega \|\|\rho )\chi )0.$$ (45) In other words, we cannot increase $`\chi `$ by mixing in an additional state. Consideration of more general changes to the ensemble leads to the same conclusion that $`\mathrm{\Delta }\chi 0`$. Thus, the ensemble must be optimal, and $`\chi =\chi ^{}`$. Conversely, suppose that the ensemble is optimal (with $`\chi =\chi ^{}`$). Could there be a state $`\omega `$ in $`𝒜`$ such that $`𝒮(\omega \|\|\rho )>\chi ^{}`$? If there were such an $`\omega `$, then by choosing $`\eta `$ small enough we could make $`𝒮(\omega \|\|\rho ^{})>\chi ^{}`$, and so $$\mathrm{\Delta }\chi \eta (𝒮(\omega \|\|\rho ^{})\chi ^{})>0.$$ (46) But this contradicts the fact that, if the original ensemble is optimal, $`\mathrm{\Delta }\chi 0`$ for any change in the ensemble. Thus, no such $`\omega `$ exists and the optimal ensemble satisfies the maximal distance property. Two corollaries follow immediately from this theorem. First, we note that the support of the average state $`\rho `$ of an optimal ensemble must contain the support of every state $`\omega `$ in $`𝒜`$. Otherwise, the relative entropy $`𝒮(\omega \|\|\rho )=\mathrm{}`$, contradicting the maximal distance property. The fact that $`\rho `$ has the largest support possible could be called the maximal support property of an optimal ensemble. Second, we recall that $`\chi ^{}`$ is just the average relative entropy distance of the members of the optimal ensemble from the average state $`\rho `$: $$\chi ^{}=\underset{a}{}p_a𝒮(\rho _a\|\|\rho ).$$ Since $`𝒮(\rho _a\|\|\rho )\chi ^{}`$ for each $`a`$, it follows that whenever $`p_a>0`$ we must have $$𝒮(\rho _a\|\|\rho )=\chi ^{}.$$ (47) We might call this the equal distance property of an optimal ensemble. We can now give an explicit formula for $`\chi ^{}`$ that does not optimize over ensembles, but only over states in $`𝒜`$. From Equation 14, for any state $`\sigma `$, $$\chi \underset{a}{}p_a𝒮(\rho _a\|\|\sigma )$$ (48) and thus $$\chi \underset{\omega }{\mathrm{max}}𝒮(\omega \|\|\sigma )$$ (49) where the maximum is taken over all $`\omega `$ in $`𝒜`$. We apply this inequality to the optimal ensemble, finding the lowest such upper bound for $`\chi ^{}`$: $$\chi ^{}\underset{\sigma }{\mathrm{min}}\left(\underset{\omega }{\mathrm{max}}𝒮(\omega \|\|\sigma )\right).$$ (50) But since the optimal ensemble has the maximal distance property, we know that $$\chi ^{}=\underset{\omega }{\mathrm{max}}𝒮(\omega \|\|\rho )$$ (51) for the optimal average state $`\rho `$. Therefore, $$\chi ^{}=\underset{\sigma }{\mathrm{min}}\left(\underset{\omega }{\mathrm{max}}𝒮(\omega \|\|\sigma )\right).$$ (52) ## 5 Additivity for quantum channels The quantity $`\chi ^{}`$ is an asymptotically achievable upper bound to the amount of classical information that can be sent using available states of the channel system $`Q`$. It is therefore tempting to identify $`\chi ^{}`$ as the classical capacity of the quantum channel. But there is a subtlety here, which involves an important unsolved problem of quantum information theory. Specifically, suppose that two quantum systems $`A`$ and $`B`$ are available for use as communication channels. The two systems evolve independently according the product map $`^A^B`$. Each system can be considered as a separate channel, or the joint system $`AB`$ can be analyzed as a single channel. It is not known whether the following holds in general: $$\chi ^{AB}\stackrel{\mathrm{?}}{=}\chi ^A+\chi ^B.$$ (53) Since separate signal ensembles for $`A`$ and $`B`$ can be combined into a product ensemble for $`AB`$, it is clear that $`\chi ^{AB}\chi ^A+\chi ^B`$. However, the joint system $`AB`$ also has other possible signal ensembles that use entangled input states and that might perhaps have a Holevo bound for the output states greater than $`\chi ^A+\chi ^B`$. Equation 53 is the “additivity conjecture” for the classical capacity of a quantum channel. If the conjecture is false, then the use of entangled input states would sometimes increase the amount of classical information that can be sent over two or more independent channels. The classical capacity of a channel (which is defined asymptotically, using many instances of the same channel) would thus be greater than $`\chi ^{}`$ for a single instance of a channel. On the other hand, if the conjecture holds, then $`\chi ^{}`$ is the classical capacity of the quantum channel. Numerical calculations to date support the additivity conjecture for a variety of channels. Recent work gives strong evidence that Equation 53 holds for various special cases, including channels described by unital maps. We present here another partial result: $`\chi ^{}`$ is additive for any “half-noisy” channel, that is, a dual channel that is represented by an map of the form $`^A^B`$, where $`^A`$ is the identity map on $`A`$. Suppose the joint system $`AB`$ evolves according to the map $`^A^B`$, and let $`\rho ^A`$ and $`\rho ^B`$ be the average output states of optimal signal ensembles for $`A`$ and $`B`$ individually. We will show that the product ensemble (with average state $`\rho ^A\rho ^B`$) is optimal by showing that this ensemble has the maximal distance property. That is, suppose we have another, possibly entangled input state of $`AB`$ that leads to the output state $`\omega ^{AB}`$. Our aim is to prove that $`𝒮(\omega ^{AB}\|\|\rho ^A\rho ^B)\chi ^A+\chi ^B`$. From the definition of $`𝒮(\|\|)`$ we can show that $`𝒮(\omega ^{AB}\|\|\rho ^A\rho ^B)`$ $`=`$ $`S\left(\omega ^{AB}\right)\text{Tr}\omega ^A\mathrm{log}\rho ^A\text{Tr}\omega ^B\mathrm{log}\rho ^B`$ (54) $`=`$ $`S\left(\omega ^A\right)+S\left(\omega ^B\right)S\left(\omega ^{AB}\right)`$ $`+𝒮(\omega ^A\|\|\rho ^A)+𝒮(\omega ^B\|\|\rho ^B).`$ (The right-hand expression has an interesting structure; $`S(\omega ^A)+S(\omega ^B)S(\omega ^{AB})`$ is clearly analogous to the mutual information defined in Equation 9.) Since $`A`$ evolves according to the identity map $`^A`$, it is easy to see that $`\chi ^A=d=dim^A`$ and $$\rho ^A=\left(\frac{1}{d}\right)\mathrm{\hspace{0.17em}1}^A.$$ (55) From this it follows that $$S\left(\omega ^A\right)+𝒮(\omega ^A\|\|\rho ^A)=\mathrm{log}d=\chi ^A$$ (56) for any $`\omega ^A`$. This accounts for two of the terms on the right-hand side of Equation 54. The remaining three terms require a more involved analysis. The final joint state $`\omega ^{AB}`$ is a mixed state, but we can always introduce a third system $`C`$ that “purifies” the state. That is, we can find $`\|\mathrm{\Omega }^{ABC}`$ such that $$\omega ^{AB}=\text{Tr}_C\|\mathrm{\Omega }^{ABC}\mathrm{\Omega }^{ABC}\|.$$ (57) Since the overall state of $`ABC`$ is a pure state, $`S(\omega ^{AB})=S(\omega ^C)`$, where $`\omega ^C`$ is the state obtained by partial trace over $`A`$ and $`B`$. Furthermore, imagine that a complete measurement is made on $`A`$, with the outcome $`k`$ occuring with probability $`p_k`$. For a given measurement outcome $`k`$, the subsequent state of the remaining system $`BC`$ will be $`\|\mathrm{\Omega }_k^{BC}`$. Letting $`\omega _k^B`$ $`=`$ $`\text{Tr}_C\|\mathrm{\Omega }_k^{BC}\mathrm{\Omega }_k^{BC}\|`$ $`\omega _k^C`$ $`=`$ $`\text{Tr}_B\|\mathrm{\Omega }_k^{BC}\mathrm{\Omega }_k^{BC}\|,`$ (58) we have that $`S(\omega _k^B)=S(\omega _k^C)`$ for all $`k`$. Furthermore, by locality, $`\omega ^B={\displaystyle \underset{k}{}}p_k\omega _k^B`$ $`\omega ^C={\displaystyle \underset{k}{}}p_k\omega _k^C.`$ (59) In other words, we have written both $`\omega ^B`$ and $`\omega ^C`$ as ensembles of states. We can apply this to get an upper bound on the remaining terms in Equation 54 $`S\left(\omega ^B\right)S\left(\omega ^{AB}\right)+𝒮(\omega ^B\|\|\rho ^B)`$ (60) $`=`$ $`S\left(\omega ^B\right){\displaystyle \underset{k}{}}p_kS\left(\omega _k^B\right)`$ $`S\left(\omega ^C\right)+{\displaystyle \underset{k}{}}p_kS\left(\omega _k^C\right)+𝒮(\omega ^B\|\|\rho ^B)`$ $``$ $`\chi _\omega ^B+𝒮(\omega ^B\|\|\rho ^B),`$ where $`\chi _\omega ^B`$ is the Holevo quantity for the ensemble of $`\omega _k^B`$ states. Donald’s identity permits us to write $$S\left(\omega ^B\right)S\left(\omega ^{AB}\right)+𝒮(\omega ^B\|\|\rho ^B)=\underset{k}{}p_k𝒮(\omega _k^B\|\|\rho ^B).$$ (61) The $`B`$ states $`\omega _k^B`$ are all available output states of the $`B`$ channel. These states are obtained by making a complete measurement on system $`A`$ when the joint system $`AB`$ is in the state $`\omega ^{AB}`$. But this state was obtained from some initial $`AB`$ state and a dynamical map $`^A^B`$. This map commutes with the measurement operation on $`A`$ alone, so we could equally well make the measurement before the action of $`^A^B`$. The $`A`$-measurement outcome $`k`$ would then determine the input state of $`B`$, which would evolve into $`\omega _k^B`$. Thus, for each $`k`$, $`\omega _k^B`$ is a possible output of the $`^B`$ map. Since $`\rho ^B`$ has the maximum distance property and the states $`\omega _k^B`$ are available outputs of the channel, $`𝒮(\omega _k^B\|\|\rho ^B)\chi ^B`$ for every $`k`$. Combining Equations 54, 56 and 61, we find the desired inequality: $$𝒮(\omega ^{AB}\|\|\rho ^A\rho ^B)\chi ^A+\chi ^B.$$ (62) This demonstrates that the product of optimal ensembles for $`A`$ and $`B`$ also has the maximum distance property for the possible outputs of the joint channel, and so this product ensemble must be optimal. It follows that $`\chi ^{AB}=\chi ^A+\chi ^B`$ in this case. Our result has been phrased for the case in which $`A`$ undergoes “trivial” dynamics $`^A`$, but the proof also works without modification if the time evolution of $`A`$ is unitary—that is, $`A`$ experiences “distortion” but not “noise”. If only one of the two systems is noisy, then $`\chi ^{}`$ is additive. The additivity conjecture for $`\chi ^{}`$ is closely related to another additivity conjecture, the “minimum output entropy” conjecture . Suppose $`A`$ and $`B`$ are systems with independent evolution described by $`^A^B`$, and let $`\rho ^AB`$ be an output state of the channel with minimal entropy $`S(\rho ^AB)`$. Is $`\rho ^{AB}`$ a product state $`\rho ^A\rho ^B`$? The answer is not known in general; but it is quite easy to show this in the half-noisy case that we consider here. ## 6 Maximizing coherent information When we turn from the transmission of classical information to the transmission of quantum information, it will be helpful to adopt an explicit description of the channel dynamics, instead of merely specifying the set of available output states $`𝒜`$. Suppose the quantum system $`Q`$ undergoes a dynamical evolution described by the map $``$. Since $``$ is a trace-preserving, completely positive map, we can always find a representation of $``$ as a unitary evolution of a larger system . In this representation, we imagine that an additonal “environment” system $`E`$ is present, initially in a pure state $`\|\stackrel{˘}{0}^E`$, and that $`Q`$ and $`E`$ interact via the unitary evolution operator $`U^{QE}`$. That is, $$\rho ^Q=(\stackrel{˘}{\rho }^Q)=\text{Tr}_EU^{QE}\left(\stackrel{˘}{\rho }^Q\|\stackrel{˘}{0}^E\stackrel{˘}{0}^E\|\right)U^{QE}.$$ (63) For convenience, we denote an initial state of a system by the breve accent (as in $`\stackrel{˘}{\rho }^Q`$), and omit this symbol for final states. The problem of sending quantum information through our channel can be viewed in one of two ways: 1. An unknown pure quantum state of $`Q`$ is to be transmitted. In this case, our criterion of success is the average fidelity $`\overline{F}`$, defined as follows. Suppose the input state $`\|\stackrel{˘}{\varphi }_k`$ occurs with probability $`p_k`$ and leads to the output state $`\rho _k`$. Then $$\overline{F}=\underset{k}{}p_k\stackrel{˘}{\varphi }_k\left\|\rho _k\right\|\stackrel{˘}{\varphi }_k.$$ (64) In general, $`\overline{F}`$ depends not only on the average input state $`\stackrel{˘}{\rho }^Q`$ but also on the particular pure state input ensemble. 2. A second “bystander” system $`R`$ is present, and the joint system $`RQ`$ is initially in a pure entangled state $`\|\stackrel{˘}{\mathrm{\Psi }}^{RQ}`$. The system $`R`$ has “trivial” dynamics described by the identity map $``$, so that the joint system evolves according to $``$, yielding a final state $`\rho ^{RQ}`$. Success is determined in this case by the entanglement fidelity $`F_e`$, defined by $$F_e=\stackrel{˘}{\mathrm{\Psi }}^{RQ}\left\|\rho ^{RQ}\right\|\stackrel{˘}{\mathrm{\Psi }}^{RQ}.$$ (65) It turns out, surprisingly, that $`F_e`$ is only dependent on $``$ and the input state $`\stackrel{˘}{\rho }^Q`$ of $`Q`$ alone. That is, $`F_e`$ is an “intrinsic” property of $`Q`$ and its dynamics. These two pictures of quantum information transfer are essentially equivalent, since $`F_e`$ approaches unity if and only if $`\overline{F}`$ approaches unity for every ensemble with the same average input state $`\stackrel{˘}{\rho }^Q`$. For now we adopt the second point of view, in which the transfer of quantum information is essentially the transfer of quantum entanglement (with the bystander system $`R`$) through the channel. The quantum capacity of a channel should be defined as the amount of entanglement that can be transmitted through the channel with $`F_e1`$, if we allow ourselves to use the channel many times and employ quantum error correction schemes . At present it is not known how to calculate this asymptotic capacity of the channel in terms of the properties of a single instance of the channel. Nevertheless, we can identify some quantities that are useful in describing the quantum information conveyed by the channel . A key quantity is the coherent information $`I^Q`$, defined by $$I^Q=S\left(\rho ^Q\right)S\left(\rho ^{RQ}\right).$$ (66) This quantity is a measure of the final entanglement between $`R`$ and $`Q`$. (The initial entanglement is measured by the entropy $`S(\stackrel{˘}{\rho }^Q)`$ of the initial state of $`Q`$, which of course equals $`S(\stackrel{˘}{\rho }^R)`$. See Section 7 below.) If we adopt a unitary representation for $``$, then the overall system $`RQE`$ including the environment remains in a pure state from beginning to end, and so $`S(\rho ^{RQ})=S(\rho ^E)`$. Thus, $$I^Q=S\left(\rho ^Q\right)S\left(\rho ^E\right).$$ (67) Despite the apparent dependence of $`I^Q`$ on the systems $`R`$ and $`E`$, it is in fact a function only of the map $``$ and the initial state $`\stackrel{˘}{\rho }^Q`$ of $`Q`$. Like the entanglement fidelity $`F_e`$, it is an “intrinsic” characteristic of the channel system $`Q`$ and its dynamics. It can be shown that the coherent information $`I^Q`$ does not increase if the map $``$ is followed by a second independent map $`^{}`$, giving an overall dynamics described by $`^{}`$. That is, the coherent information cannot be increased by any “quantum data processing” on the channel outputs. The coherent information is also closely related to quantum error correction. Perfect quantum error correction—resulting in $`F_e=1`$ for the final state—is possible if and only if the channel loses no coherent information, so that $`I^Q=S(\stackrel{˘}{\rho }^Q)`$. These and other properties lead us to consider $`I^Q`$ as a good measure of the quantum information that is transmitted through the channel . The coherent information has an intriguing relation to the Holevo quantity $`\chi `$, and thus to classical information transfer (and to relative entropy) . Suppose we describe that the input state $`\stackrel{˘}{\rho }^Q`$ by an ensemble of pure states $`\|\stackrel{˘}{\varphi }_k^Q`$: $$\stackrel{˘}{\rho }^Q=\underset{k}{}p_k\|\stackrel{˘}{\varphi }_k^Q\stackrel{˘}{\varphi }_k^Q\|.$$ (68) We adopt a unitary representation for the evolution and note that the initial pure state $`\|\stackrel{˘}{\varphi }_k^Q\|\stackrel{˘}{0}^E`$ evolves into a pure, possibly entangled state $`\|\varphi _k^{QE}`$. Thus, for each $`k`$ the entropies of the final states of $`Q`$ and $`E`$ are equal: $$S\left(\rho _k^Q\right)=S\left(\rho _k^E\right).$$ (69) It follows that $`I^Q`$ $`=`$ $`S\left(\rho ^Q\right)S\left(\rho ^E\right)`$ $`=`$ $`S\left(\rho ^Q\right){\displaystyle \underset{k}{}}p_kS\left(\rho _k^Q\right)S\left(\rho ^E\right)+{\displaystyle \underset{k}{}}p_kS\left(\rho _k^E\right)`$ $`I^Q`$ $`=`$ $`\chi ^Q\chi ^E.`$ (70) Remarkably, the difference $`\chi ^Q\chi ^E`$ depends only on $``$ and the average input state $`\stackrel{˘}{\rho }^Q`$, not the details of the environment $`E`$ or the exact choice of pure state input ensemble. The quantities $`\chi ^Q`$ and $`\chi ^E`$ are related to the classical information transfer to the output system $`Q`$ and to the environment $`E`$, respectively. Thus, Equation 70 relates the classical and quantum information properties of the channel. This relation has been used to analyze the privacy of quantum cryptographic channels . We will use it here to give a relative entropy characterization of the the input state $`\stackrel{˘}{\rho }^Q`$ that maximizes the coherent information of the channel. Let us suppose that $`\stackrel{˘}{\rho }^Q`$ is an input state that maximizes the coherent information $`I^Q`$. If we change the input state to $$\stackrel{˘}{\rho }^Q{}_{}{}^{}=(1\eta )\stackrel{˘}{\rho }^Q+\eta \stackrel{˘}{\omega }^Q,$$ (71) for some pure state $`\stackrel{˘}{\omega }^Q`$, we produces some change $`\mathrm{\Delta }I^Q`$ in the coherent information. Viewing $`\stackrel{˘}{\rho }^Q`$ as an ensemble of pure states, this change amounts to a modification of that ensemble; and such a modification leads to changes in the output ensembles for both system $`Q`$ and system $`E`$. Thus, $$\mathrm{\Delta }I^Q=\mathrm{\Delta }\chi ^Q\mathrm{\Delta }\chi ^E.$$ (72) We can apply Equation 43 to bound both $`\mathrm{\Delta }\chi ^Q`$ and $`\mathrm{\Delta }\chi ^E`$ and obtain a lower bound for $`\mathrm{\Delta }I^Q`$: $`\mathrm{\Delta }I^Q`$ $``$ $`\eta (𝒮(\omega ^Q\|\|\rho _{}^{Q}{}_{}{}^{})\chi ^Q)\eta (𝒮(\omega ^E\|\|\rho ^E)\chi ^E)`$ $`\mathrm{\Delta }I^Q`$ $``$ $`\eta (𝒮(\omega ^Q\|\|\rho _{}^{Q}{}_{}{}^{})𝒮(\omega ^E\|\|\rho ^E)I^Q).`$ (73) Since we assume that $`I^Q`$ is maximized for the input $`\stackrel{˘}{\rho }^Q`$, then $`\mathrm{\Delta }I^Q0`$ when we modify the input state. This must be true for every value of $`\eta `$ in the relation above. Whenever $`𝒮(\omega ^Q\|\|\rho ^Q)`$ is finite, we can conclude that $$𝒮(\omega ^Q\|\|\rho ^Q)𝒮(\omega ^E\|\|\rho ^E)I^Q.$$ (74) This is analogous to the maximum distance property for optimal signal ensembles, except that it is the difference of two relative entropy distances that is bounded above by the maximum of $`I^Q`$. Let us write Equation 70 in terms of relative entropy, imagining that the input state $`\stackrel{˘}{\rho }^Q`$ is written in terms of an ensemble of pure states $`\|\stackrel{˘}{\varphi }_k^Q`$: $$I^Q=\underset{k}{}p_k(𝒮(\rho _k^Q\|\|\rho ^Q)𝒮(\rho _k^E\|\|\rho ^E)).$$ (75) Every input pure state $`\|\stackrel{˘}{\varphi }_k^Q`$ in the input ensemble with $`p_k>0`$ will be in the support of $`\stackrel{˘}{\rho }^Q`$, and so Equation 74 holds. Therefore, we can conclude that $$I^Q=𝒮(\rho _k^Q\|\|\rho ^Q)𝒮(\rho _k^E\|\|\rho ^E)$$ (76) for every such state in the ensemble. Furthermore, any pure state in the support of $`\stackrel{˘}{\rho }^Q`$ is a member of some pure state ensemble for $`\stackrel{˘}{\rho }^Q`$. This permits us to draw a remarkable conclusion. If $`\stackrel{˘}{\rho }^Q`$ is the input state that maximizes the coherent information $`I^Q`$ of the channel, then for any pure state $`\stackrel{˘}{\omega }^Q`$ in the support of $`\stackrel{˘}{\rho }^Q`$, $$I^Q=𝒮(\omega ^Q\|\|\rho ^Q)𝒮(\omega ^E\|\|\rho ^E).$$ (77) This result is roughly analogous to the equal distance property for optimal signal ensembles. Together with Equation 74, it provides a strong characterization of the state that maximizes coherent information. The additivity problem for $`\chi ^{}`$ leads us to ask whether the maximum of the coherent information is additive when independent channels are combined. In fact, there are examples known where $`\mathrm{max}I^{AB}>\mathrm{max}I^A+\mathrm{max}I^B`$; in other words, entanglement between independent channels can increase the amount of coherent information that can be sent through them . The asymptotic behavior of coherent information and its precise connection to quantum channel capacities are questions yet to be resolved. ## 7 Indeterminate length quantum coding In the previous section we saw that the relative entropy can be used to analyze the coherent information “capacity” of a quantum channel. Another issue in quantum information theory is quantum data compression , which seeks to represent quantum information using the fewest number of qubits. In this section we will see that the relative entropy describes the cost of suboptimal quantum data compression. One approach to classical data compression is to use variable length codes, in which the codewords are finite binary strings of various lengths . The best-known examples are the Huffman codes. The Shannon entropy $`H(X)`$ of a random variable $`X`$ is a lower bound to the average codeword length in such codes, and for Huffman codes this average codeword length can be made arbitrarily close to $`H(X)`$. Thus, a Huffman code optimizes the use of a communication resources (number of bits required) in classical communication without noise. There are analogous codes for the compression of quantum information. Since coherent superpositions of codewords must be allowed as codewords, these are called indeterminate length quantum codes . A quantum analogue to Huffman coding was recently described by Braunstein et al. An account of the theory of indeterminate length quantum codes, including the quantum Kraft inequality and the condensability condition (see below), will be presented in a forthcoming paper . Here we will outline a few results and demonstrate a connection to the relative entropy. The key idea in constructing an indeterminate length code is that the codewords themselves must carry their own length information. For a classical variable length code, this requirement can be phrased in two ways. A uniquely decipherable code is one in which any string of $`N`$ codewords can be correctly separated into its individual codewords, while a prefix-free code is one in which no codeword is an initial segment of another codeword. The lengths of the codewords in each case satisfy the Kraft-McMillan inequality: $$\underset{k}{}2^{l_k}1,$$ (78) where is the sum is over the codewords and $`l_k`$ is the length of the $`k`$th codeword. Every prefix-free code is uniquely decipherable, so the prefix-free property is a more restrictive property. Nevertheless, it turns out that any uniquely decipherable code can be replaced by a prefix-free code with the same codeword lengths. There are analogous conditions for indeterminate length quantum codes, but these properties must be phrased carefully because we allow coherent superpositions of codewords. For example, a classical prefix-free code is sometimes called an “instantaneous” code, since as soon as a complete codeword arrives we can recognize it at once and decipher it immediately. However, if an “instantaneous” decoding procedure were to be attempted for a quantum prefix-free code, it would destroy coherences between codewords of different lengths. Quantum codes require that the entire string of codewords be deciphered together. The property of an indeterminate length quantum code that is analogous to unique decipherability is called condensability. We digress briefly to describe the condensability condition. We focus on zero-extended forms (zef ) of our codewords. That is, we cosider that our codewords occupy an initial segment of a qubit register of fixed length $`n`$, with $`\|0`$’s following. (Clearly $`n`$ must be chosen large enough to contain the longest codeword.) The set of all zef codewords spans a subspace of the Hilbert space of register states. We imagine that the output of a quantum information source has been mapped unitarily to the zef codeword space of the register. Our challenge is to take $`N`$ such registers and “pack” them together in a way that can exploit the fact that some of the codewords are shorter than others. If codeword states must carry their own length information, there must be a length observable $`\mathrm{\Lambda }`$ on the zef codeword space with the following two properties: * The eigenvalues of $`\mathrm{\Lambda }`$ are integers $`1,\mathrm{},n`$, where $`n`$ is the length of the register. * If $`\|\psi _{\text{zef}}`$ is an eigenstate of $`\mathrm{\Lambda }`$ with eigenvalue $`l`$, then it has the form $$\|\psi _{\text{zef}}=\|\psi ^{1\mathrm{}l}0^{l+1\mathrm{}n}.$$ (79) That is, the last $`nl`$ qubits in the register are in the state $`\|0`$ for a zef codeword of length $`l`$. For register states not in the zef subspace, we can take $`\mathrm{\Lambda }=\mathrm{}`$. A code is condensable if the following condition holds: For any $`N`$, there is a unitary operator $`U`$ (depending on $`N`$) that maps $$\underset{Nn\text{qubits}}{\underset{}{\|\psi _{1,\text{zef}}\mathrm{}\|\psi _{N,\text{zef}}}}\underset{Nn\text{qubits}}{\underset{}{\|\mathrm{\Psi }_{1\mathrm{}N}}}$$ with the property that, if the individual codewords are all length eigenstates, then $`U`$ maps the codewords to a zef string of the $`Nn`$ qubits—that is, one with $`\|0`$’s after the first $`L=l_1+\mathrm{}+l_N`$ qubits: $$\|\psi _1^{1\mathrm{}l_1}0^{l_1+1\mathrm{}n}\mathrm{}\|\psi _N^{1\mathrm{}l_N}0^{l_N+1\mathrm{}n}\|\mathrm{\Psi }^{1\mathrm{}L}0^{L+1\mathrm{}Nn}.$$ The unitary operator $`U`$ thus “packs” $`N`$ codewords, given in their zef forms, into a “condensed” string that has all of the trailing $`\|0`$’s at the end. The unitary character of the packing protocol automatically yields an “unpacking” procedure given by $`U^1`$. Thus, if the quantum code is condensable, a packed string of $`N`$ codewords can be coherently sorted out into separated zef codewords. The quantum analogue of the Kraft-McMillan inequality states that, for any indeterminate length quantum code that is condensable, the length observable $`\mathrm{\Lambda }`$ on the subspace of zef codewords must satisfy $$\text{Tr}\mathrm{\hspace{0.17em}2}^\mathrm{\Lambda }1,$$ (80) where we have restricted our trace to the zef subspace. We can construct a density operator $`\omega `$ (a positive operator of unit trace) on the zef subspace by letting $`K=\text{Tr}\mathrm{\hspace{0.17em}2}^\mathrm{\Lambda }1`$ and $$\omega =\frac{1}{K}\mathrm{\hspace{0.17em}2}^\mathrm{\Lambda }.$$ (81) The density operator $`\omega `$ is generally not the same as the actual density operator $`\rho `$ of the zef codewords produced by the quantum information source. The average codeword length is $`\overline{l}`$ $`=`$ $`\text{Tr}\rho \mathrm{\Lambda }`$ $`=`$ $`\text{Tr}\rho \mathrm{log}\left(2^\mathrm{\Lambda }\right)`$ $`=`$ $`\text{Tr}\rho \mathrm{log}\omega \mathrm{log}K`$ $`\overline{l}`$ $`=`$ $`S(\rho )+𝒮(\rho \|\|\omega )\mathrm{log}K.`$ (82) Since $`\mathrm{log}K0`$ and the relative entropy is positive definite, $$\overline{l}S(\rho ).$$ (83) The average codeword length must always be at least as great as the von Neuman entropy of the information source. Equality for Equation 83 can be approached asymptotically using block coding and a quantum analogue of Huffman (or Shannon-Fano) coding. For special cases in which the eigenvalues of $`\rho `$ are of the form $`2^m`$, then a code exists for which $`\overline{l}=S(\rho )`$, without the asymptotic limit. In either case, we say that a code satisfying $`\overline{l}=S(\rho )`$ is a length optimizing quantum code. Equation 82 tells us that, if we have a length optimizing code, $`K=1`$ and $$\rho =\omega =2^\mathrm{\Lambda }.$$ (84) The condensed string of $`N`$ codewords has $`Nn`$ qubits, but we can discard all but about $`N\overline{l}`$ of them and still retain high fidelity. That is, $`\overline{l}`$ is the asymptotic number of qubits that must be used per codeword to represent the quantum information faithfully. Suppose that we have an indeterminate length quantum code that is designed for the wrong density operator. That is, our code is length optimizing for some other density operator $`\omega `$, but $`\rho \omega `$. Then (recalling that $`K=1`$ for a length optimizing code, even if it is optimizing for the wrong density operator), $$\overline{l}=S(\rho )+𝒮(\rho \|\|\omega ).$$ (85) $`S(\rho )`$ tells us the number of qubits necessary to represent the quantum information if we used a length optimizing code for $`\rho `$. (As we have mentioned, such codes always exist in an asymptotic sense.) However, to achieve high fidelity in the situation where we have used a code designed for $`\omega `$, we have to use at least $`\overline{l}`$ qubits per codeword, an additional cost of $`𝒮(\rho \|\|\omega )`$ qubits per codeword. This result gives us an interpretation of the relative entropy function $`𝒮(\rho \|\|\omega )`$ in terms of the physical resources necessary to accomplish some task—in this case, the additional cost (in qubits) of representing the quantum information described by $`\rho `$ using a coding scheme optimized for $`\omega `$. This is entirely analogous to the situation for classical codes and classical relative entropy . A fuller development of this analysis will appear in . ## 8 Relative entropy of entanglement One recent application of relative entropy has been to quantify the entanglement of a mixed quantum state of two systems . Suppose Alice and Bob share a joint quantum system $`AB`$ in the state $`\rho ^{AB}`$. This state is said to be separable if it is a product state or else a probabilistic combination of product states: $$\rho ^{AB}=\underset{k}{}p_k\rho _k^A\rho _k^B.$$ (86) Without loss of generality, we can if we wish take the elements in this ensemble of product states to be pure product states. Systems in separable states display statistical correlations having perfectly ordinary “classical” properties—that is, they do not violate any sort of Bell inequality. A separable state of $`A`$ and $`B`$ could also be created from scratch by Alice and Bob using only local quantum operations (on $`A`$ and $`B`$ separately) and the exchange of classical information. States which are not separable are said to be entangled. These states cannot be made by local operations and classical communication; in other words, their creation requires the exchange of quantum information between Alice and Bob. The characterization of entangled states and their possible transformations has been a central issue in much recent work on quantum information theory. A key question is the quantification of entanglement, that is, finding numerical measures of the entanglement of a quantum state $`\rho ^{AB}`$ that have useful properties. If the joint system $`AB`$ is in a pure state $`\|\mathrm{\Psi }^{AB}`$, so that the subsystem states are $$\begin{array}{c}\rho ^A=\text{Tr}_B\|\mathrm{\Psi }^{AB}\mathrm{\Psi }^{AB}\|\\ \rho ^B=\text{Tr}_A\|\mathrm{\Psi }^{AB}\mathrm{\Psi }^{AB}\|\end{array}$$ (87) then the entropy $`S(\rho ^A)=S(\rho ^B)`$ can be used to measure the entanglement of $`A`$ and $`B`$. This measure has many appealing properties. It is zero if and only if $`\|\mathrm{\Psi }^{AB}`$ is separable (and thus a product state). For an “EPR pair” of qubits—that is, a state of the general form $$\|\varphi ^{AB}=\frac{1}{\sqrt{2}}\left(\|0^A0^B+\|1^A1^B\right),$$ (88) the susbsystem entropy $`S(\rho ^A)=1`$ bit. The subsystem entropy is also an asymptotic measure, both of the resources necessary to create the particular entangled pure state, and of the value of the state as a resource . That is, for sufficiently large $`N`$, * approximately $`NS(\rho ^A)`$ EPR pairs are required to create $`N`$ copies of $`\|\mathrm{\Psi }^{AB}`$ by local operations and classical communication; and * approximately $`NS(\rho ^A)`$ EPR pairs can be created from $`N`$ copies of $`\|\mathrm{\Psi }^{AB}`$ by local operations and classical communication. For mixed entangled states $`\rho ^{AB}`$ of the joint system $`AB`$, things are not so well-established. Several different measures of entanglement are known, including * the entanglement of formation $`E(\rho ^{AB})`$, which is the minimum asymptotic number of EPR pairs required to create $`\rho ^{AB}`$ by local operations and classical communication; and * the distillable entanglement $`D(\rho ^{AB})`$, the maximum asymptotic number of EPR pairs that can be created from $`\rho ^{AB}`$ by entanglement purification protocols involving local operations and classical communication. Bennett et al. further distinguish $`D_1`$ and $`D_2`$, the distillable entanglements with respect to purification protocols that allow one-way and two-way classical communication, respectively. All of these measures reduce to the subsystem entropy $`S(\rho ^A)`$ if $`\rho ^{AB}`$ is a pure entangled state. These entanglement measures are not all equal; furthermore, explicit formulas for their calculation are not known in most cases. This motivates us to consider alternate measures of entanglement with more tractable properties and which have useful relations to the asymptotic measures $`E`$, $`D_1`$ and $`D_2`$. A state $`\rho ^{AB}`$ is entangled inasmuch as it is not a separable state, so it makes sense to adopt as a measure of entanglement a measure of the distance of $`\rho ^{AB}`$ from the set $`\mathrm{\Sigma }^{AB}`$ of separable states of $`AB`$. Using relative entropy as our “distance”, we define the relative entropy of entanglement $`E_r`$ to be $$E_r\left(\rho ^{AB}\right)=\underset{\sigma ^{AB}\mathrm{\Sigma }^{AB}}{\mathrm{min}}𝒮(\rho ^{AB}\|\|\sigma ^{AB}).$$ (89) The relative entropy of entanglement has several handy properties. First of all, it reduces to the subsystem entropy $`S(\rho ^A)`$ whenever $`\rho ^{AB}`$ is a pure state. Second, suppose we write $`\rho ^{AB}`$ as an ensemble of pure states $`\|\psi _k^{AB}`$. Then $$E_r\left(\rho ^{AB}\right)\underset{k}{}p_kS\left(\rho _k^A\right)$$ (90) where $`\rho _k^A=\text{Tr}_B\|\psi _k^{AB}\psi _k^{AB}\|`$. It follows from this that $`E_rE`$ for any state $`\rho ^{AB}`$. Even more importantly, the relative entropy of entanglement $`E_r`$ can be shown to be non-increasing on average under local operations by Alice and Bob together with classical communication between them. The quantum version of Sanov’s theorem gives the relative entropy of entanglement an interpretation in terms of the statistical distinguishability of $`\rho ^{AB}`$ and the “least distinguishable” separable state $`\sigma ^{AB}`$. The relative entropy of entanglement is thus a useful and well-motivated measure of the entanglement of a state $`\rho ^{AB}`$ of a joint system, both on its own terms and as a surrogate for less tractable asymptotic measures. ## 9 Manipulating multiparticle entanglement The analysis in this section closely follows that of Linden et al. , who provides a more detailed discussion of the main result here and its applications. Suppose Alice, Bob and Claire initially share three qubits in a “GHZ state” $$\|\mathrm{\Psi }^{ABC}=\frac{1}{\sqrt{2}}\left(\|0^A0^B0^C+\|1^A1^B0^C\right).$$ (91) The mixed state $`\rho ^{BC}`$ shared by Bob and Claire is, in fact, not entangled at all: $$\rho ^{BC}=\frac{1}{2}\left(\|0^B0^C0^B0^C\|+\|1^B1^C1^B1^C\|\right).$$ (92) No local operations performed by Bob and Claire can produce an entangled state from this starting point. However, Alice can create entanglement for Bob and Claire. Alice measures her qubit in the basis $`\{\|+^A,\|^A\}`$, where $$\|\pm ^A=\frac{1}{\sqrt{2}}\left(\|0^A\pm \|1^A\right).$$ (93) It is easy to verify that the state of Bob and Claire’s qubits after this measurement, depending on the measurement outcome, must be one of the two states $$\|\varphi _\pm ^{BC}=\frac{1}{\sqrt{2}}\left(\|0^A0^B\pm \|1^A1^B\right),$$ (94) both of which are equivalent (up to a local unitary transformation by either Bob or Claire) to an EPR pair. In other words, if Alice makes a local measurement on her qubit and then announces the result by classical communication, the GHZ triple can be converted into an EPR pair for Bob and Claire. When considering the manipulation of quantum entanglement shared among several parties, we must therefore bear in mind that the entanglement between subsystems can both increase and decrease, depending on the situation. This raises several questions: Under what circumstances can Alice increase Bob and Claire’s entanglement? How much can she do so? Are there any costs involved in the process? To study these questions, we must give a more detailed account of “local operations and classical communication”. It turns out that Alice, Bob and Claire can realize any local operation on their joint system $`ABC`$ by a combination of the following: * Local unitary transformations on the subsystems $`A`$, $`B`$ and $`C`$; * Adjoining to a subsystem additional local “ancilla” qubits in a standard state $`\|0`$; * Local ideal measurements on the (augmented) subsystems $`A`$, $`B`$ and $`C`$; and * Discarding local ancilla qubits. Strictly speaking, though, we do not need to include the last item. That is, any protocol that involves discarding ancilla qubits can be replaced by one in which the ancillas are simply “set aside”—not used in future steps, but not actually gotten rid of. In a similar vein, we can imagine that the ancilla qubits required are already present in the subsystems $`A`$, $`B`$ and $`C`$, so the second item in our list is redundant. We therefore need to consider only local unitary transformations and local ideal measurements. What does classical communication add to this? It is sufficient to suppose that Alice, Bob and Claire have complete information—that is, they are aware of all operations and the outcomes of all measurements performed by each of them, and thus know the global state of $`ABC`$ at every stage. Any protocol that involved an incomplete sharing of information could be replaced by one with complete sharing, simply by ignoring some of the messages that are exchanged. Our local operations (local unitary transformations and local ideal measurements) always take an initial pure state to a final pure state. That is, if $`ABC`$ starts in the joint state $`\|\mathrm{\Psi }^{ABC}`$, then the final state will be a pure state $`\|\mathrm{\Psi }_k^{ABC}`$ that depends on the joint outcome $`k`$ of all the measurements performed. Thus, $`ABC`$ is always in a pure state known to all parties. It is instructive to consider the effect of local operations on the entropies of the various subsystems of $`ABC`$. Local unitary transformations leave $`S(\rho ^A)`$, $`S(\rho ^B)`$ and $`S(\rho ^C)`$ unchanged. But suppose that Alice makes an ideal measurement on her subsystem, obtaining outcome $`k`$ with probability $`p_k`$. The initial global state is $`\|\mathrm{\Psi }^{ABC}`$ and the final global state is $`\|\mathrm{\Psi }_k^{ABC}`$, depending on $`k`$. For the initial subsystem states, we have that $$S\left(\rho ^A\right)=S\left(\rho ^{BC}\right)$$ (95) since the overall state is a pure state. Similarly, the various final subsystem states satisfy $$S\left(\rho _k^A\right)=S\left(\rho _k^{BC}\right).$$ (96) But an operation on $`A`$ cannot change the average state of $`BC`$: $$\rho ^{BC}=\underset{k}{}p_k\rho _k^{BC}.$$ (97) Concavity of the entropy gives $$S\left(\rho ^{BC}\right)\underset{k}{}p_kS\left(\rho _k^{BC}\right)$$ (98) and therefore $$S\left(\rho ^A\right)\underset{k}{}p_kS\left(\rho _k^A\right).$$ (99) Concavity also tells us that $`S(\rho ^B){\displaystyle \underset{k}{}}p_kS(\rho _k^B)`$, etc., and similar results hold for local measurements performed by Bob or Claire. We now return to the question of how much Alice can increase the entanglement shared by Bob and Claire. Let us measure the bipartite entanglement of the system $`BC`$ (which may be in a mixed state) by the relative entropy of entanglement $`E_r(\rho ^{BC})`$, and let $`\sigma ^{BC}`$ be the separable state of $`BC`$ for which $$E_r(\rho ^{BC})=𝒮(\rho ^{BC}\|\|\sigma ^{BC}).$$ (100) No local unitary operation can change $`E_r(\rho ^{BC})`$; furthermore, no local measurement by Bob or Claire can increase $`E_r(\rho ^{BC})`$ on average. We need only consider an ideal measurement performed by Alice on system $`A`$. Once again we suppose that outcome $`k`$ of this measurement occurs with probability $`p_k`$, and once again Equation 97 holds. Donald’s identity tells us that $$\underset{k}{}p_k𝒮(\rho _k^{BC}\|\|\sigma ^{BC})=\underset{k}{}p_k𝒮(\rho _k^{BC}\|\|\rho ^{BC})+𝒮(\rho ^{BC}\|\|\sigma ^{BC}).$$ (101) But $`E_r(\rho _k^{BC})𝒮(\rho _k^{BC}\|\|\sigma ^{BC})`$ for every $`k`$, leading to the following inequality: $$\underset{k}{}p_kE_r(\rho _k^{BC})E_r(\rho ^{BC})\underset{k}{}p_k𝒮(\rho _k^{BC}\|\|\rho ^{BC}).$$ (102) We recognize the left-hand side of this inequality $`\chi `$ for the ensemble of post-measurement states of $`BC`$, which we can rewrite using the definition of $`\chi `$ in Equation 11. This yields: $`{\displaystyle \underset{k}{}}p_kE_r(\rho _k^{BC})E_r(\rho ^{BC})`$ $``$ $`S\left(\rho ^{BC}\right){\displaystyle \underset{k}{}}p_kS\left(\rho _k^{BC}\right)`$ (103) $`=`$ $`S\left(\rho ^A\right){\displaystyle \underset{k}{}}p_kS\left(\rho _k^A\right),`$ since the overall state of $`ABC`$ is pure at every stage. To summarize, in our model (in which all measurements are ideal, all classical information is shared, and no classical or quantum information is ever discarded), the following principles hold: * The entropy of any subsystem $`A`$ cannot be increased on average by any local operations. * The relative entropy of entanglement of two subsystems $`B`$ and $`C`$ cannot be increased on average by local operations on those two subsystems. * The relative entropy of entanglement of $`B`$ and $`C`$ can be increased by a measurement performed on a third subsystem $`A`$, but the average increase in $`E_r^BC`$ is no larger than the average decrease in the entropy of $`A`$. We say that a joint state $`\|\mathrm{\Psi }_1^{ABC}`$ can be transformed reversibly into $`\|\mathrm{\Psi }_2^{ABC}`$ if, for sufficiently large $`N`$, $`N`$ copies of $`\|\mathrm{\Psi }_1^{ABC}`$ can be transformed with high probability (via local operations and classical communication) to approximately $`N`$ copies of $`\|\mathrm{\Psi }_2^{ABC}`$, and vice versa. The qualifiers in this description are worth a comment or two. “High probability” reflects the fact that, since the local operations may involve measurements, the actual final state may depend on the exact measurement outcomes. “Approximately $`N`$ copies” means more than $`(1ϵ)N`$ copies, for some suitably small $`ϵ`$ determined in advance. We denote this reversibility relation by $$\|\mathrm{\Psi }_1^{ABC}\|\mathrm{\Psi }_2^{ABC}.$$ Two states that are related in this way are essentially equivalent as “entanglement resources”. In the large $`N`$ limit, they may be interconverted with arbitrarily little loss. Our results for entropy and relative entropy of entanglement allow us to place necessary conditions on the reversible manipulation of multiparticle entanglement. For example, if $`\|\mathrm{\Psi }_1^{ABC}\|\mathrm{\Psi }_2^{ABC}`$, then the two states must have exactly the same subsystem entropies. Suppose instead that $`S(\rho _1^A)<S(\rho _2^A)`$. Then the transformation of $`N`$ copies of $`\|\mathrm{\Psi }_1^{ABC}`$ into about $`N`$ copies of $`\|\mathrm{\Psi }_2^{ABC}`$ would involve an increase in the entropy of subsystem $`A`$, which cannot happen on average. In a similar way, we can see that $`\|\mathrm{\Psi }_1^{ABC}`$ and $`\|\mathrm{\Psi }_2^{ABC}`$ must have the same relative entropies of entanglement for every pair of subsystems. Suppose instead that $`E_{r,1}^{BC}<E_{r,2}^{BC}`$. Then the transformation of $`N`$ copies of $`\|\mathrm{\Psi }_1^{ABC}`$ into about $`N`$ copies of $`\|\mathrm{\Psi }_2^{ABC}`$ would require an increase in $`E_r^{BC}`$. This can take place if a measurement is performed on $`A`$, but as we have seen this would necessarily involve a decrease in $`S(\rho ^A)`$. Therefore, reversible transformations of multiparticle entanglement must preserve both subsystem entropies and the entanglement (measured by $`E_r`$) of pairs of subsystems. As a simple example of this, suppose Alice, Bob and Claire share two GHZ states. Each subsystem has an entropy of 2.0 bits. This would also be the case if Alice, Bob and Claire shared three EPR pairs, one between each pair of participants. Does it follow that two GHZs can be transformed reversibly (in the sense described above) into three EPRs? No. If the three parties share two GHZ triples, then Bob and Claire are in a completely unentangled state, with $`E_r^{BC}=0`$. But in the “three EPR” situation, the relative entropy of entanglement $`E_r^{BC}`$ is 1.0 bits, since they share an EPR pair. Thus, two GHZs cannot be reversibly transformed into three EPRs; indeed, $`2N`$ GHZs are inequivalent to $`3N`$ EPRs. Though we have phrased our results for three parties, they are obviously applicable to situations with four or more separated subsystems. In reversible manipulations of multiparticle entanglement, all subsystem entropies (including the entropies of clusters of subsystems) must remain constant, as well as the relative entropies of entanglement of all pairs of subsystems (or clusters of subsystems). ## 10 Remarks The applications discussed here show the power and the versatility of relative entropy methods in attacking problems of quantum information theory. We have derived useful fundamental results in classical and quantum information transfer, quantum data compression, and the manipulation of quantum entanglement. In particular, Donald’s identity proves to be an extremely useful tool for deriving important inequalities. One of the insights provided by quantum information theory is that the von Neumann entropy $`S(\rho )`$ has an interpretation (actually several interpretations) as a measure of the resources necessary to perform an information task. We have seen that the relative entropy also supports such interpretations. We would especially like to draw attention to the results in Sections 3 on the cost of communication and Section 7 on quantum data compression, which are presented here for the first time. We expect that relative entropy techniques will be central to further work in quantum information theory. In particular, we think that they show promise in resolving the many perplexing additivity problems that face the theory at present. Section 5, though not a very strong result in itself, may point the way along this road. The authors wish to acknowledge the invaluable help of many colleagues. T. Cover, M. Donald, M. Neilsen, M. Ruskai, A. Uhlmann and V. Vedral have given us indispensible guidance about the properties and meaning of the relative entropy function. Our work on optimal signal ensembles and the additivity problem was greatly assisted by conversations with C. Fuchs, A. Holevo, J. Smolin, and W. Wootters. Results described here on reversibility for transformations of multiparticle entanglement were obtained in the course of joint work with N. Linden and S. Popescu. We would like to thank the organizers of the AMS special session on “Quantum Information and Computation” for a stimulating meeting and an opportunity to pull together several related ideas into the present paper. We hope it will serve as a spur for the further application of relative entropy methods to problems of quantum information theory.
warning/0004/hep-lat0004012.html	ar5iv	text	# 1 Introduction ## 1 Introduction Recent developments in the treatment of fermions in lattice gauge theory are based on a hermitian lattice Dirac operator $`\gamma _5D`$ which satisfies the Ginsparg-Wilson relation <sup>1</sup><sup>1</sup>1To be precise, the general relation $`\gamma _5D+D\gamma _5=2aD\gamma _5\alpha D`$, where $`\alpha `$ is a local operator, has been proposed in Ref., although the authors in Ref. analyzed “only the simplest case where the matrix $`\alpha `$ is proportional to the unit matrix in Dirac space”. With this qualification in mind, we refer to (1.1) as the “ordinary Ginsparg-Wilson relation” in this paper. The original Ginsparg- Wilson relation is more general as stated above. $$\gamma _5D+D\gamma _5=2aD\gamma _5D$$ (1.1) where the lattice spacing $`a`$ is utilized to make a dimensional consideration transparent, and $`\gamma _5`$ is a hermitian chiral Dirac matrix. An explicit example of the operator satisfying (1.1) and free of species doubling has been given by Neuberger. The relation (1.1) led to an interesting analysis of the notion of index in lattice gauge theory. This index theorem in turn led to a new form of chiral symmetry, and the chiral anomaly is obtained as a non-trivial Jacobian factor under this modified chiral transformation. This chiral Jacobian is regarded as a lattice generalization of the continuum path integral. The very detailed analyses of the lattice chiral Jacobian have been performed-. It is also possible to formulate the lattice index theorem in a manner analogous to the continuum index theorem. An interesting chirality sum rule, which relates the number of zero modes to that of the heaviest states, has also been noticed. In this paper we discuss a generalization of the relation (1.1), which is characterized by a non-negative integer $`k`$. It is shown that the explicit construction of an infinite tower of lattice Dirac operators which satisfy the index theorem is possible, but a large enough lattice is required to accomodate a Dirac operator with a large value of $`k`$. ## 2 Generalized algebra and its representation We discuss a generalization of the algebra (1.1) to the form<sup>2</sup><sup>2</sup>2This relation is obtained from the proposal in Ref. ,$`\gamma _5D+D\gamma _5=2aD\gamma _5\alpha D`$, by choosing $`\alpha `$ as an operator containing $`D`$ itself (and thus Dirac matrices). From a view point of algebra, the original construction in contains two unknown operators and one relation. In our construction, we have a closed algebraic relation for one unknown operator $`D`$, which allows a neat analyis of representation in this Section. This specific algebraically closed realization ,which is characterized by a non-negative integer, has not been discussed in Ref.. $$\gamma _5(\gamma _5D)+(\gamma _5D)\gamma _5=2a^{2k+1}(\gamma _5D)^{2k+2}$$ (2.1) where $`k`$ stands for a non-negative integer and $`k=0`$ corresponds to the ordinary Ginsparg-Wilson relation. When one defines $$H\gamma _5aD$$ (2.2) (2.1) is rewritten as $$\gamma _5H+H\gamma _5=2H^{2k+2}$$ (2.3) or equivalently $$\mathrm{\Gamma }_5H+\mathrm{\Gamma }_5H=0$$ (2.4) where we defined $$\mathrm{\Gamma }_5\gamma _5H^{2k+1}.$$ (2.5) Note that both of $`H`$ and $`\mathrm{\Gamma }_5`$ are hermitian operators. We now discuss a general representation of the algebraic relation (2.4) following the analysis in Appendix of Ref..(In Ref., the algebra was normalized as $`\gamma _5(\gamma _5D)+(\gamma _5D)\gamma _5=a(\gamma _5D)^2`$, but here we use the normalization (2.1) to simplify various expressions.) The relation (2.4) suggests that if $$H\varphi _n=a\lambda _n\varphi _n,(\varphi _n,\varphi _n)=1$$ (2.6) with a real eigenvalue $`a\lambda _n`$ for the hermitian operator $`H`$, then $$H(\mathrm{\Gamma }_5\varphi _n)=a\lambda _n(\mathrm{\Gamma }_5\varphi _n).$$ (2.7) Namely, the eigenvalues $`\lambda _n`$ and $`\lambda _n`$ are always paired if $`\lambda _n0`$ and $`(\mathrm{\Gamma }_5\varphi _n,\mathrm{\Gamma }_5\varphi _n)0`$. We also note the relation, which is derived by sandwiching the relation (2.3) by $`\varphi _n`$, $$(\varphi _n,\gamma _5\varphi _n)=(a\lambda _n)^{2k+1}for\lambda _n0.$$ (2.8) Consequently $$\|(a\lambda _n)^{2k+1}\|=\|(\varphi _n,\gamma _5\varphi _n)\|\varphi _n\gamma _5\varphi _n=1.$$ (2.9) Namely, all the possible eigenvalues are bounded by $$\|\lambda _n\|\frac{1}{a}.$$ (2.10) We thus evaluate the norm of $`\mathrm{\Gamma }_5\varphi _n`$ $`(\mathrm{\Gamma }_5\varphi _n,\mathrm{\Gamma }_5\varphi _n)`$ $`=`$ $`(\varphi _n,(\gamma _5H^{2k+1})(\gamma _5H^{2k+1})\varphi _n)`$ (2.11) $`=`$ $`(\varphi _n,(1H^{2k+1}\gamma _5\gamma _5H^{2k+1}+H^{2(2k+1)})\varphi _n)`$ $`=`$ $`[1(a\lambda _n)^{2(2k+1)}]`$ $`=`$ $`[1(a\lambda _n)^2][1+(a\lambda _n)^2+\mathrm{}+(a\lambda _n)^{4k}]`$ where we used (2.8). By remembering that all the eigenvalues are real, we find that $`\varphi _n`$ is a “highest” state $$\mathrm{\Gamma }_5\varphi _n=0$$ (2.12) only if $$[1(a\lambda _n)^2]=(1a\lambda _n)(1+a\lambda _n)=0$$ (2.13) for the Euclidean positive definite inner product $`(\varphi _n,\varphi _n)_x\varphi _n^{}(x)\varphi _n(x)`$. We thus conclude that the states $`\varphi _n`$ with $`\lambda _n=\pm \frac{1}{a}`$ are not paired by the operation $`\mathrm{\Gamma }_5\varphi _n`$ and $$\gamma _5D\varphi _n=\pm \frac{1}{a}\varphi _n,\gamma _5\varphi _n=\pm \varphi _n$$ (2.14) respectively. These eigenvalues are in fact the maximum or minimum of the possible eigenvalues of $`H/a`$ due to (2.10). As for the vanishing eigenvalues $`H\varphi _n=0`$, we find from (2.4) that $`H\gamma _5\varphi _n=0`$, namely, $`H[(1\pm \gamma _5)/2]\varphi _n=0`$. We thus have $$\gamma _5D\varphi _n=0,\gamma _5\varphi _n=\varphi _nor\gamma _5\varphi _n=\varphi _n.$$ (2.15) To summarize the analyses so far, all the normalizable eigenstates $`\varphi _n`$ of $`\gamma _5D=H/a`$ are categorized into the following 3 classes: (i) $`n_\pm `$ (“zero modes”), $$\gamma _5D\varphi _n=0,\gamma _5\varphi _n=\pm \varphi _n,$$ (2.16) (ii) $`N_\pm `$ (“highest states”), $$\gamma _5D\varphi _n=\pm \frac{1}{a}\varphi _n,\gamma _5\varphi _n=\pm \varphi _n,respectively,$$ (2.17) (iii)“paired states” with $`0<\|\lambda _n\|<1/a`$, $$\gamma _5D\varphi _n=\lambda _n\varphi _n,\gamma _5D(\mathrm{\Gamma }_5\varphi _n)=\lambda _n(\mathrm{\Gamma }_5\varphi _n).$$ (2.18) Note that $`\mathrm{\Gamma }_5(\mathrm{\Gamma }_5\varphi _n)\varphi _n`$ for $`0<\|\lambda _n\|<1/a`$. We thus obtain the index relation $`Tr\mathrm{\Gamma }_5`$ $``$ $`{\displaystyle \underset{n}{}}(\varphi _n,\mathrm{\Gamma }_5\varphi _n)`$ (2.19) $`=`$ $`{\displaystyle \underset{\lambda _n=0}{}}(\varphi _n,\mathrm{\Gamma }_5\varphi _n)+{\displaystyle \underset{0<\|\lambda _n\|<1/a}{}}(\varphi _n,\mathrm{\Gamma }_5\varphi _n)+{\displaystyle \underset{\|\lambda _n\|=1/a}{}}(\varphi _n,\mathrm{\Gamma }_5\varphi _n)`$ $`=`$ $`{\displaystyle \underset{\lambda _n=0}{}}(\varphi _n,\mathrm{\Gamma }_5\varphi _n)`$ $`=`$ $`{\displaystyle \underset{\lambda _n=0}{}}(\varphi _n,(\gamma _5H^{2k+1})\varphi _n)`$ $`=`$ $`{\displaystyle \underset{\lambda _n=0}{}}(\varphi _n,\gamma _5\varphi _n)`$ $`=`$ $`n_+n_{}=index`$ where $`n_\pm `$ stand for the number of normalizable zero modes with $`\gamma _5\varphi _n=\pm \varphi _n`$ in the classification (i) above. We here used the fact that $`\mathrm{\Gamma }_5\varphi _n=0`$ for the “highest states” and that $`\varphi _n`$ and $`\mathrm{\Gamma }_5\varphi _n`$ are orthogonal to each other for $`0<\|\lambda _n\|<1/a`$ since they have eigenvalues with opposite signatures. On the other hand, the relation $`Tr\gamma _5=0`$, which is expected to be valid in (finite) lattice theory, leads to ( by using (2.8)) $`Tr\gamma _5`$ $`=`$ $`{\displaystyle \underset{n}{}}(\varphi _n,\gamma _5\varphi _n)`$ (2.20) $`=`$ $`{\displaystyle \underset{\lambda _n=0}{}}(\varphi _n,\gamma _5\varphi _n)+{\displaystyle \underset{\lambda _n0}{}}(\varphi _n,\gamma _5\varphi _n)`$ $`=`$ $`n_+n_{}+{\displaystyle \underset{\lambda _n0}{}}(a\lambda _n)^{2k+1}=0.`$ In the last line of this relation, all the states except for the “highest states” with $`\lambda _n=\pm 1/a`$ cancel pairwise for $`\lambda _n0`$. We thus obtain a chirality sum rule $`n_+n_{}+N_+N_{}=0`$ or $$n_++N_+=n_{}+N_{}$$ (2.21) where $`N_\pm `$ stand for the number of “highest states” with $`\gamma _5\varphi _n=\pm \varphi _n`$ in the classification (ii) above. These relations show that the chirality asymmetry at vanishing eigenvalues is balanced by the chirality asymmetry at the largest eigenvalues with $`\|\lambda _n\|=1/a`$. It was argued in Ref. that $`N_\pm `$ states are the topological (instanton-related) excitations of the would-be species doublers. All the $`n_\pm `$ and $`N_\pm `$ states are the eigenstates of $`D`$, $`D\varphi _n=0`$ and $`D\varphi _n=(1/a)\varphi _n`$, respectively. If one denotes the number of states in the classification (iii) above by $`2N_0`$, the total number of states (the dimension of the representation) $`N`$ is given by $$N=2(n_++N_++N_0)$$ (2.22) which is expected to be common to all the algebraic relations in (2.1) and to be a constant independent of background gauge field configurations. We note that all the states $`\varphi _n`$ with $`0<\|\lambda _n\|<1/a`$, which appear pairwise with $`\lambda _n=\pm \|\lambda _n\|`$, can be normalized to satisfy the relations $`\mathrm{\Gamma }_5\varphi _n`$ $`=`$ $`[1(a\lambda _n)^{2(2k+1)}]^{1/2}\varphi _n,`$ $`\gamma _5\varphi _n`$ $`=`$ $`(a\lambda _n)^{2k+1}\varphi _n+[1(a\lambda _n)^{2(2k+1)}]^{1/2}\varphi _n.`$ (2.23) Here $`\varphi _n`$ stands for the eigenstate with an eigenvalue opposite to that of $`\varphi _n`$. These states $`\varphi _n`$ cannot be the eigenstates of $`\gamma _5`$ since $`\|(\varphi _n,\gamma _5\varphi _n)\|=\|(a\lambda _n)^{2k+1}\|<1`$. We have thus established that the representation of all the algebraic relations (2.1) has a similar structure. In the next Section, we show that the index $`n_+n_{}`$ is identical to all these algebraic relations if the operator $`\gamma _5D`$ satisfies suitable conditions. ## 3 Chiral Jacobian and the index relation The Euclidean path integral for a fermion is defined by $$𝒟\overline{\psi }𝒟\psi \mathrm{exp}[\overline{\psi }D\psi ]$$ (3.1) where $$\overline{\psi }D\psi \underset{x,y}{}\overline{\psi }(x)D(x,y)\psi (y)$$ (3.2) and the summation runs over all the points on the lattice. The relation (2.4) is re-written as $$\gamma _5\mathrm{\Gamma }_5\gamma _5D+D\mathrm{\Gamma }_5=0$$ (3.3) and thus the Euclidean action is invariant under the global “chiral” transformation $`\overline{\psi }(x)\overline{\psi }^{}(x)=\overline{\psi }(x)+i{\displaystyle \underset{z}{}}\overline{\psi }(z)ϵ\gamma _5\mathrm{\Gamma }_5(z,x)\gamma _5`$ $`\psi (y)\psi ^{}(y)=\psi (y)+i{\displaystyle \underset{w}{}}ϵ\mathrm{\Gamma }_5(y,w)\psi (w)`$ (3.4) with an infinitesimal constant parameter $`ϵ`$. Under this transformation, one obtains a Jacobian factor $$𝒟\overline{\psi }^{}𝒟\psi ^{}=J𝒟\overline{\psi }𝒟\psi $$ (3.5) with $$J=\mathrm{exp}[2iTrϵ\mathrm{\Gamma }_5]=\mathrm{exp}[2iϵ(n_+n_{})]$$ (3.6) where we used the index relation (2.19). We now relate this index appearing in the Jacobian to the Pontryagin index of the gauge field in a smooth continuum limit by following the procedure in Ref.. We start with $$Tr\{\mathrm{\Gamma }_5f(\frac{(\gamma _5D)^2}{M^2})\}=Tr\{\mathrm{\Gamma }_5f(\frac{(H/a)^2}{M^2})\}=n_+n_{}$$ (3.7) Namely, the index is not modified by any regulator $`f(x)`$ with $`f(0)=1`$ and $`f(x)`$ rapidly going to zero for $`x\mathrm{}`$, as can be confirmed by using (2.19). This means that you can use any suitable $`f(x)`$ in the evaluation of the index by taking advantage of this property. We then consider a local version of the index $$tr\{\mathrm{\Gamma }_5f(\frac{(\gamma _5D)^2}{M^2})\}(x,x)=tr\{(\gamma _5H^{2k+1})f(\frac{(\gamma _5D)^2}{M^2})\}(x,x)$$ (3.8) where trace stands for Dirac and Yang-Mills indices; Tr in (3.7) includes a sum over the lattice points $`x`$. A local version of the index is not sensitive to the precise boundary condition , and one may take an infinite volume limit of the lattice in the above expression. We now examine the continuum limit $`a0`$ of the above local expression (3.8)<sup>3</sup><sup>3</sup>3This continuum limit corresponds to the so-called “naive” continuum limit in the context of lattice gauge theory.. We first observe that the term $$tr\{H^{2k+1}f(\frac{(\gamma _5D)^2}{M^2})\}$$ (3.9) goes to zero in this limit. The large eigenvalues of $`H=a\gamma _5D`$ are truncated at the value $`aM`$ by the regulator $`f(x)`$ which rapidly goes to zero for large $`x`$. In other words, the global index of the operator $`TrH^{2k+1}f(\frac{(\gamma _5D)^2}{M^2})O(aM)^{2k+1}`$. We thus examine the small $`a`$ limit of $$tr\{\gamma _5f(\frac{(\gamma _5D)^2}{M^2})\}.$$ (3.10) The operator appearing in this expression is well regularized by the function $`f(x)`$ , and we evaluate the above trace by using the plane wave basis to extract an explicit gauge field dependence. We consider a square lattice where the momentum is defined in the Brillouin zone $$\frac{\pi }{2a}k_\mu <\frac{3\pi }{2a}.$$ (3.11) We assume that the operator $`D`$ is free of species doubling; in other words, the operator $`D`$ blows up rapidly ($`\frac{1}{a}`$) for small $`a`$ in the momentum region corresponding to species doublers. The contributions of doublers are eliminated by the regulator $`f(x)`$ in the above expression, since $$tr\{\gamma _5f(\frac{(\gamma _5D)^2}{M^2})\}(\frac{1}{a})^4f(\frac{1}{(aM)^2})0$$ (3.12) for $`a0`$ if one chooses $`f(x)=e^x`$, for example. We thus examine the above trace in the momentum range of the physical species $$\frac{\pi }{2a}k_\mu <\frac{\pi }{2a}.$$ (3.13) We obtain the limiting $`a0`$ expression $`\underset{a0}{lim}tr\{\gamma _5f({\displaystyle \frac{(\gamma _5D)^2}{M^2}})\}(x,x)`$ (3.14) $`=`$ $`\underset{a0}{lim}tr{\displaystyle _{\frac{\pi }{2a}}^{\frac{\pi }{2a}}}{\displaystyle \frac{d^4k}{(2\pi )^4}}e^{ikx}\gamma _5f({\displaystyle \frac{(\gamma _5D)^2}{M^2}})e^{ikx}`$ $`=`$ $`\underset{L\mathrm{}}{lim}\underset{a0}{lim}tr{\displaystyle _L^L}{\displaystyle \frac{d^4k}{(2\pi )^4}}e^{ikx}\gamma _5f({\displaystyle \frac{(\gamma _5D)^2}{M^2}})e^{ikx}`$ $`=`$ $`\underset{L\mathrm{}}{lim}tr{\displaystyle _L^L}{\displaystyle \frac{d^4k}{(2\pi )^4}}e^{ikx}\gamma _5f({\displaystyle \frac{(i\gamma _5\overline{)}D)^2}{M^2}})e^{ikx}`$ $``$ $`tr\{\gamma _5f({\displaystyle \frac{\overline{)}D^2}{M^2}})\}`$ where we first take the limit $`a0`$ with fixed $`k_\mu `$ in $`Lk_\mu L`$, and then take the limit $`L\mathrm{}`$. This procedure is justified if the integral is well convergent <sup>4</sup><sup>4</sup>4 To be precise, we deal with an integral of the structure $`_{\frac{\pi }{2a}}^{\frac{\pi }{2a}}𝑑xf_a(x)=_L^{\frac{\pi }{2a}}𝑑xf_a(x)+_L^L𝑑xf_a(x)+_{\frac{\pi }{2a}}^L𝑑xf_a(x)`$ where $`f_a(x)`$ depends on the parameter $`a`$. (A generalization to a 4-dimensional integral is straightforward.) We thus have to prove that both of $`lim_{a0}_L^{\frac{\pi }{2a}}𝑑xf_a(x)`$ and $`lim_{a0}_{\frac{\pi }{2a}}^L𝑑xf_a(x)`$ can be made arbitrarily small if one lets $`L`$ to be large. A typical integral we encounter in lattice theory has a generic structure $`lim_{a0}_{\pi /2a}^{\pi /2a}𝑑xe^{[\mathrm{sin}^2ax+(1\mathrm{cos}2ax)^2]/(a^2M^2)}=lim_{a0}_{ϵ\pi /2a}^{\pi /2a}𝑑xe^{[\mathrm{sin}^2ax+(1\mathrm{cos}2ax)^2]/(a^2M^2)}+lim_{a0}_{\pi /2a}^{ϵ\pi /2a}𝑑xe^{[\mathrm{sin}^2ax+(1\mathrm{cos}2ax)^2]/(a^2M^2)}+lim_{a0}_{ϵ\pi /2a}^{ϵ\pi /2a}𝑑xe^{[\mathrm{sin}^2ax+(1\mathrm{cos}2ax)^2]/(a^2M^2)}=lim_{a0}_{ϵ\pi /2a}^{ϵ\pi /2a}𝑑xe^{[\mathrm{sin}^2ax+(1\mathrm{cos}2ax)^2]/(a^2M^2)}=lim_L\mathrm{}_L^L𝑑xe^{x^2/M^2}`$ and satisfies the above criterion, if one chooses the regulator $`f(x)=e^x`$: Here $`ϵ`$ is an arbitrary small fixed parameter, and the left-hand side of this relation stands for a conventional lattice calculation and the right-hand side stands for a continuum calculation.. We also assumed that the operator $`D`$ satisfies the following relation in the limit $`a0`$ $`De^{ikx}h(x)`$ $``$ $`e^{ikx}(\overline{)}k+i\overline{)}g\overline{)}A)h(x)`$ (3.15) $`=`$ $`i(\overline{)}+ig\overline{)}A)(e^{ikx}h(x))i\overline{)}D(e^{ikx}h(x))`$ for any fixed $`k_\mu `$, ($`\frac{\pi }{2a}<k_\mu <\frac{\pi }{2a}`$), and a sufficiently smooth function $`h(x)`$. The function $`h(x)`$ corresponds to the gauge potential in our case, which in turn means that the gauge potential $`A_\mu (x)`$ is assumed to vary very little over the distances of the elementary lattice spacing. Our final expression (3.14) in the limit $`M\mathrm{}`$ reproduces the Pontryagin number in the continuum formulation $`\underset{M\mathrm{}}{lim}tr\{\gamma _5f(\overline{)}D^2/M^2)\}`$ (3.16) $`=`$ $`\underset{M\mathrm{}}{lim}tr{\displaystyle \frac{d^4k}{(2\pi )^4}e^{ikx}\gamma _5f(\overline{)}D^2/M^2)e^{ikx}}`$ $`=`$ $`\underset{M\mathrm{}}{lim}tr{\displaystyle \frac{d^4k}{(2\pi )^4}\gamma _5f\{(ik_\mu +D_\mu )^2/M^2+\frac{ig}{4}[\gamma ^\mu ,\gamma ^\nu ]F_{\mu \nu }/M^2\}}`$ $`=`$ $`\underset{M\mathrm{}}{lim}trM^4{\displaystyle \frac{d^4k}{(2\pi )^4}\gamma _5f\{(ik_\mu +D_\mu /M)^2+\frac{ig}{4}[\gamma ^\mu ,\gamma ^\nu ]F_{\mu \nu }/M^2\}}`$ where the remaining trace stands for Dirac and Yang-Mills indices. We also used the relation $$\overline{)}D^2=D_\mu D^\mu +\frac{ig}{4}[\gamma ^\mu ,\gamma ^\nu ]F_{\mu \nu }$$ (3.17) and the rescaling of the variable $`k_\mu Mk_\mu `$. By noting $`tr\gamma _5=tr\gamma _5[\gamma ^\mu ,\gamma ^\nu ]=0`$, the above expression ( after expansion in powers of $`1/M`$) is written as (with $`ϵ^{1234}=1`$) $`\underset{M\mathrm{}}{lim}tr\gamma _5f(\overline{)}D^2/M^2)`$ $`=`$ $`tr\gamma _5{\displaystyle \frac{1}{2!}}\{{\displaystyle \frac{ig}{4}}[\gamma ^\mu ,\gamma ^\nu ]F_{\mu \nu }\}^2{\displaystyle \frac{d^4k}{(2\pi )^4}f^{\prime \prime }(k_\mu k^\mu )}`$ (3.18) $`={\displaystyle \frac{g^2}{32\pi ^2}}trϵ^{\mu \nu \alpha \beta }F_{\mu \nu }F_{\alpha \beta }`$ where we used $`{\displaystyle \frac{d^4k}{(2\pi )^4}f^{\prime \prime }(k_\mu k^\mu )}`$ $`=`$ $`{\displaystyle \frac{1}{16\pi ^2}}{\displaystyle _0^{\mathrm{}}}f^{\prime \prime }(x)x𝑑x`$ (3.19) $`=`$ $`{\displaystyle \frac{1}{16\pi ^2}}`$ with $`x=k_\mu k^\mu >0`$ in our metric. When one combines (3.7) and (3.18), one reproduces the Atiyah-Singer index theorem (in continuum $`R^4`$ space). We note that a local version of the index (anomaly) is valid for Abelian theory also. The global index (3.7) as well as a local version of the index (3.8) are both independent of the regulator $`f(x)`$ provided $$f(0)=1,f(\mathrm{})=0,f^{}(x)x\|_{x=0}=f^{}(x)x\|_{x=\mathrm{}}=0.$$ (3.20) We have thus established that the lattice index in (3.7) for any algebraic relation in (2.1) is related to the Pontryagin index in a smooth continuum limit as $$n_+n_{}=d^4x\frac{g^2}{32\pi ^2}trϵ^{\mu \nu \alpha \beta }F_{\mu \nu }F_{\alpha \beta }$$ (3.21) by assuming the quite general properties of the basic operator $`D`$ only: The basic relation (2.1) with hermitian $`\gamma _5D`$ and the continuum limit property (3.15) without species doubling in the limit $`a0`$. This shows that the instanton-related topological property is identical for all the algebraic relations in (2.1), and the Jacobian factor (3.6) in fact contains the correct chiral anomaly. (We are implicitly assuming that the index (3.7) does not change in the process of taking a continuum limit.) Our result is naturally consistent with the calculation of chiral anomaly by different methods in and . ## 4 Explicit example of the lattice Dirac operator with $`k`$=1 We now discuss an explicit construction of the lattice Dirac operator which satisfies the generalized algebraic relation (2.1) with $`k=1`$, though a generalization to an arbitrary $`k`$ is straightforward as is described in Section 5 later. For this purpose, we first briefly review the construction of the Neuberger’s overlap Dirac operator for the ordinary Ginsparg-Wilson relation. We start with the conventional Wilson fermion operator $`D_W`$ defined by $`D_W(x,y)`$ $``$ $`i\gamma ^\mu C_\mu (x,y)+B(x,y){\displaystyle \frac{1}{a}}m_0\delta _{x,y},`$ $`C_\mu (x,y)`$ $`=`$ $`{\displaystyle \frac{1}{2a}}[\delta _{x+\widehat{\mu }a,y}U_\mu (y)\delta _{x,y+\widehat{\mu }a}U_\mu ^{}(x)],`$ $`B(x,y)`$ $`=`$ $`{\displaystyle \frac{r}{2a}}{\displaystyle \underset{\mu }{}}[2\delta _{x,y}\delta _{y+\widehat{\mu }a,x}U_\mu ^{}(x)\delta _{y,x+\widehat{\mu }a}U_\mu (y)],`$ $`U_\mu (y)`$ $`=`$ $`\mathrm{exp}[iagA_\mu (y)],`$ (4.1) where we added a constant mass term to $`D_W`$ for later convenience. The parameter $`r`$ stands for the Wilson parameter. Our matrix convention is that $`\gamma ^\mu `$ are anti-hermitian, $`(\gamma ^\mu )^{}=\gamma ^\mu `$, and thus $`\overline{)}C\gamma ^\mu C_\mu (n,m)`$ is hermitian $$\overline{)}C^{}=\overline{)}C.$$ (4.2) The operator $`D`$ introduced by Neuberger, which satisfies the conventional Ginsparg-Wilson relation (1.1), has an explicit expression $$aD=\frac{1}{2}[1+\gamma _5\frac{H_W}{\sqrt{H_W^2}}]=\frac{1}{2}[1+D_W\frac{1}{\sqrt{D_W^{}D_W}}]$$ (4.3) where $`D_W=\gamma _5H_W`$ is the Wilson operator defined above, and $`H_W`$ is hermitian $`H_W^{}=H_W`$. The physical meaning of this construction becomes more transparent if one considers (naive) near continuum configurations specified by a small $`a`$ limit with the parameters $`r/a`$ and $`m_0/a`$ kept finite. We can then approximate the operator $`D_W`$ by $$D_Wi\overline{)}D+M_n$$ (4.4) for each species doubler, where the mass parameters $`M_n`$ stand for $`M_0=\frac{m_0}{a}`$ and one of $`{\displaystyle \frac{2r}{a}}{\displaystyle \frac{m_0}{a}},(4,1);{\displaystyle \frac{4r}{a}}{\displaystyle \frac{m_0}{a}},(6,1)`$ $`{\displaystyle \frac{6r}{a}}{\displaystyle \frac{m_0}{a}},(4,1);{\displaystyle \frac{8r}{a}}{\displaystyle \frac{m_0}{a}},(1,1)`$ (4.5) for $`n=115`$; we denoted ( multiplicity, chiral charge ) in the bracket for species doublers. Here we used the relation valid in the near continuum configurations for the physical species, for example, $`D_W(k)`$ $`=`$ $`{\displaystyle \underset{\mu }{}}\gamma ^\mu {\displaystyle \frac{\mathrm{sin}ak_\mu }{a}}+{\displaystyle \frac{r}{a}}{\displaystyle \underset{\mu }{}}(1\mathrm{cos}ak_\mu ){\displaystyle \frac{m_0}{a}}`$ (4.6) $``$ $`\gamma ^\mu k_\mu {\displaystyle \frac{m_0}{a}}`$ in the momentum representation with vanishing gauge field. In a symbolic notation, one can then write the overlap Dirac operator as $`aD`$ $``$ $`{\displaystyle \underset{n=0}{\overset{15}{}}}{\displaystyle \frac{1}{2}}[1+(i\overline{)}D+M_n){\displaystyle \frac{1}{\sqrt{\overline{)}D^2+M_n^2}}}]\|nn\|,`$ $`a\gamma _5D`$ $``$ $`{\displaystyle \underset{n=0}{\overset{15}{}}}(1)^n\gamma _5{\displaystyle \frac{1}{2}}[1+(i\overline{)}D+M_n){\displaystyle \frac{1}{\sqrt{\overline{)}D^2+M_n^2}}}]\|nn\|.`$ (4.7) Here we explicitly write the projection $`\|nn\|`$ for each species doubler. If one chooses the mass parameters so that $$M_0=\frac{m_0}{a}<0,M_n>0forn0$$ (4.8) namely $$0<m_0<2r$$ (4.9) and if one lets all the mass parameters $`\|M_n\|`$ become large, one obtains $`a\gamma _5D`$ $``$ $`\gamma _5{\displaystyle \frac{1}{2}}[{\displaystyle \frac{i\overline{)}D}{\|M_0\|}}+{\displaystyle \frac{1}{2}}{\displaystyle \frac{\overline{)}D^2}{M_0^2}}]forn=0,`$ $`a\gamma _5D`$ $``$ $`(1)^n\gamma _5{\displaystyle \frac{1}{2}}[2+{\displaystyle \frac{i\overline{)}D}{M_n}}{\displaystyle \frac{1}{2}}{\displaystyle \frac{\overline{)}D^2}{M_n^2}}]forn0.`$ (4.10) If one chooses $`m_0`$ to satisfy $$2a\|M_0\|=2m_0=1$$ (4.11) one recovers the correctly normalized continuum Dirac operator for the physical species and $`\gamma _5D(1)^n\gamma _5\frac{1}{a}`$ for unphysical species doublers. In particular, the first relation in (4.10) can then be written as $$Ha\gamma _5D\gamma _5ai\overline{)}D+\gamma _5(\gamma _5ai\overline{)}D)^2$$ (4.12) which ensures the conventional Ginsparg-Wilson relation in the leading order. These properties become important in the following discussion. ### 4.1 Generalized algebra with $`k=1`$ We now come back to the generalized algebra (2.1) with $`k=1`$ $$H\gamma _5+\gamma _5H=2H^4$$ (4.13) where $`H=a\gamma _5D`$ and $`\mathrm{\Gamma }_5=\gamma _5H^3`$. This algebraic relation implies that $$\gamma _5H^2=[\gamma _5H+H\gamma _5]HH[\gamma _5H+H\gamma _5]+H^2\gamma _5=H^2\gamma _5$$ (4.14) Namely, the algebraic relation (4.13) is equivalent to the two relations $`H^3\gamma _5+\gamma _5H^3=2H^6,`$ $`\gamma _5H^2H^2\gamma _5=0.`$ (4.15) If one defines $`H_{(3)}H^3`$, the first relation of (4.15) becomes $$H_{(3)}\gamma _5+\gamma _5H_{(3)}=2H_{(3)}^2$$ (4.16) with $`\mathrm{\Gamma }_5=\gamma _5H_{(3)}`$, which is identical to the conventional Ginsparg-Wilson relation (1.1). We utilize this property to construct a solution to (4.15). Note that the operator $`\mathrm{\Gamma }_5`$ is identical in these three ways of writing in (4.13), (4.15), and (4.16). The physical condition for the operator $`H`$ in (4.13) in the near continuum configuration is (Cf.(4.12)) $$H\gamma _5ai\overline{)}D+\gamma _5(\gamma _5ai\overline{)}D)^4$$ (4.17) and thus $`H_{(3)}`$ in (4.16) should satisfy $`H_{(3)}`$ $``$ $`[\gamma _5ai\overline{)}D+\gamma _5(\gamma _5ai\overline{)}D)^4]^3`$ (4.18) $``$ $`(\gamma _5ai\overline{)}D)^3+\gamma _5(\gamma _5ai\overline{)}D)^6`$ as can be confirmed by noting $`\gamma _5\overline{)}D+\overline{)}D\gamma _5=0`$. Here only the leading terms in chiral symmetric and chiral symmetry breaking terms respectively are written. One can thus construct a solution for $`H_{(3)}`$ by $$H_{(3)}=\frac{1}{2}\gamma _5[1+D_W^{(3)}\frac{1}{\sqrt{(D_W^{(3)})^{}D_W^{(3)}}}]$$ (4.19) where we defined $`D_W^{(3)}`$ by<sup>5</sup><sup>5</sup>5It is also possible to use $`D_W^{(3)}i(\overline{)}C)^3+(B\frac{m_0}{a})^3`$, or any suitable (ultra-local) operator which satisfies $`\gamma _5D_W^{(3)}=(\gamma _5D_W^{(3)})^{}`$ and (4.22). $$D_W^{(3)}i(\overline{)}C)^3+(B)^3(\frac{m_0}{a})^3$$ (4.20) The operators $`\overline{)}C`$,$`B`$ and the parameter $`m_0/a`$ are the same as in the original Wilson fermion operator (4.1). By rewriting (4.19) as $$H_{(3)}=\frac{1}{2}\gamma _5[1+\gamma _5H_W^{(3)}\frac{1}{\sqrt{H_W^{(3)}H_W^{(3)}}}]$$ (4.21) in terms of the hermitian $`H_W^{(3)}\gamma _5D_W^{(3)}=(H_W^{(3)})^{}`$ and comparing it with (4.3), one can confirm that our operator $`H_{(3)}`$ satisfies the relation (4.16). The condition (4.18) is satisfied by noting $$D_W^{(3)}i(\overline{)}D)^3+(M_n^{(3)})^3$$ (4.22) in the near continuum configuration, where the mass parameters are given by $`(M_0^{(3)})^3`$ $``$ $`({\displaystyle \frac{m_0}{a}})^3`$ $`(M_n^{(3)})^3`$ $``$ $`\{({\displaystyle \frac{2r}{a}})^3({\displaystyle \frac{m_0}{a}})^3,({\displaystyle \frac{4r}{a}})^3({\displaystyle \frac{m_0}{a}})^3,({\displaystyle \frac{6r}{a}})^3({\displaystyle \frac{m_0}{a}})^3,({\displaystyle \frac{8r}{a}})^3({\displaystyle \frac{m_0}{a}})^3\}`$ (4.23) $`forn0.`$ Although we have the same condition on the parameters as before $$0<m_0<2r$$ (4.24) to avoid the species doublers, the value of $`m_0`$ itself is now required to satisfy $$2(m_0)^3=1$$ (4.25) to ensure the properly normalized physical condition (4.18). ### 4.2 Reconstruction of $`H`$ from $`H_{(3)}`$ We now discuss how to reconstruct $`H`$, which satisfies (4.13), from $`H_{(3)}`$ defined above. The basic idea is to take a real cubic root of $`H_{(3)}`$ as $$H=(H_{(3)})^{1/3}$$ (4.26) in such a manner that $`H`$ thus obtained satisfies the second constraint in (4.15). For this purpose, we first recall the essence of the general representation of the algebra (2.1) analyzed in Section 2, which is applicable to (4.16) as well. If one defines the eigenvalue problem $$H_{(3)}\varphi _n=(a\lambda _n)^3\varphi _n,(\varphi _n,\varphi _n)=1$$ (4.27) one can classify the eigenstates into the 3 classes: (i) $`n_\pm `$ (“zero modes”), $$H_{(3)}\varphi _n=0,\gamma _5\varphi _n=\pm \varphi _n,$$ (4.28) (ii) $`N_\pm `$ (“highest states”), $$H_{(3)}\varphi _n=\pm \varphi _n,\gamma _5\varphi _n=\pm \varphi _n,respectively,$$ (4.29) (iii)“paired states” with $`0<\|(a\lambda _n)^3\|<1`$, $$H_{(3)}\varphi _n=(a\lambda _n)^3\varphi _n,H_{(3)}(\mathrm{\Gamma }_5\varphi _n)=(a\lambda _n)^3(\mathrm{\Gamma }_5\varphi _n).$$ (4.30) where $$\mathrm{\Gamma }_5=\gamma _5H_{(3)}.$$ (4.31) Note that $`\mathrm{\Gamma }_5(\mathrm{\Gamma }_5\varphi _n)\varphi _n`$ for $`0<\|(a\lambda _n)^3\|<1`$. We obtain the index relation $`Tr\mathrm{\Gamma }_5`$ $``$ $`{\displaystyle \underset{n}{}}(\varphi _n,\mathrm{\Gamma }_5\varphi _n)`$ (4.32) $`=`$ $`{\displaystyle \underset{\lambda _n=0}{}}(\varphi _n,\gamma _5\varphi _n)`$ $`=`$ $`n_+n_{}=index`$ where $`n_\pm `$ stand for the number of normalizable zero modes in the classification (i) above. We also have a chirality sum rule $$n_++N_+=n_{}+N_{}$$ (4.33) where $`N_\pm `$ stand for the number of “highest states” in the classification (ii) above. If one denotes the number of states in the classification (iii) above by $`2N_0`$, the total number of states (the dimension of the representation) $`N`$ is given by $$N=2(n_++N_++N_0)$$ (4.34) which is expected to be common to all the fermion operators defined on the same lattice. Also, all the states $`\varphi _n`$ with $`0<\|(a\lambda _n)^3\|<1`$, which appear pairwise with $`(a\lambda _n)^3=\pm \|(a\lambda _n)^3\|`$, can be normalized to satisfy the relations $`\mathrm{\Gamma }_5\varphi _n`$ $`=`$ $`[1(a\lambda _n)^6]^{1/2}\varphi _n,`$ $`\gamma _5\varphi _n`$ $`=`$ $`(a\lambda _n)^3\varphi _n+[1(a\lambda _n)^6]^{1/2}\varphi _n,`$ (4.35) where $`\varphi _n`$ stands for the eigenstate with an eigenvalue opposite to that of $`\varphi _n`$. Based on these general results in Section 2, we first observe that the index $`n_+n_{}`$ in (4.32) is identical to the index of the expected solution of (4.13), although $`H_{(3)}`$ satisfies (4.18). This observation is based on the relation $$n_+n_{}\underset{n}{}(\varphi _n,\mathrm{\Gamma }_5f((H_{(3)})^2/(aM)^6)\varphi _n)$$ (4.36) which is valid for any regulator with $`f(0)=1`$. One can perform the same analysis as in (3.7) in Section 3: The basic ingredient is the condition (4.18) for a physical momentum region in the smooth continuum limit and the absence of species doublers. The calculation analogous to (3.14) then gives $$n_+n_{}=\underset{M\mathrm{}}{lim}Tr\gamma _5f(\frac{\overline{)}D^6}{M^6})=\underset{M\mathrm{}}{lim}Tr\gamma _5g(\frac{\overline{)}D^2}{M^2})$$ (4.37) with $`g(x)f(x^3)`$ and $`g(0)=1`$. The right-hand side of this relation shows that the present index is identical to the index of the general operator in (2.1), which includes an expected solution of (4.13). Due to the chirality sum rule (4.33), we also obtain the same value of $`N_+N_{}`$ as for an expected solution of (4.13). The agreement of the index of $`H_{(3)}`$ with the index of the expected solution $`H`$ of (4.13) suggests that we can define $`H`$ operationally by $$H\varphi _na\lambda _n\varphi _n$$ (4.38) by using the same set of eigenfunctions and (the cubic roots of) eigenvalues $$\{\varphi _n\},\{a\lambda _n\}$$ (4.39) as for $`H_{(3)}`$ in (4.27). Note that the operator $`\mathrm{\Gamma }_5=\gamma _5H_{(3)}=\gamma _5H^3`$, which reverses the signature of eigenvalues of “paired states” and defines the index, is consistently chosen to be identical for (4.16) and for (4.38)<sup>6</sup><sup>6</sup>6This means that an explicit calculation of the chiral Jacobian (and chiral anomaly) for the theory defined by (4.13) is performed by $`Tr\mathrm{\Gamma }_5=Tr(\gamma _5H_{(3)})`$ in terms of $`H_{(3)}`$ in (4.19).. We can then confirm the second constraint in (4.15) and the defining algebraic relation (4.13) for any “paired state” $`\varphi _n`$, $`[H^2\gamma _5\gamma _5H^2]\varphi _n`$ $`=`$ $`H^2\gamma _5\varphi _n\gamma _5(a\lambda _n)^2\varphi _n`$ (4.40) $`=`$ $`H^2\{(a\lambda _n)^3\varphi _n+[1(a\lambda _n)^6]^{1/2}\varphi _n\}`$ $`(a\lambda _n)^2\{(a\lambda _n)^3\varphi _n+[1(a\lambda _n)^6]^{1/2}\varphi _n\}`$ $`=`$ $`0`$ and $$[\mathrm{\Gamma }_5H+H\mathrm{\Gamma }_5]\varphi _n=\mathrm{\Gamma }_5(a\lambda _n)\varphi _na\lambda _n(\mathrm{\Gamma }_5\varphi _n)=0$$ (4.41) where we used the relations in (4.35) and the definition (4.38). For “zero modes” and the “highest states”, which are the eigenstates of $`\gamma _5`$, the condition $`[H^2\gamma _5\gamma _5H^2]\varphi _n=0`$ obviously holds, and the relation $`[\mathrm{\Gamma }_5H+H\mathrm{\Gamma }_5]\varphi _n=0`$ is also confirmed. The general representation of the algebra (4.13) is obtained from the standard representation, which is defined by $`H`$ in (4.38), $`\gamma _5`$ in (4.35), and the state vectors $`\{\varphi _n\}`$ in (4.39), by applying a suitable unitary transformation. ## 5 Discussion When one considers the algebraic relation with a constant $`R`$ $$\gamma _5(\gamma _5D)+(\gamma _5D)\gamma _5=2Ra^{2k+1}(\gamma _5D)^{2k+2}$$ (5.1) instead of (2.1), one can eliminate the paprameter $`R`$ by a scale transformation $$DD^{}=R^{1/(2k+1)}D.$$ (5.2) The path integral $$𝒟\overline{\psi }𝒟\psi \mathrm{exp}[\overline{\psi }D^{}\psi ]$$ (5.3) is equivalent to $$𝒟\overline{\psi }𝒟\psi \mathrm{exp}[\overline{\psi }D\psi ]$$ (5.4) after absorbing the parameter $`R^{1/(2k+1)}`$ into $`\overline{\psi }`$, at least in a well regularized lattice path integral. Consequently, the parameter $`R`$ and also the factor $`a^{2k+1}`$ do not have an intrinsic physical significance<sup>7</sup><sup>7</sup>7 However, when one includes a Yukawa interaction, for example, this scaling argument need to be refined.. In contrast, the power of $`(\gamma _5D)^{2k+2}`$ in the right-hand side of (5.1) has an intrinsic physical meaning. One may recall the near continuum expressions (4.12) and (4.17) $`H`$ $``$ $`\gamma _5ai\overline{)}D+\gamma _5(\gamma _5ai\overline{)}D)^2fork=0,`$ $`H`$ $``$ $`\gamma _5ai\overline{)}D+\gamma _5(\gamma _5ai\overline{)}D)^4fork=1`$ (5.5) respectively. The first terms in these expressions stand for the leading terms in chiral symmetric terms, and the second terms in these expressions stand for the leading terms in chiral symmetry breaking terms. This shows that one can improve the chiral symmetry<sup>8</sup><sup>8</sup>8To avoid the mis-understanding, we note that the improvement of chiral symmetry here is meant in the sense of Wilsonian renormalization group. The chiral symmetry breaking term becomes more irrelevant for larger $`k`$, and this should be interesting from a view point of regularization of field theory in general. Also, the approach to the continuum Dirac operator is controlled by two parameters, for example, by letting $`klarge`$ and $`asmall`$ simultaneously. by choosing a large parameter $`k`$. The Dirac operator for such a general value of $`k`$ is constructed by rewriting (2.1) as a set of relations (see (4.14)) $`H^{2k+1}\gamma _5+\gamma _5H^{2k+1}=2H^{2(2k+1)},`$ $`H^2\gamma _5\gamma _5H^2=0,`$ (5.6) with $`H=a\gamma _5D`$. The first of these relations (5.6) becomes identical to the ordinary Ginsparg-Wilson relation (1.1) if one defines $`H_{(2k+1)}H^{2k+1}`$. One can construct a solution to (5.6) by following the prescription in Section 4 $$H_{(2k+1)}=\frac{1}{2}\gamma _5[1+D_W^{(2k+1)}\frac{1}{\sqrt{(D_W^{(2k+1)})^{}D_W^{(2k+1)}}}]$$ (5.7) where $$D_W^{(2k+1)}i(\overline{)}C)^{2k+1}+B^{2k+1}(\frac{m_0}{a})^{2k+1}$$ (5.8) The operator $`H`$ is then finally defined by (in the representation where $`H_{(2k+1)}`$ is diagonal) $$H=(H_{(2k+1)})^{1/2k+1}$$ (5.9) in such a manner that the second relation of (5.6) is satisfied. This condition is in deed satisfied as a generalization of (4.40) in the representation where $`H_{(2k+1)}`$ is diagonal. We use the relation (2.23) in this proof. Also the conditions $`0<m_0<2r`$ and $$2m_0^{2k+1}=1$$ (5.10) ensure a proper normalization of the Dirac operator $`H`$. However, one need to use a large enough lattice to accomodate the operator $`H`$ with a large $`k`$, since the operator (5.8) correlates lattice points far apart from each other for a large $`k`$. An explicit analysis of the locality property of our operator $`H`$ as in Ref. is left as an important problem. In the context of lattice simulation, it would be interesting to see how the chiral properties are modified if one uses the operator with $`k=1`$, which has been analyzed in detail in this paper, instead of the conventional overlap Dirac operator with $`k=0`$. To detect the possible effects of $`k0`$ in a reliable way, it is expected that one would have to consider a sufficiently large lattice and those observables which are sensitive to low energy excitations. As for the chiral fermions on the lattice, our general algebra (2.1) satisfies the decomposition $$D=\frac{(1+\gamma _5)}{2}D\frac{(1\widehat{\gamma }_5)}{2}+\frac{(1\gamma _5)}{2}D\frac{(1+\widehat{\gamma }_5)}{2}$$ (5.11) with $$\widehat{\gamma }_5\gamma _52a^{2k+1}(\gamma _5D)^{2k+1},(\widehat{\gamma }_5)^2=1$$ (5.12) by noting $`\gamma _5(\gamma _5D)^2=(\gamma _5D)^2\gamma _5`$. This decomposition has the same form as for the overlap operator $`D`$ satisfying the ordinary Ginsparg-Wilson relation. It is thus expected that one can apply the same considerations as in Refs. and to our general Dirac operator also. In particular, the fermion number non-conservation of the chiral theory defined by $`{\displaystyle 𝒟\overline{\psi }𝒟\psi \mathrm{exp}\{\overline{\psi }D_L\psi \}}`$ $`{\displaystyle 𝒟\overline{\psi }𝒟\psi \mathrm{exp}\{\overline{\psi }\frac{(1+\gamma _5)}{2}D\frac{(1\widehat{\gamma }_5)}{2}\psi \}}`$ (5.13) follows from the fermion number transformation $$\psi e^{i\alpha }\psi ,\overline{\psi }\overline{\psi }e^{i\alpha }.$$ (5.14) If one remembers that the functional spaces of the variables $`\psi `$ and $`\overline{\psi }`$ are specified by the projection operators $`(1\widehat{\gamma }_5)/2`$ and $`(1+\gamma _5)/2`$, respectively, the Jacobian factor for the transformation (5.14) is given by $`J`$ $`=`$ $`\mathrm{exp}\{i\alpha Tr[{\displaystyle \frac{(1+\gamma _5)}{2}}{\displaystyle \frac{(1\widehat{\gamma }_5)}{2}}]\}`$ (5.15) $`=`$ $`\mathrm{exp}\{i\alpha Tr[\gamma _5(\gamma _5aD)^{2k+1}]\}=\mathrm{exp}\{i\alpha [n_+n_{}]\}`$ where the index is defined in (2.19). In conclusion, we have shown that the general idea of Ginsparg and Wilson can be precisely realized as a closed algebraic relation (2.1) and it admits of the explicit construction of an infinite tower of new lattice Dirac operators as a generalization of the overlap Dirac operator. This should be interesting in the context of the regularization of field theory in general. Acknowledgement The present work was initiated when I was visiting at Center for Subatomic Structure of Matter(CSSM), University of Adelaide. I am grateful to David Adams and T-W. Chiu for stimulating discussions, and to Anthony Williams and David Adams for their hospitality at CSSM.
warning/0004/astro-ph0004034.html	ar5iv	text	# CURRENT STATUS OF COSMOLOGICAL MDM MODELS ## 1 Model description We considered a family of models with the following free parameters: * $`\sigma _8[0.47,0.61]`$, (15 models with step 0.01); * $`n[0.8,1.4]`$, (7 models with step 0.1); * $`\mathrm{\Omega }_\nu [0,0.4]`$, (5 models with step 0.1); * $`\mathrm{\Omega }_b[0.01,0.11]`$, (6 models with step 0.02); * $`h[0.45,0.7]`$, (6 models with step 0.05). We use the analytic approximation of the transfer function by Novosyadlyj et al (1998)$`^\mathrm{?}`$. Altogether we have 18900 variants of the model. The derived parameters are abundance of the cold matter, $`\mathrm{\Omega }_{cm}=1\mathrm{\Omega }_\nu \mathrm{\Omega }_b`$, and the contribution of tensor mode to large-scale CMBR anisotropy, T/S. Our goal is to constrain the model parameters by data of the mass function of galaxy clusters and $`\mathrm{\Delta }T/T`$ anisotropy in both large ($`\mathrm{}10`$) and small ($`\mathrm{}200`$) angular scales. ## 2 Mass function of galaxy clusters The number of massive halos with the mass larger then $`M`$ is calculated with help of Press-Schechter formalism$`^\mathrm{?}`$. Observational data are taken from Bahcall & Cen paper $`^\mathrm{?}`$. The $`\chi ^2`$ analysis allows us to delimit the amplitude of the power spectrum at cluster scale with high accuracy, $`\sigma _8=0.52\pm 0.01`$; taking into account the uncertainties of the Press-Schechter approximation and experimental systematics enhances the total errorbar by a value of $`0.04`$ $`^\mathrm{?}`$, $`^\mathrm{?}`$. Other parameters ($`n,\mathrm{\Omega }_\nu ,\mathrm{\Omega }_b,h`$) are not constrained within their ranges by the cluster data. ## 3 CMBR Anisotropy The contribution of cosmic gravitational waves into the large-scale CMBR anisortopy is estimated by the T/S parameter: $$\left(\frac{\mathrm{\Delta }T}{T}\right)^2_{10^o}=\mathrm{S}+\mathrm{T}=\mathrm{S}\left(1+\frac{\mathrm{T}}{\mathrm{S}}\right)1.1\times 10^{10},$$ where S is the contribution of the perturbations of matter density normalized by $`\sigma _8=0.52`$: $$\mathrm{S}=\underset{\mathrm{}=2}{\overset{\mathrm{}}{}}\mathrm{S}_{\mathrm{}}\mathrm{W}_{\mathrm{}},\mathrm{S}_{\mathrm{}}=\frac{2\mathrm{}+1}{64\pi }\mathrm{AH}_0^{\mathrm{n}+3}\frac{\mathrm{\Gamma }(3\mathrm{n})\mathrm{\Gamma }(\mathrm{}+(\mathrm{n}1)/2)}{\mathrm{\Gamma }^2(2\mathrm{n}/2)\mathrm{\Gamma }(\mathrm{}+(5\mathrm{n})/2)},\mathrm{W}_{\mathrm{}}=\mathrm{exp}\left[\left(\frac{2\mathrm{}+1}{27}\right)^2\right],$$ $`A`$ and $`W_{\mathrm{}}`$ are the normalization constant and DMR window function, respectively. The accuracy of this approximation is better than $`3\%`$, the harmonics with $`\mathrm{}{}_{}{}^{<}10`$ ensure the dominant contribution. The result of calculation of T/S$`[0,3]`$ is presented in Fig.1. The value T/S icreases linearly with $`h`$ and decreases with $`\mathrm{\Omega }_\nu `$ growing, therefore the curves T/S with the maximum parameter $`\mathrm{\Omega }_\nu =0.4`$ can be used to put the lower limit on T/S (see Conclusions). Taking a moderate T/S$`<0.5`$ and nearly flat power spectrum ($`0.92n1.02`$), we put an upper limit on the Hubble constant, $`h<0.6`$, and lower limit on the hot dark matter abundance, $`\mathrm{\Omega }_\nu >0.1`$. However, the hardest constraint for the parameter $`\mathrm{\Omega }_\nu `$ can be got when we confront the amplitude of the first acoustic peak in $`\mathrm{\Delta }T/T`$ ($`\mathrm{}200`$) with the observational data. We compare the hight of the acoustic peak generated in our models with its value in the standard CDM (without gravitational waves) normalized by the COBE data. The parameter for such a comparison is the relative amplitude of the peak, $`\mathrm{}\mathrm{}_{\mathrm{}=200}/1.1\times 10^{10}`$, where $`\mathrm{}_{\mathrm{}}\mathrm{}(\mathrm{}+1)S_{\mathrm{}}/(\mathrm{}+0.5)`$; $`\mathrm{}=5.1`$ for sCDM. Evidently, $`\mathrm{}`$ decreases with T/S growing (and other parameters fixed). E.g. for CDM models ($`\mathrm{\Omega }_\nu =0`$ and ’standard’ values for $`n,\mathrm{\Omega }_b`$ and $`h`$) the relative amplitude of the acoustic peak decays by a factor T/S$`+14`$. The more efficient ways to enhance the acoustic peak is a transition to the ’red’ power spectra and/or high abundance of the hot matter. The role of the ’blue’ spectra becomes important when the parameter $`\mathrm{\Omega }_\nu `$ rises up (since the ’red’ spectra will violate the condition T/S$`0`$). The results of the derivation are presented in Fig.2. The condition for a ’considerable’ acoustic peak ($`\mathrm{}4`$) with the standard BBN constraint for the baryonic density, leaves us with a broad set of the power spectra ($`n[0.9,1.2]`$) but requires high fraction of the hot matter ($`\mathrm{\Omega }_\nu [0.2,0.4]`$) in the class of the models considered. ## 4 Conclusions * The data on the galaxy cluster abundance determine the value $`\sigma _8`$ with a high accuracy, the other parameters ($`n,\mathrm{\Omega }_\nu ,\mathrm{\Omega }_b,h`$) remain free within their ranges. * None of the MDM models with $`n=1`$ and T/S$`=0`$ satisfies both normalizations, on the galaxy cluster abundance and large-scale $`\mathrm{\Delta }T/T`$ anisotropy, which leads either to rejection from the flat spectrum or to the introduction of a non-zero T/S (or both). * Small values of T/S are realised for the red spectra ($`n<1`$) and moderate $`h(<0.6)`$, the violation of these conditions leads to a high T/S($`{}_{}{}^{>}1`$). * Increasing $`\mathrm{\Omega }_\nu `$ weaken the requirement to the value of T/S, however even for $`\mathrm{\Omega }_\nu 0.4`$ the models with $`h+n1.5`$ suggest considerable abundance of gravitational waves: T/S$`{}_{}{}^{>}0.3`$. * In models with $`\mathrm{\Omega }_\nu 0.4`$ and scale-invariant spectrum of density perturbations ($`n=1`$): T/S$`{}_{}{}^{>}10(h0.47)`$. * In models with $`\mathrm{\Omega }_b=0.05`$ and $`h=0.5`$ we have the following approximation for the primordial gravitational waves (the accuracy is better than 11 % for $`0.1`$T/S$`3`$): $$\frac{\mathrm{T}}{\mathrm{S}}=\frac{30(n0.7)^2}{10\mathrm{\Omega }_\nu +1}+10\mathrm{\Omega }_\nu \left(n^{3/2}1.06\right).$$ * In double-normalized models with T/S$`>`$0 the hight of the acoustic peak is less than its ‘standard’ value ($`\mathrm{}=5.1`$). The deacrease of the parameter $`\mathrm{}`$ does not exceed 30% in models with large $`\mathrm{\Omega }_\nu [0.2,0.4]`$ and any spectrum slope, $`n[0.9,1.2]`$. Any condition, $`n<0.9`$ or $`\mathrm{\Omega }_\nu <0.2`$, decreases the relative amplitude of the first acoustic peak for more than 30$`\%`$ (i.e. $`\mathrm{}<3.5`$ in models with $`\mathrm{\Omega }_b=0.05`$, $`h=0.5`$). The acoustic peak practically disappears in CDM models. * When increasing the baryonic abundance the difference between $`\mathrm{}`$ in our models and that in sCDM decreases. The amplitude of the acoustic peak coincides with its ’standard’ value ($`\mathrm{}4.5`$) in models with $`\mathrm{\Omega }_b=0.1`$ and either ’blue’ spectrum $`n[1,1.2]`$ and $`\mathrm{\Omega }_\nu 0.3`$, or moderate ’red’ spectrum $`n[0.9,1]`$ and $`\mathrm{\Omega }_\nu 0.2`$. * Thus, rising the parameter $`\mathrm{\Omega }_\nu `$ up to the values in the interval $`[0.2,0.4]`$ solves the problem of the first acoustic peak in $`\mathrm{\Delta }T/T`$, leaving the baryonic density within the primordial nucleosynthesis constraints. ## Acknowledgments The work was partly supported by Swiss National Foundation (SNSF 7IP 050163.96/1), INTAS grant (97-1192), and Russian Foundation “Development and Support of Radioastronomy Scientific and Educational Center” (N 315). E.V.M., V.N.L. and N.A.A. are grateful to the Organizing Commeittee for the hospitality. ## References
warning/0004/math0004157.html	ar5iv	text	# Mirror symmetry for concavex vector bundles on projective spaces ## 1. Introduction Let $`V^+=_{iI}𝒪(k_i)`$ and $`V^{}=_{jJ}𝒪(l_j)`$ be vector bundles on $`^s`$ with $`k_i`$ and $`l_j`$ positive integers. Suppose $`X\stackrel{\iota }{}^s`$ is the zero locus of a generic section of $`V^+`$ and $`Y`$ is a projective manifold such that $`X\stackrel{j}{}Y`$ with normal bundle $`𝒩_{X/Y}=\iota ^{}(V^{})`$. The relations between Gromov-Witten theories of $`X`$ and $`Y`$ are studied here by means of a suitably defined equivariant Gromov-Witten theory in $`^s`$. We apply mirror symmetry to the latter to evaluate the gravitational descendants of $`Y`$ supported in $`X`$. Section $`2`$ is a collection of definitions and techniques that will be used throughout this paper. In section $`3`$, using an idea from Kontsevich, we introduce a modified equivariant Gromov-Witten theory in $`^s`$ corresponding to $`V=V^+V^{}`$. The corresponding $`𝒟`$-module structure (,,) is computed in section $`4`$. It is generated by a single function $`\stackrel{~}{J}_V`$. In general, the equivariant quantum product does not have a nonequivariant limit. It is shown in Lemma $`\mathrm{4.1.1}`$ that the generator $`\stackrel{~}{J}_V`$ does have a limit $`J_V`$ which takes values in $`H^{}^m[[q,t]]`$. It is this limit that plays a crucial role in this work. Let $`Y`$ be a smooth, projective manifold. The generator $`J_Y`$ of the pure $`𝒟`$-module structure of $`Y`$ encodes one-pointed gravitational descendents of $`Y`$. It takes values in the completion of $`H^{}Y`$ along the semigroup (Mori cone) of the rational curves of $`Y`$. The pullback map $`j^{}:H^{}YH^{}X`$ extends to a map between the respective completions. In Theorem 4.2.2 we describe one aspect of the relation between pure Gromov-Witten theory of $`X\stackrel{j}{}Y`$ and the modified Gromov-Witten theory of $`^s`$. Under natural restrictions, the pull back $`j^{}(J_Y)`$ pushes forward to $`J_V`$. It follows that although defined on $`^s`$, $`J_V`$ encodes the gravitational descendants of $`Y`$ supported in $`X`$, hence the contribution of $`X`$ to the Gromov-Witten invariants of $`Y`$. The only way that $`X`$ remembers the ambient variety $`Y`$ in this context is by the normal bundle. $`Y`$ can therefore be substituted by a local manifold. This suggests that there should be a local version of mirror symmetry (see the Remark at the end of section $`4`$). This was first realized by Katz, Klemm, and Vafa . The principle of local mirror symmetry in general has yet to be understood. Some interesting calculations that contribute toward this goal can be found in . In section $`5`$ we give a proof of the Mirror Theorem which allows us to compute $`J_V`$. A hypergeometric series $`I_V`$ that corresponds to the total space of $`V`$ is defined. The Mirror Theorem 5.0.1 states that $`I_V=J_V`$ up to a change of variables. Hence, the gravitational descendants of $`Y`$ supported on $`X`$ can be computed in $`^s`$. Two examples of local Calabi-Yau threefolds are considered in section $`6`$. For $`X=^1`$ and $`V=𝒪(1)𝒪(1)`$, we obtain the Aspinwall-Morrison formula for multiple covers. If $`X=^2`$ and $`V=𝒪(3)`$, the quantum product of $`Y`$ pulls back to the modified quantum product in $`^2`$. The mirror theorem in this case yields the virtual number of plane curves on a Calabi-Yau threefold. The rich history of mirror symmetry started in 1990 with a surprising conjecture by Candelas, de la Ossa, Green and Parkes () which predicts the number $`n_d`$ of degree $`d`$ rational curves on a quintic threefold. In , Givental presented a clever argument which, as shown later by Bini et al. in and Pandharipande in , yields a proof of the mirror conjecture for Fano and Calabi-Yau (convex) complete intersections in projective spaces. Meanwhile, in a very well written paper , Lian, Liu and Yau used a different approach to obtain a complete proof of mirror theorem for concavex complete intersections on projective spaces. An alternative proof of the convex Mirror Theorem has been given by Bertram (). In this paper we use Givental’s approach to study the local nature of mirror symmetry and to present a proof of the concavex Mirror Theorem. Acknowledgements. This work is part of author’s Ph.D. thesis at Oklahoma State University. The author would like to thank Sheldon Katz for his passionate and tireless work in advising with this project. Special thanks to Bumsig Kim, Carel Faber, Ionut-Ciocan Fontanine and Zhenbo Qin who were very helpful throughout this work. At various times the author has benefited from conversations with Tom Graber, Rahul Pandharipande and Ravi Vakil, to whom the author is very grateful. We would like to thank also the referees whose help in improving this manuscript was invaluable. ## 2. Stable maps and localization ### 2.1. Genus zero stable maps. Let $`\overline{M}_{0,n}(X,\beta )`$ be the Deligne-Mumford moduli stack of pointed stable maps to $`X`$. For an excellent reference on the construction and its properties we refer the reader to . We recall some of the features on $`\overline{M}_{0,n}(X,\beta )`$ and establish some notation. For each marking point $`x_i`$ let $`e_i:\overline{M}_{0,n}(X,\beta )X`$ be the evaluation map at $`x_i`$ and $`_i`$ the cotangent line bundle at $`x_i`$. The fiber of this line bundle over a moduli point $`(C,x_1,\mathrm{},x_n,f)`$ is the cotangent space of the curve $`C`$ at $`x_i`$. Let $`\pi _k:\overline{M}_{0,n}(X,\beta )\overline{M}_{0,n1}(X,\beta )`$ be the morphism that forgets the $`k`$-th marked point. The obstruction theory of the moduli stack $`\overline{M}_{0,n}(X,\beta )`$ is described locally by the following exact sequence $$0\text{Ext}^0(\mathrm{\Omega }_C(\underset{i=1}{\overset{n}{}}x_i),𝒪_C)\text{H}^0(C,f^{}TX)𝒯_M$$ (1) $$\text{Ext}^1(\mathrm{\Omega }_C(\underset{i=1}{\overset{n}{}}x_i),𝒪_C)\text{H}^1(C,f^{}TX)\mathrm{{\rm Y}}0.$$ (Here and thereafter we are naming sheaves after their fibres). To understand the geometry behind this exact sequence we note that $`𝒯_M=\text{Ext}^1(f^{}\mathrm{\Omega }_X\mathrm{\Omega }_C,𝒪_C)`$ and $`\mathrm{{\rm Y}}=\text{Ext}^2(f^{}\mathrm{\Omega }_X\mathrm{\Omega }_C,𝒪_C)`$ are respectively the tangent space and the obstruction space at the moduli point $`(C,x_1,\mathrm{},x_n,f)`$. The spaces $`\text{Ext}^0(\mathrm{\Omega }_C(_{i=1}^nx_i),𝒪_C)`$ and $`\text{Ext}^1(\mathrm{\Omega }_C(_{i=1}^nx_i),𝒪_C)`$ describe respectively the infinitesimal automorphisms and infinitesimal deformations of the marked source curve. It follows that the expected dimension of $`\overline{M}_{0,n}(X,\beta )`$ is $`K_X\beta +\text{dim}X+n3`$. A smooth projective manifold $`X`$ is called convex if $`H^1(^1,f^{}TX)=0`$ for any morphism $`f:^1X`$. For a convex $`X`$, the obstruction bundle $`\mathrm{{\rm Y}}`$ vanishes and the moduli stack is unobstructed and of the expected dimension. Examples of convex varieties are homogeneous spaces $`G/P`$. In general this moduli stack may behave badly and have components of larger dimensions. In this case, a Chow homology class of the expected dimension has been constructed . It is called the virtual fundamental class and denoted by $`[\overline{M}_{0,n}(X,\beta )]^{\text{virt}}`$. Although its construction is quite involved, we will be using mainly two relatively easy properties. The virtual fundamental class is preserved when pulled back by the forgetful map $`\pi _n`$. A proof of this fact can be found in section $`\mathrm{7.1.5}`$ of . If the obstruction sheaf $`\mathrm{{\rm Y}}`$ is free, the virtual fundamental class refines the top Chern class of $`\mathrm{{\rm Y}}`$. This fact is proven in Proposition $`5.6`$ of . ### 2.2. Equivariant cohomology and localization theorem The notion of equivariant cohomology and the localization theorem is valid for any compact connected Lie group. For a detailed exposition on this subject we suggest Chapter $`9`$ of . Below we state without proof the results that will be used in this work. The complex torus $`T=(^{})^{s+1}`$ is classified by the principal $`T`$-bundle (2) $$ET=(^{\mathrm{}+1}\{0\})^{s+1}BT=(^{\mathrm{}})^{s+1}.$$ Let $`\lambda _i=c_1(\pi _i^{}(𝒪(1)))`$ and $`\lambda :=(\lambda _0,\mathrm{},\lambda _s)`$. We will use $`𝒪(\lambda _i)`$ for the line bundle $`\pi _i^{}(𝒪(1))`$. Clearly $`H^{}(BT)=[\lambda ]`$. If $`T`$ acts on a variety $`X`$, we let $`X_T:=X\times _TET.`$ ###### Definition 2.2.1. The equivariant cohomology of $`X`$ is (3) $$H_T^{}(X):=H^{}(X_T).$$ If $`X=x`$ is a point then $`X_T=BT`$ and $`H_T^{}(x)=[\lambda ]`$. For an arbitrary $`X`$, the equivariant cohomology $`H_T^{}(X)`$ is a $`[\lambda ]`$-module via the equivariant morphism $`Xx`$. Let $`𝒰`$ be a vector bundle over $`X`$. If the action of $`T`$ on $`X`$ can be lifted to an action on $`𝒰`$ which is linear on the fibers, $`𝒰`$ is an equivariant vector bundle and $`𝒰_T`$ is a vector bundle over $`X_T`$. The equivariant chern classes of $`E`$ are $`c_k^T(𝒰):=c_k(𝒰_T)`$. We will use $`\text{E}(𝒰)`$ ($`\text{E}_T(𝒰)`$) to denote the nonequivariant (equivariant) top chern class of $`𝒰`$. Let $`X^T=_{jJ}X_j`$ be the decomposition of the fixed point locus into its connected components. $`X_j`$ is smooth for all $`j`$ and the normal bundle $`N_j`$ of $`X_j`$ in $`X`$ is equivariant. Let $`i_j:X_jX`$ be the inclusion. The following corollary of the localization theorem will be used extensively here: ###### Theorem 2.2.1. Let $`\alpha H_T^{}(X)(\lambda )`$. Then (4) $$_{X_T}\alpha =\underset{jJ}{}_{(X_j)_T}\frac{i_j^{}(\alpha )}{\text{E}_T(N_j)}.$$ A basis for the characters of the torus is given by $`\epsilon _i(t_0,\mathrm{},t_s)=t_i`$. There is an isomorphism between the character group of the torus and $`H^2(BT)`$ sending $`\epsilon _i`$ to $`\lambda _i`$. We will say that the weight of the character $`\epsilon _i`$ is $`\lambda _i`$. For an equivariant vector bundle $`𝒰`$ over $`X`$ it may happen that the restriction of $`𝒰`$ on a fixed point component $`X_j`$ is trivial (for example if $`X_j`$ is an isolated point). In that case $`𝒰`$ decomposes as a direct sum $`_{i=1}^m\mu _i`$ of characters of the torus. If the weight of $`\mu _i`$ is $`\rho _i`$, then the restriction of $`c_k^T(𝒰)`$ on $`X_j`$ is the symmetric polynomial $`\sigma _k(\rho _1,\mathrm{},\rho _m)`$. Our interest here is for $`X=^s`$. For any action of $`T`$ on $`^s`$ we will denote (5) $`𝒫:=H_T^{}^s`$ (6) $`=:𝒫(\lambda )`$ Consider the diagonal action of $`T=(^{})^{s+1}`$ on $`^s`$ with weights $`(\lambda _0,\mathrm{},\lambda _s)`$ i.e. (7) $$(t_0,t_1,\mathrm{},t_s)(z_0,z_1,\mathrm{},z_s)=(t_0^1z_0,\mathrm{},t_s^1z_s).$$ Then $`_T^s=(_i𝒪(\lambda _i))`$. There is an obvious lifting of the action of $`T`$ on the tautological line bundle $`𝒪(1)`$. It follows that $`𝒪(k)`$ is equivariant for all $`k`$. Let $`p=c_1^T(𝒪_^s(1))`$ be the equivariant hyperplane class. We obtain $`𝒫=[\lambda ,p]/_i(p\lambda _i)`$ and $`=(\lambda )[p]/_i(p\lambda _i)`$. The locus of the fixed points consists of points $`p_j`$ for $`j=0,1,\mathrm{},s`$ where $`p_j`$ is the point whose $`j`$-th coordinate is $`1`$ and all the other ones are $`0`$. On the level of the cohomology the map $`i_j^{}`$ sends $`p`$ to $`\lambda _j`$. A basis for $``$ as a $`(\lambda )`$-vector space is given by $`\varphi _j=_{kj}(p\lambda _k)`$ for $`j=0,1,\mathrm{},s`$. Also $`i_j^{}(\varphi _j)=_{kj}(\lambda _j\lambda _k)=\text{Euler}_T(N_j)`$. The localization theorem $`\mathrm{2.2.1}`$ says that for any polynomial $`F(p)(\lambda )[p]/_{i=1}^s(p\lambda _i)`$ (8) $$_{_T^s}F(p)=\underset{j}{}\frac{F(\lambda _j)}{_{kj}(\lambda _j\lambda _k)}.$$ By translating the target of a stable map we get an action of $`T`$ on $`\overline{M}_{0,n}(^s,d)`$. In Kontsevich identified the fixed point components of this action in terms of decorated graphs. If $`f:(C,x_1,\mathrm{},x_n)^s`$ is a fixed stable map then $`f(C)`$ is a fixed curve in $`^s`$. The marked points, collapsed components and nodes are mapped to the fixed points $`p_i`$ of the $`T`$-action on $`^s`$. A noncontracted component must be mapped to a fixed line $`\overline{p_ip_j}`$ on $`^s`$. The only branch points are the two fixed points $`p_i`$ and $`p_j`$ and the restriction of the map $`f`$ to this component is determined by its degree. The graph $`\mathrm{\Gamma }`$ corresponding to the fixed point component containing such a map is constructed as follows. The vertices correspond to the connected components of $`f^1\{p_0,p_1,\mathrm{},p_s\}`$. The edges correspond to the noncontracted components of the map. The graph is decorated as follows. Edges are marked by the degree of the map on the corresponding component, and vertices are marked by the fixed point of $`^s`$ where the corresponding component is mapped to. To each vertex we associate a leg for each marked point that belongs to the corresponding component. For a vertex $`v`$, let $`n(v)`$ be the number of legs or edges incident to that vertex. Also for an edge $`e`$ let $`d_e`$ be the degree of the stable map on the corresponding component. Let (9) $$\overline{}_\mathrm{\Gamma }:=\underset{v}{}\overline{M}_{0,n(v)}.$$ There is a finite group of automorphisms $`\text{G}_\mathrm{\Gamma }`$ acting on $`\overline{M}_\mathrm{\Gamma }`$ . The order of the automorphism group $`\text{G}_\mathrm{\Gamma }`$ is (10) $$a_\mathrm{\Gamma }=\underset{e}{}d_e\|\text{Aut}(\mathrm{\Gamma })\|.$$ The fixed point component corresponding to the decorated graph $`\mathrm{\Gamma }`$ is (11) $$\overline{M}_\mathrm{\Gamma }=\overline{}_\mathrm{\Gamma }/\text{G}.$$ Let $`i_\mathrm{\Gamma }:\overline{M}_\mathrm{\Gamma }\overline{M}_{0,n}(^s,d)`$ be the inclusion of the fixed point component corresponding to $`\mathrm{\Gamma }`$ and $`N_\mathrm{\Gamma }`$ its normal bundle. This bundle is $`T`$-equivariant. Let $`\alpha `$ be an equivariant cohomology class in $`H_T^{}(\overline{M}_{0,n}(^s,d))`$ and $`\alpha _\mathrm{\Gamma }:=i_\mathrm{\Gamma }^{}(\alpha )`$. Theorem $`\mathrm{2.2.1}`$ says: (12) $$_{\overline{M}_{0,n}(^s,d)_T}\alpha =\underset{\mathrm{\Gamma }}{}_{(\overline{M}_\mathrm{\Gamma })_T}\frac{\alpha _\mathrm{\Gamma }}{a_\mathrm{\Gamma }\text{Euler}_T(N_\mathrm{\Gamma })}.$$ Explicit formulas for $`\text{Euler}_T(N_\mathrm{\Gamma })`$ in terms of chern classes of cotangent line bundles in $`H_T^{}(\overline{M}_\mathrm{\Gamma })`$ have been found by Kontsevich in . ### 2.3. Linear and nonlinear sigma models for a projective space. Two compactifications of the space of degree $`d`$ maps $`^1^s`$ will be very important in this paper. $`M_d:=\overline{M}_{0,0}(^s\times ^1,(d,1))`$ is called the degree $`d`$ nonlinear sigma model of $`^s`$ and $`N_d:=(H^0(^1,𝒪_^1(d))^{s+1})`$ is called the degree $`d`$ linear sigma model of the projective space $`^s`$. An element in $`H^0(^1,𝒪_^1(d))^{s+1}`$ is an $`s+1`$-tuple of degree $`d`$ homogeneous polynomials in two variables $`w_0`$ and $`w_1`$. As a vector space, $`H^0(^1,𝒪_^1(d))^{s+1}`$ is generated by the vectors $`v_{ir}=(0,\mathrm{},0,w_0^rw_1^{dr},0\mathrm{},0)`$ for $`i=0,1,\mathrm{},s`$ and $`r=0,1,\mathrm{},d`$. The only nonzero component of $`v_{ir}`$ is the $`i`$-th one. The action of $`T^{}:=T\times ^{}`$ in $`^s\times ^1`$ with weights $`(\lambda _0,\mathrm{},\lambda _s)`$ in the $`^s`$ factor and $`(\mathrm{},0)`$ in the $`^1`$ factor gives rise to an action of $`T^{}`$ in $`M_d`$ by translation of maps. $`T^{}`$ also acts in $`N_d`$ as follows: for $`\overline{t}=(t_0,\mathrm{},t_s)T`$ and $`t^{}`$ (13) $$(\overline{t},t)[P_0(w_0,w_1),\mathrm{},P_s(w_0,w_1)]=[t_0P_0(tw_0,w_1),\mathrm{},t_sP_s(tw_0,w_1)]$$ There is a $`T^{}`$-equivariant morphism $`\psi :M_dN_d.`$ Here is a set-theoretical description of this map (for a proof that it is a morphism see or ). Let $`q_i`$ for $`i=1,2`$ be the projection maps on $`^s\times ^1`$. For a stable map $`(C,f)M_d`$ let $`C_0`$ be the unique component of $`C`$ such that $`q_2f:C_0^1`$ is an isomorphism. Let $`C_1,\mathrm{},C_n`$ be the irreducible components of $`CC_0`$ and $`d_i`$ the degree of the restriction of $`q_1f`$ on $`C_i`$. Choose coordinates on $`C_0^1`$ such that $`q_2f(y_0,y_1)=(y_1,y_0)`$. Let $`C_0C_i=(a_i,b_i)`$ and $`q_1f=[f_0:f_1:\mathrm{}:f_s]:C_0^s`$. Then (14) $$\psi (C,f):=\underset{i=1}{\overset{n}{}}(b_iw_0a_iw_1)^{d_i}[f_0:f_1:\mathrm{}:f_s].$$ Let $`p_{ir}`$ be the points of $`N_d`$ corresponding to the vectors $`v_{ir}`$. The fixed point loci of the $`T^{}`$-action on $`N_d`$ consists of the points $`p_{ir}`$. We write $`\kappa `$ for the equivariant hyperplane class of $`N_d`$. The restriction of $`\kappa `$ at the fixed point $`p_{ir}`$ is $`\lambda _i+r\mathrm{}`$. The restriction of the equivariant Euler class of the tangent space $`TN_d`$ at $`p_{ir}`$ is (15) $$E_{ir}=\underset{(j,t)(i,r)}{}(\lambda _i\lambda _j+r\mathrm{}t\mathrm{}).$$ Fixed point components of $`M_d`$ are obtained as follows. Let $`\mathrm{\Gamma }_{d_j}^i`$ be the graph of a $`T`$-fixed point component in $`\overline{M}_{0,1}(^s,d_j)`$ where the marking is mapped to $`p_i`$ and $`d_1+d_2=d`$. Let $`(d_1,d_2)`$ be a partition of $`d`$. We identify $`\overline{M}_{\mathrm{\Gamma }_{d_1}^i}\times \overline{M}_{\mathrm{\Gamma }_{d_2}^i}`$ with a $`T^{}`$-fixed point component $`M_{d_1d_2}^i`$ in $`M_d`$ in the following manner. Let $`(C_1,x_1,f_1)\overline{M}_{\mathrm{\Gamma }_{d_1}^i}`$ and $`(C_2,x_2,f_2)\overline{M}_{\mathrm{\Gamma }_{d_2}^i}`$. Let $`C`$ be the nodal curve obtained by gluing $`C_1`$ with $`^1`$ at $`x_1`$ and $`0^1`$ and $`C_2`$ with $`^1`$ at $`x_2`$ and $`\mathrm{}^1`$. Let $`f:C^s\times ^1`$ map $`C_1`$ to the slice $`^s\times \mathrm{}`$ by means of $`f_1`$ and $`C_2`$ to $`^s\times 0`$ by means of $`f_2`$. Finally $`f`$ maps $`^1`$ to $`p_i\times ^1`$ by permuting coordinates. $`\psi `$ maps $`M_{d_1d_2}^i`$ to $`p_{id_2}N_d`$, hence the equivariant restriction of $`\psi ^{}(\kappa )`$ in $`M_{d_1d_2}^i`$ is $`\lambda _i+d_2\mathrm{}`$. The normal bundle $`N_{\mathrm{\Gamma }_{d_1d_2}^i}`$ of this component in the above identification can be found by splitting it in five pieces: smoothing the nodes $`x_1`$ and $`x_2`$ and deforming the restriction of the map to $`C_1,C_2,^1`$. Using Kontsevich’s calculations, Givental obtained $$\frac{1}{\text{E}_T(N_{\mathrm{\Gamma }_{d_1d_2}^i})}=\frac{1}{_{ki}(\lambda _i\lambda _k)}\frac{1}{\text{E}_T(N_{\mathrm{\Gamma }_{d_1}^i})}\frac{1}{\text{E}_T(N_{\mathrm{\Gamma }_{d_2}^i})}\frac{e_1^{}(\varphi _i)}{\mathrm{}(\mathrm{}c_1)}\frac{e_1^{}(\varphi _i)}{\mathrm{}(\mathrm{}c_2)}$$ where $`c_j,j=1,2`$ is the first Chern class of the cotangent line bundle on $`\overline{M}_{\mathrm{\Gamma }_{d_j}^i}`$. ## 3. A Gromov-Witten theory induced by a vector bundle ### 3.1. The obstruction class of a concavex vector bundle. The notion of concavex vector bundle is due to Lian, Liu and Yau and is central to this work. ###### Definition 3.1.1. 1. A line bundle $``$ on $`X`$ is called convex if $`H^1(C,f^{}())=0`$ for any genus zero stable map $`(C,x_1,\mathrm{}x_n,f)`$. 2. A line bundle $``$ on $`X`$ is called concave if $`H^0(C,f^{}())=0`$ for any nonconstant genus zero stable map $`(C,x_1,\mathrm{}x_n,f)`$. 3. A direct sum of convex and concave line bundles on $`X`$ is called a concavex vector bundle. A concavex vector bundle $`V`$ in a projective space $`^s`$ has the form (16) $$V=V^+V^{}=\left(_{iI}𝒪(k_i)\right)\left(_{jJ}𝒪(l_j)\right)$$ where $`k_i`$ and $`l_j`$ are positive numbers. Denote $`𝔼^+:=\text{E}(V^+)`$ and $`𝔼^{}:=\text{E}(V^{})`$. Let $`d>0`$. Consider the following diagram $$\begin{array}{ccc}\overline{M}_{0,n+1}(^s,d)& \stackrel{e_{n+1}}{}& ^s\\ \pi _{n+1}& & \\ \overline{M}_{0,n}(^s,d)\end{array}$$ Since $`V`$ is concavex, the sheaf (17) $$V_d:=V_d^+V_d^{}=\pi _{n+1}^{}{}_{}{}^{}e_{n+1}^{}(V^+)R^1\pi _{n+1}^{}{}_{}{}^{}e_{n+1}^{}(V^{})$$ is locally free. ###### Definition 3.1.2. The obstruction class corresponding to $`V`$ is defined to be: (18) $$𝔼_d:=\text{E}(V_d)=\text{E}(V_d^+)\text{E}(V_d^{}):=𝔼_d^+𝔼_d^{}$$ For a $`T`$-action on $`^s`$ that lifts to a linear action on the fibers of $`V=V^+V^{}`$, let $`E^+:=\text{E}_T(V^+)`$ and $`E^{}:=\text{E}_T(V^{})`$. Assume that $`E^{}`$ is invertible. ###### Definition 3.1.3. The modified equivariant integral $`\omega _V:(\lambda )`$ corresponding to $`V`$ is defined as follows (19) $$\omega _V(\alpha ):=_{_T^m}\alpha \frac{E^+}{E^{}}.$$ Consider the trivial action of $`T=(_{iI}C^{})(_{jJ}C^{})`$ on $`^s`$. In this case $`_T^s=^s\times (_{iI}^{\mathrm{}})\times (_{jJ}^{\mathrm{}})`$ and $`\overline{M}_{0,n}(^s,d)_T=\overline{M}_{0,n}(^s,d)\times (_{iI}^{\mathrm{}})\times (_{jJ}^{\mathrm{}})`$. It follows that $`𝒫=H^{}(^s,[\lambda ])`$ and $`=H^{}(^s,(\lambda ))`$. Let $`p`$ denotes the equivariant hyperplane class. The $`T`$-action lifts to a linear action on the fibers of $`V`$ with weights $`((\lambda _i)_{iI},(\lambda _j)_{jJ})`$. Let $`q_i`$ and $`q_j`$ denote the projection maps on $`\overline{M}_{0,n}(^s,d)_T`$. Both $`V_d^+`$ and $`V_d^{}`$ are $`T`$-equivariant bundles and $`(V_d^+)_T=V_d^+(_{iI}q_i^{}𝒪_{^{\mathrm{}}}(\lambda _i))`$ (20) $`(V_d^{})_T=V_d^{}(_{jJ}q_j^{}𝒪_{^{\mathrm{}}}(\lambda _j))`$ The equivariant obstruction class is (21) $$E_d:=\text{E}_T(V_d)=\text{E}_T(V_d^+)\text{E}_T(V_d^{})=E_d^+E_d^{}.$$ The modified equivariant integral for the trivial action of $`T`$ on $`^s`$ gives rise to a modified perfect pairing in $``$ $$a,b_V:=\omega _V(ab).$$ Let $`T_0=1,T_1=p,\mathrm{},T^s=p^s`$ be a basis of $``$ as a $`(\lambda )`$-vector space. The intersection matrix $`(g_{rt}):=(T_r,T_t_V)`$ has an inverse $`(g^{rt})`$. Let $`T^i=_{j=0}^sg^{ij}T_j`$ be the dual basis with respect to this pairing. Clearly (22) $$T^i=T_{mi}\left(\frac{E^{}}{E^+}\right).$$ This implies that in $`H^{}(^s\times ^s)(\lambda )`$ we have (23) $$\underset{i=1}{\overset{s}{}}T_iT^i=\mathrm{\Delta }\left(1\frac{E^{}}{E^+}\right)$$ where $`\mathrm{\Delta }=_{i=0}^sT_iT_{si}`$ is the class of the diagonal in $`^s\times ^s`$. Recall that the morphism $`\pi _k:\overline{M}_{0,n}(^s,d)\overline{M}_{0,n1}(^s,d)`$ forgets the $`k`$-th marked point. ###### Lemma 3.1.1. $`\pi _{k}^{}{}_{}{}^{}(E_d)=E_d`$ and $`\pi _{k}^{}{}_{}{}^{}(𝔼_d)=𝔼_d.`$ Proof. For simplicity we will consider the case $`V=𝒪(k)𝒪(l)`$ and $`k=n`$. The general case is similar. Let $`M_k=\overline{M}_{0,k}(^s,d)`$ and $`M_{n,n}=M_n\times _{M_{n1}}M_n`$. Consider the following equivariant commutative diagram: We compute: (24) $$\pi _{n+1}^{}{}_{}{}^{}e_{n+1}^{}𝒪(k)=\pi _{n+1}^{}{}_{}{}^{}\pi _n^{}e_{n}^{}{}_{}{}^{}𝒪(k)=\beta _{}\mu _{}\mu ^{}\alpha ^{}e_{n}^{}{}_{}{}^{}𝒪(k).$$ By the projection formula (25) $$\mu _{}\mu ^{}\alpha ^{}e_{n}^{}{}_{}{}^{}𝒪(k)=\alpha ^{}e_{n}^{}{}_{}{}^{}𝒪(k)\mu _{}(𝒪_{M_{n+1}}).$$ Since the map $`\mu `$ is birational and $`M_{n+1}`$ is normal $`\mu _{}(𝒪_{M_{n+1}})=𝒪_{M_{n,n}}`$ hence (26) $$\mu _{}\mu ^{}\alpha ^{}e_{n}^{}{}_{}{}^{}𝒪(k)=\alpha ^{}e_{n}^{}{}_{}{}^{}𝒪(k).$$ Substituting in (24) and applying base extension properties ($`\pi _n`$ is flat) yields (27) $$\pi _{n+1}^{}{}_{}{}^{}e_{n+1}^{}𝒪(k)=\beta _{}\alpha ^{}e_{n}^{}{}_{}{}^{}𝒪(k)=\pi _{n}^{}{}_{}{}^{}(\pi _{n}^{}{}_{}{}^{}e_{n}^{}{}_{}{}^{}𝒪(k)).$$ For the case of a negative line bundle we have (28) $$R^1\pi _{n+1}^{}{}_{}{}^{}e_{n+1}^{}𝒪(l)=R^1\pi _{n+1}^{}{}_{}{}^{}\pi _{n}^{}{}_{}{}^{}e_{n}^{}{}_{}{}^{}𝒪(l)=R^1\pi _{n+1}^{}{}_{}{}^{}\mu ^{}\alpha ^{}e_{n}^{}{}_{}{}^{}𝒪(l).$$ We now use the spectral sequence (29) $$R^p\beta _{}(R^q\mu _{})R^{p+q}\pi _{n+1}^{}{}_{}{}^{}$$ where $``$ is a sheaf of $`𝒪_{_{n+1}}`$-modules. The map $`\mu `$ is birational. If we think of $`M_n`$ as the universal map of $`M_{n1}`$, then the map $`\mu `$ has nontrivial fibers only over pairs of stable maps in $`M_n`$ that represent the same special point (i.e. node or marked point) of a stable map in $`M_{n1}`$. These nontrivial fibers are isomorphic to $`^1`$. Since $`=e_{n+1}^{}𝒪(l)`$ we obtain $`R^q\mu _{}=0`$ for $`q>0`$. It follows that this spectral sequence degenerates, giving (30) $$R^1\pi _{n+1}^{}{}_{}{}^{}e_{n+1}^{}𝒪(l)=R^1\beta _{}\mu _{}\mu ^{}\alpha ^{}e_{n}^{}{}_{}{}^{}𝒪(l).$$ Now we proceed as in (26) to conclude (31) $$R^1\pi _{n+1}^{}{}_{}{}^{}e_{n+1}^{}𝒪(l)=\pi _{n}^{}{}_{}{}^{}(R^1\pi _{n}^{}{}_{}{}^{}e_{n}^{}{}_{}{}^{}𝒪(l)).$$ The lemma is proven.$``$ ###### Remark 3.1.1. The previous lemma justifies the omission of $`n`$ from the notation of the obstruction class. ### 3.2. Modified equivariant correlators and quantum cohomology. Let $`\gamma _i`$ for $`i=1,\mathrm{},n`$ and $`d>0`$. Introduce the following modified equivariant Gromow-Witten invariants: (32) $$\stackrel{~}{I_d}(\gamma _1,\mathrm{},\gamma _n):=_{\overline{M}_{0,n}(^m,d)_T}e_1^{}(\gamma _1)\mathrm{}e_n^{}(\gamma _n)E_d(\lambda )$$ Now $`\overline{M}_{0,n}(^s,0)=\overline{M}_{0,n}\times ^s`$ and all the evaluation maps equal the projection $`q_2`$ to the second factor. The integrals in this case are defined as follows (33) $$\stackrel{~}{I_0}(\gamma _1,\mathrm{},\gamma _n):=_{\overline{M}_{0,n}(^s,0)}e_1^{}(\gamma _1)\mathrm{}e_n^{}(\gamma _n)q_2^{}(\text{E}(V))(\lambda )$$ The modified equivariant gravitational descendants are defined similarly to Gromov-Witten invariants: (34) $$\stackrel{~}{I}_d(\tau _{k_1}\gamma _1,\mathrm{},\tau _{k_n}\gamma _n):=_{\overline{M}_{0,n}(^s,d)_T}c_1^{k_1}(_1)e_1^{}(\gamma _1)\mathrm{}c_1^{k_n}(_n)e_n^{}(\gamma _n)E_d.$$ Lemma (3.1.1) is essential in proving that the modified correlators satisfy the same properties, such as fundamental class property, divisor property, point mapping axiom etc., that the usual Gromov-Witten invariants do. The proofs are similar to the ones in pure Gromov-Witten theory. As an illustration, we prove one of these properties. Fundamental class property. Let $`\gamma _n=1`$ and $`d0`$. The forgetful morphism $`\pi _n:\overline{M}_{0,n}(^s,d)\overline{M}_{0,n1}(^s,d)`$ is equivariant. Using Lemma 3.1.1 we obtain: $$e_1^{}(\gamma _1)\mathrm{}e_{n1}^{}(\gamma _{n1})e_n^{}(1)E_d=\pi ^{}(e_1^{}(\gamma _1)\mathrm{}e_{n1}^{}(\gamma _{n1})E_d).$$ Therefore: $$\stackrel{~}{I_d}(\gamma _1,\mathrm{},\gamma _{n1},1)=_{\overline{M}_{0,n}(^s,d)}\pi ^{}(e_1^{}(\gamma _1)\mathrm{}e_{n1}^{}(\gamma _{n1})E_d)=$$ $$=_{\pi _{n}^{}{}_{}{}^{}(\overline{M}_{0,n}(^s,d))}e_1^{}(\gamma _1)\mathrm{}e_{n1}^{}(\gamma _{n1})E_d=0.$$ The last equality is because the fibers of $`\pi _n`$ are positive dimensional. If $`d=0`$, by the point mapping property we know that the integral is zero unless $`n=3`$. In that case: $`\stackrel{~}{I_0}(\gamma _1,\gamma _2,1)=\gamma _1,\gamma _2`$.$``$ We will now prove a technical lemma which will be very useful later. Let $`AB`$ be a partition of the set of markings and $`d=d_1+d_2`$. Let $`D=D(A,B,d_1,d_2)`$ be the closure in $`\overline{M}_{0,n}(^s,d)`$ of stable maps of the following type. The source curve is a union $`C=C_1C_2`$ of two lines meeting at a node $`x`$. The marked points corresponding to $`A`$ are on $`C_1`$ and those corresponding to $`B`$ are on $`C_2`$. The restriction of the map $`f`$ on $`C_i`$ has degree $`d_i`$ for $`i=1,2`$. $`D`$ is a boundary divisor in $`\overline{M}_{0,n}(^s,d)`$. Let $`M_1:=\overline{M}_{0,\|A\|+1}(^s,d_1)`$ and $`M_2:=\overline{M}_{0,\|B\|+1}(^s,d_2)`$. Let $`e_x`$ and $`\stackrel{~}{e_x}`$ be the evaluation maps at the additional marking in $`M_1`$ and $`M_2`$ and $`\mu :=(e_x,\stackrel{~}{e_x})`$. The boundary divisor $`D`$ is obtained from the following fibre diagram $$\begin{array}{ccc}D& \stackrel{\iota }{}& M_1\times M_2\\ \nu & & \mu & & \\ ^s& \stackrel{\delta }{}& ^s\times ^s\end{array}$$ where $`\nu `$ is the “evaluation map at the node $`x`$” and $`\delta `$ is the diagonal map. ###### Lemma 3.2.1. For any classes $`\gamma _1,\mathrm{},\gamma _n`$ in $``$: $$_D\underset{i=1}{\overset{n}{}}e_i^{}(\gamma _i)E_d=\underset{a=0}{\overset{s}{}}\left(_{M_1}\underset{iA}{}e_i^{}(\gamma _i)e_x^{}(T_a)E_{d_1}\right)\times \left(_{M_2}\underset{jB}{}e_j^{}(\gamma _j)\stackrel{~}{e}_x^{}(T^a)E_{d_2}\right).$$ Proof. This lemma is the analogue of the Lemma $`16`$ in . The proof needs a minor modification. Let $`\alpha :D\overline{M}_{0,n}(^s,d)`$. Consider the normalization sequence at $`x`$: (35) $$0𝒪_C𝒪_C^{}𝒪_{C^{\prime \prime }}𝒪_x0.$$ Twisting it by $`f^{}(V^+)`$ and $`f^{}(V^{})`$ and taking the cohomology sequence yields the following identities on $`D`$: (36) $$\alpha ^{}(E_d^+)\nu ^{}(E^+)=\iota ^{}(E_{d_1}^+\times E_{d_2}^+).$$ and (37) $$\alpha ^{}(E_d^{})=\iota ^{}(E_{d_1}^{}\times E_{d_2}^{})\nu ^{}(E^{}).$$ By combining equations (36) and (37) we obtain the restriction of $`E_d`$ in the divisor $`D`$: (38) $$\alpha ^{}(E_d)=\iota ^{}(E_{d_1}\times E_{d_2})\nu ^{}\left(\frac{E^{}}{E^+}\right).$$ Using formula $`(23)`$ we obtain (39) $$\iota _{}\nu ^{}\left(\frac{E^{}}{E^+}\right)=\mu ^{}\left(1\frac{E^{}}{E^+}\right)\mu ^{}(\mathrm{\Delta }).$$ Therefore $`{\displaystyle _D}{\displaystyle \underset{i=1}{\overset{n}{}}}e_i^{}(\gamma _i)E_d={\displaystyle _{M_1\times M_2}}{\displaystyle \underset{i=1}{\overset{n}{}}}e_i^{}(\gamma _i)E_{d_1}E_{d_2}\mu ^{}\left(1{\displaystyle \frac{E^{}}{E^+}}\right)`$ $`\mu ^{}(\mathrm{\Delta })={\displaystyle _{M_1\times M_2}}{\displaystyle \underset{i=1}{\overset{n}{}}}e_i^{}(\gamma _i)E_{d_1}E_{d_2}\mu ^{}\left({\displaystyle \underset{a}{}}T_aT^a\right)=`$ (40) $`{\displaystyle \underset{a=0}{\overset{m}{}}}\left({\displaystyle _{M_1}}{\displaystyle \underset{i=1}{\overset{n_1}{}}}e_i^{}(\gamma _i)e_x^{}(T_a)E_{d_1}\right)\times \left({\displaystyle _{M_2}}{\displaystyle \underset{j=1}{\overset{n_2}{}}}e_j^{}(\gamma _j)\stackrel{~}{e}_x^{}(T^a)E_{d_2}\right).`$ The lemma is proven.$``$ The same proof can be used to show that the previous splitting lemma is true for gravitational descendants as well. ###### Corollary 3.2.1. The following modified topological recursion relations hold: $`\stackrel{~}{I}_d(\tau _{k_1+1}\gamma _1,\tau _{k_2}\gamma _2,\tau _{k_3}\gamma _3,{\displaystyle \underset{i=4}{\overset{n}{}}}\tau _{s_i}\omega _i)=`$ (41) $`{\displaystyle \stackrel{~}{I}_{d_1}(\tau _{k_1}\gamma _1,\underset{iI_1}{}\tau _{s_i}\omega _i,T_a)\stackrel{~}{I}_{d_2}(T^a,\tau _{k_2}\gamma _2,\tau _{k_3}\gamma _3,\underset{iI_2}{}\tau _{s_i}\omega _i)}`$ where the sum is over all splittings $`d_1+d_2=d`$ and partitions $`I_1I_2=\{4,\mathrm{},n\}`$ and over all indices $`a`$. Proof. Let $`A`$ and $`B`$ be two disjoint subsets of $`\{1,2,\mathrm{},n\}`$. We will denote by $`D(A,B)`$ the sum of boundary divisors $`D(E,F,d_1,d_2)`$ such that $`E,F`$ is a partition of $`\{1,2,\mathrm{},n\}`$ and $`AE`$ and $`BF`$ and $`d_1+d_2=d`$. The notation $`D(A,B)`$ does not reflect neither the number $`n`$ of marked points nor the degree $`d`$ of the maps but they will be clear from the context. Consider the morphism $`\pi :\overline{M}_{0,n}(^s,d)\overline{M}_{0,3}`$ that forgets the map and all but the first $`3`$ markings. Since $`\overline{M}_{0,3}`$ is a point, the cotangent line bundle at the first marking is trivial. But $`\pi ^{}(_1)=_1D(\{1\},\{2,3\})`$ therefore $`_1=D(\{1\},\{2,3\})`$ in $`\overline{M}_{0,n}(^s,d)`$. Multiply both sides of the previous equation by $`_{j=1}^3c_1(_j)^{k_j}e_j^{}(\gamma _j)_{i=4}^nc_1(_i)^{s_i}e_i^{}(\omega _i)E_d`$ and integrate. The corollary follows from the splitting lemma for gravitational descendents.$``$ In the process of finding solutions to the WDVV equations, Kontsevich suggested the following modified equivariant Gromov-Witten potential (42) $$\stackrel{~}{\mathrm{\Phi }}(t_0,t_1,\mathrm{},t_m):=\underset{n3}{}\underset{d0}{}\frac{1}{n!}\stackrel{~}{I_d}(\gamma ^n)$$ where $`\gamma =t_0+t_1p+\mathrm{}+t_sp^s`$ and $`t_i(\lambda )`$. Let $`\stackrel{~}{\mathrm{\Phi }}_{ijk}={\displaystyle \frac{^3\stackrel{~}{\mathrm{\Phi }}}{t_it_jt_k}}`$. ###### Definition 3.2.1. The modified, equivariant quantum product on $``$ is defined to be the linear extension of (43) $$T_i_VT_j:=\underset{k=0}{\overset{m}{}}\stackrel{~}{\mathrm{\Phi }}_{ijk}T^k.$$ ###### Theorem 3.2.1. $`QH_V^{}_T^s:=(,_V)`$ is a commutative, associative algebra with unit $`T_0`$. Proof. A simple calculation shows that: $$\stackrel{~}{\mathrm{\Phi }}_{ijk}=\underset{n0}{}\underset{d0}{}\frac{1}{n!}\stackrel{~}{I_d}(T_i,T_j,T_k,\gamma ^n).$$ The commutativity of the modified, equivariant quantum product follows from the symmetry of the new integrals. $`T_0`$ is the unit due to the fundamental class property for the modified Gromov-Witten invariants. To proving the associativity we proceed as in Theorem 4 in . Let $`\stackrel{~}{\mathrm{\Phi }}_{ijk}={\displaystyle \frac{^3\stackrel{~}{\mathrm{\Phi }}}{t_it_jt_k}}.`$ We compute $$(T_i_VT_j)_VT_k=\stackrel{~}{\mathrm{\Phi }}_{ije}g^{ef}\stackrel{~}{\mathrm{\Phi }}_{fkl}g^{ld}T_d$$ $$T_i_V(T_j_VT_k)=\stackrel{~}{\mathrm{\Phi }}_{jke}g^{ef}\stackrel{~}{\mathrm{\Phi }}_{fil}g^{ld}T_d.$$ Since the matrix $`(g^{ld})`$ is nonsingular, $`(T_i_VT_j)_VT_k=T_i_V(T_j_VT_k)`$ is equivalent to (44) $$\underset{e,f}{}\stackrel{~}{\mathrm{\Phi }}_{ije}g^{ef}\stackrel{~}{\mathrm{\Phi }}_{fkl}=\underset{e,f}{}\stackrel{~}{\mathrm{\Phi }}_{jke}g^{ef}\stackrel{~}{\mathrm{\Phi }}_{fil}.$$ Equation (44) is the WDVV equation for the modified potential $`\stackrel{~}{\mathrm{\Phi }}`$. To prove this equation let $`q,r,s,t`$ be four different integers in $`\{1,2,\mathrm{}n\}`$. There exists an equivariant morphism: $$\pi :\overline{M}_{0,n}(^s,d)\overline{M}_{0,4}=^1$$ that forgets the map and all the marked points but $`q,r,s,t`$. Obviously the divisors $`D(\{q,r\},\{s,t\})`$ and $`D(\{q,s\},\{r,t\})`$ are linearly equivalent in $`\overline{M}_{0,4}`$ hence, via the pullback $`\pi ^{}`$, they are linearly equivalent in $`\overline{M}_{0,n}(^s,d)`$. Now integrate the class $$\underset{i=1}{\overset{n4}{}}(e_i^{}(\gamma ))e_{n3}^{}(T_i)e_{n2}^{}(T_j)e_{n1}^{}(T_k)e_n^{}(T_l)E_d$$ over $`D(\{q,r\},\{s,t\})`$ and use Lemma 3.2.1 to obtain WDVV equation hence the associativity.$``$ If we restrict $`\stackrel{~}{\mathrm{\Phi }}_{ijk}`$ to the divisor classes $`\gamma =tp`$, and use the divisor property for the modified Gromov-Witten invariants, we obtain the small product: (45) $$T_i_VT_j:=T_iT_j+\underset{d>0}{}q^d\underset{k=0}{\overset{m}{}}\stackrel{~}{I_d}(T_i,T_j,T_k)T^k.$$ Here $`q=e^t`$. We extend this product to $`_{}[[q]]`$ to obtain the small equivariant quantum cohomology ring $`SQH_V^{}_{}^{s}{}_{T}{}^{}`$. We will use $`_V`$ to denote both the small and the big quantum product. The difference will be clear from the context. ###### Remark 3.2.1. * Equation (38) and Lemma 3.1.1 are the basis for building a modified equivariant Gromow-Witten theory similar to pure Gromov-Witten theory. * One can see from (22) that the only potential problem with the existence of the nonequivariant limit of (45) is the presence of $`E^+`$ in the denominator of $`T^k`$. Hence if $`V=V^{}`$ is a pure negative line bundle the nonequivariant limit of this product exists. An example of this situation is treated in the last section. ## 4. A $`𝒟`$-module structure induced by $`V`$ ### 4.1. Equivariant quantum differential equations. Recall from section $`2.3`$ the generator $`\mathrm{}`$ of $`H^2(BC^{})`$. Consider the system of first order differential equations on the modified, big quantum cohomology ring $`QH_V^{}(_T^s)`$ (46) $$\mathrm{}\frac{}{t_i}=T_i_V:i=1,\mathrm{},m.$$ ###### Theorem 4.1.1. The space of solutions of these equations has the following basis: $`s_a=T_a+{\displaystyle \underset{j=0}{\overset{s}{}}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{d=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\mathrm{}^{(k+1)}}{n!}}\stackrel{~}{I}_d(\tau _kT_a,T_j,\gamma ^n)T^j=`$ (47) $`T_a+{\displaystyle \underset{j=0}{\overset{m}{}}}{\displaystyle \underset{d=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{n!}}\stackrel{~}{I}_d({\displaystyle \frac{T_a}{\mathrm{}c}},T_j,\gamma ^n)T^j`$ where $`c`$ is a formal symbol that stands for $`c_{1}^{}{}_{}{}^{T}(_1)`$ and $`{\displaystyle \frac{T_a}{\mathrm{}c}}`$ should be expanded in powers of $`{\displaystyle \frac{c}{\mathrm{}}}`$. Proof. On one hand (48) $$\mathrm{}\frac{s_a}{t_i}=\underset{j=0}{\overset{m}{}}\underset{n=0}{\overset{\mathrm{}}{}}\underset{d=0}{\overset{\mathrm{}}{}}\underset{k=0}{\overset{\mathrm{}}{}}\frac{\mathrm{}^k}{n!}\stackrel{~}{I}_d(\tau _kT_a,T_j,T_i,\gamma ^n)T^j.$$ On the other hand $`T_is_a=T_iT_a+{\displaystyle \underset{j=0}{\overset{m}{}}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{d=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\mathrm{}^{(k+1)}}{n_1!}}\stackrel{~}{I}_d(\tau _kT_a,T_j,\gamma ^n)(T_iT^j)=`$ $`{\displaystyle \underset{n,d,e}{}}{\displaystyle \frac{1}{n!}}\stackrel{~}{I}_d(T_i,T_a,T_e,\gamma ^n)T^e+{\displaystyle \underset{j=0}{\overset{m}{}}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{d_1}{}}{\displaystyle \frac{\mathrm{}^{(k+1)}}{n!}}\stackrel{~}{I}_{d_1}(\tau _kT_a,T_j,\gamma ^n)`$ (49) $`\times {\displaystyle \underset{m,d_2,e}{}}{\displaystyle \frac{1}{m!}}\stackrel{~}{I}_{d_2}(T_i,T^j,T_e,\gamma ^m)T^e.`$ The theorem follows from the topological recursion relations (41).$``$ Restrictions $`\stackrel{~}{s}_a`$ of the sections $`s_a`$ to $`\gamma H^0(^m)H^2(^m)`$ are solutions of $$\mathrm{}\frac{}{t_i}=T_i_V:i=0,1.$$ Repeated use of the divisor axiom yields (50) $$\stackrel{~}{s}_a=e^{\frac{t_0+pt_1}{\mathrm{}}}T_a+\underset{d=1}{\overset{\mathrm{}}{}}\underset{j=0}{\overset{m}{}}q^d\stackrel{~}{I}_d(\frac{e^{\frac{t_0+pt_1}{\mathrm{}}}T_a}{\mathrm{}c},T_j)T^j,$$ where $`q:=e^{t_1}`$ ###### Definition 4.1.1. The module of differential operators that annihilate $`\stackrel{~}{s}_a,1_V`$ for all $`a`$ is called the modified equivariant $`𝒟`$-module of $`^s`$ induced by $`V`$. This module is generated by the following $`[[t_0,t_1,q]]`$-valued function (51) $$\stackrel{~}{J}_V=\underset{a=0}{\overset{s}{}}\stackrel{~}{s}_a,1_VT^a.$$ Recall that $`e_1:\overline{M}_{0,2}(^s,d)^s`$ is the evaluation map at the first marked point and $`c`$ is the chern class of the cotangent line bundle at the first marked point. Substituting (50) in (51) and using the projection formula we obtain: (52) $`\stackrel{~}{J}_V=\text{exp}\left({\displaystyle \frac{t_0+pt_1}{\mathrm{}}}\right)`$ $`\left(1+{\displaystyle \underset{d>0}{}}q^dPD^1\left(e_{1}^{}{}_{}{}^{}\left(\left({\displaystyle \frac{E_d}{\mathrm{}c}}[\overline{M}_{0,2}(^s,d)]\right)\right)\left({\displaystyle \frac{E^{}}{E^+}}\right)\right)\right).`$ In the above expression $`PD:H^{}(\overline{M}_{0,2}(^s,d))H_{s+d+sd1}\overline{M}_{0,2}(^s,d))`$ is the Poincaré duality isomorphism. It will be convenient for us to work with the moduli space of one pointed stable maps. To that end we note that (53) $$e_{1}^{}{}_{}{}^{}\left(\frac{E_d}{\mathrm{}c}[\overline{M}_{0,2}(^s,d)]\right)=e_{1}^{}{}_{}{}^{}\left(\frac{E_d}{\mathrm{}(\mathrm{}c)}[\overline{M}_{0,1}(^s,d)]\right).$$ This identity follows easily from the fact that if $`\pi _2:\overline{M}_{0,2}(^s,d)\overline{M}_{0,1}(^s,d)`$ forgets the second marked point and $`D`$ is the image of the universal section of $`\pi `$ induced by the marked point, then $`c=\pi _{2}^{}{}_{}{}^{}(c)+D`$ and $`𝔼_d=\pi _{2}^{}{}_{}{}^{}(𝔼_d)`$. The final expression for $`\stackrel{~}{J}_V`$ is (54) $`\stackrel{~}{J}_V=\text{exp}\left({\displaystyle \frac{t_0+pt_1}{\mathrm{}}}\right)`$ $`\left(1+{\displaystyle \underset{d>0}{}}q^dPD^1\left(e_{1}^{}{}_{}{}^{}\left({\displaystyle \frac{E_d}{\mathrm{}(\mathrm{}c)}}[\overline{M}_{0,1}(^s,d)]\right)\right)\left({\displaystyle \frac{E^{}}{E^+}}\right)\right).`$ From this presentation we see that the presence of the equivariant class $`E^+`$ in the denominator of $`\stackrel{~}{J}_V`$ is a potential problem for the existence of the nonequivariant limit. ###### Lemma 4.1.1. $`\stackrel{~}{J}_V𝒫[[q]]`$ therefore it has a nonequivariant limit. Proof. Let $`V_d^{}`$ be the subbundle of $`V_d^+`$ whose fiber consists of those sections of $`H^0(C,f^{}(V^+))`$ that vanish at the marked point. Let $`E_d^{}:=c_{\text{top}}^T(V_d^{})`$. There is an exact sequence of equivariant bundles on $`\overline{M}_{0,1}(^s,d)`$: $$0V_d^{}V_d^+e_{1}^{}{}_{}{}^{}(V^+)0.$$ Taking the top chern classes we obtain (55) $$E_d^+=E_d^{}e_1^{}(E^+).$$ We compute $`PD^1e_{1}^{}{}_{}{}^{}\left({\displaystyle \frac{E_d}{\mathrm{}(\mathrm{}c)}}[\overline{M}_{0,1}(^s,d)]\right)=`$ $`PD^1\left(e_{1}^{}{}_{}{}^{}\left({\displaystyle \frac{E_d^{}e_1^{}(E^+)E_d^{}}{\mathrm{}(\mathrm{}c)}}[\overline{M}_{0,1}(^s,d)]\right)\right)=`$ $`PD^1\left(E^+e_{1}^{}{}_{}{}^{}\left({\displaystyle \frac{E_d^{}E_d^{}}{\mathrm{}(\mathrm{}c)}}[\overline{M}_{0,1}(^s,d)]\right)\right)=`$ (56) $`E^+PD^1\left(e_{1}^{}{}_{}{}^{}\left({\displaystyle \frac{E_d^{}E_d^{}}{\mathrm{}(\mathrm{}c)}}[\overline{M}_{0,1}(^s,d)]\right)\right).`$ Therefore $`\stackrel{~}{J}_V=\text{exp}\left({\displaystyle \frac{t_0+pt_1}{\mathrm{}}}\right)`$ (57) $`\left(1+{\displaystyle \underset{d>0}{}}q^dPD^1e_{1}^{}{}_{}{}^{}\left({\displaystyle \frac{E_d^{}E_d^{}}{\mathrm{}(\mathrm{}c)}}[\overline{M}_{0,1}(^s,d)]\right)E^{}\right).`$ It is now visible from this presentation that $`\stackrel{~}{J}_V𝒫[[q]]`$ and $`J_V:=\underset{\lambda 0}{lim}\stackrel{~}{J}_V=`$ (58) $`\text{exp}\left({\displaystyle \frac{t_0+Ht_1}{\mathrm{}}}\right)\left(1+{\displaystyle \underset{d>0}{}}q^dPD^1e_{1}^{}{}_{}{}^{}\left({\displaystyle \frac{𝔼_d^{}𝔼_d^{}}{\mathrm{}(\mathrm{}c)}}[\overline{M}_{0,1}(^s,d)]\right)𝔼^{}\right).`$ The lemma is proven.$``$ ### 4.2. A local property of the $`J`$-function Let $`Y`$ be a smooth projective variety and $`j:^sY`$ an embedding. Suppose that $`𝒩_{^s/Y}=V^{}=𝒪(l)`$ for some $`l>0`$. Let $`C`$ be a curve in $`^s`$. The map $`j`$ gives rise to an embedding (59) $`\overline{M}_{0,n}(^s,[C])\overline{M}_{0,n}(Y,j_{}([C]))`$ ###### Lemma 4.2.1. Let $`C`$ be a degree $`d`$ rational curve in $`^s`$. Then $`\overline{M}_{0,n}(^s,d)=\overline{M}_{0,n}(Y,j_{}([C])).`$ Proof. Let $`(C^{},x_1,\mathrm{},x_n,f)\overline{M}_{0,n}(Y,j_{}([C]))`$ and $`f(C^{})=C_1C_2\mathrm{}C_p`$ be the irreducible decomposition. Then $`d[\text{line}]=[C_1]+\mathrm{}[C_p]`$. Let $`I_1=\{i:C_i^s\}`$ and $`I_2=\{1,2,\mathrm{},n\}I_1`$. Assume that $`I_2`$ is nonempty. If $$d[\text{line}]\underset{iI_1}{}[C_i]$$ has nonpositive degree in $`^s`$, we intersect with an ample divisor in $`Y`$ to see that $$d[\text{line}]\underset{iI_1}{}[C_i]=\underset{iI_2}{}[C_i]$$ is impossible. Otherwise, we intersect with $`[^s]`$ to get the same contradiction. Hence $`I_2`$ is empty and all the curves $`C_i`$ lie in $`^s`$. It follows that $`f`$ factors through $`^s`$ and therefore $`(C^{},x_1,\mathrm{}x_n,f)\overline{M}_{0,n}(^s,d)`$. On the other hand $`\overline{M}_{0,n}(^s,d)`$ is a component of $`\overline{M}_{0,n}(Y,j_{}([C]))`$ (see for example section $`\mathrm{7.4.4}`$ in ). These two arguments imply the lemma.$``$ Denote $`\overline{M}_{0,n}(Y,d):=\overline{M}_{0,n}(Y,j_{}([C]))`$, where $`C`$ is any rational curve of degree $`d`$ in $`^s`$. The following Lemma is a special case of a conjecture by Cox, Katz and Lee in which was proved in . ###### Lemma 4.2.2. $`[\overline{M}_{0,n}(Y,d)]^{\text{virt}}=𝔼_d[\overline{M}_{0,n}(^s,d)].`$ At this point we introduce a new object. For any smooth projective variety $`Y`$ and any ring $`𝒜`$, we define the formal completion of $`𝒜`$ along the semigroup of the Mori cone of $`Y`$ to be (60) $`𝒜[[q^\beta ]]:=\{{\displaystyle \underset{\beta }{}}a_\beta q^\beta ,`$ $`a_\beta 𝒜,`$ $`\beta \text{effective}\}.`$ where $`\beta H_2(Y,)`$ is effective if it is a positive linear combination of algebraic curves. This new ring behaves like a power series since for each $`\beta `$, the set of $`\alpha `$ such that $`\alpha `$ and $`\beta \alpha `$ are both effective is finite. For example, in the case of $`^s`$ we obtain the power series $`𝒜[[q]]`$. Choose generators $`D_1,\mathrm{},D_r`$ of $`H^2(Y,)`$ such that $`j^{}(D_1)=H`$ and $`j^{}(D_i)=0`$ for $`i2`$. Elements of $`H^0(Y,)H^2(Y,)`$ are of the form $`t_0+tD:=t_0+t_1D_1+\mathrm{}+t_rD_r`$. It is shown in that the generator of the quantum $`𝒟`$-module for the pure Gromov-Witten theory of $`Y`$ is (61) $$J_Y=\text{exp}\left(\frac{t_0+tD}{\mathrm{}}\right)\underset{\beta H_2(Y,)}{}q^\beta PD^1\left(e_{1}^{}{}_{}{}^{}\left(\frac{[\overline{M}_{0,1}(Y,\beta )]^{\text{virt}}}{\mathrm{}(\mathrm{}c)}\right)\right)$$ The moduli spaces $`\overline{M}_{0,1}(Y,\beta )`$ are empty unless $`\beta `$ is effective. Hence we will consider $`J_Y`$ as an element of the ring $`H^{}Y[[t_0,t_1,\mathrm{},t_r]][[q^\beta ]]`$. We extend the map $`j^{}:H^{}YH^{}^s`$ to a homomorphism (62) $$j^{}:H^{}Y[[t_0,t_1,\mathrm{},t_r]][[q^\beta ]]H^{}^s[[t_0,t_1]][[q]]$$ by defining $`j^{}(t_i)=0`$ for $`i>1`$ and $`j^{}(q^\beta )=q^\beta `$ for $`\beta j_{}(H_2(^s,))`$ and $`j^{}(q^\beta )=0`$ for $`\beta H_2(Y,)j_{}(H_2(^s,))`$. The following results show that $`J`$-function is local. ###### Theorem 4.2.1. $`j^{}(J_Y)=J_V`$. Proof. Notice that (63) $$j^{}J_Y=\text{exp}\left(\frac{t_0+t_1H}{\mathrm{}}\right)\underset{d_1=0}{\overset{\mathrm{}}{}}q_1^{d_1}j^{}PD^1\left(e_{1}^{}{}_{}{}^{}\left(\frac{[\overline{M}_{0,1}(Y,d_1)]^{\text{virt}}}{\mathrm{}(\mathrm{}c)}\right)\right)$$ Consider the following fiber diagram $$\begin{array}{ccc}\overline{M}_{0,1}(^s,d_1)& \stackrel{=}{}& \overline{M}_{0,1}(Y,d_1)\\ e_1& & e_1& & \\ ^s& \stackrel{j}{}& Y\end{array}$$ By excess intersection theory and the previous lemma $$j^{}\left(e_{1}^{}{}_{}{}^{}\left(\frac{[\overline{M}_{0,1}(Y,d_1)]^{\text{virt}}}{\mathrm{}(\mathrm{}c)}\right)\right)=𝔼^{}e_{1}^{}{}_{}{}^{}\left(\frac{𝔼_{d_1}}{\mathrm{}(\mathrm{}c)}[\overline{M}_{0,1}(^s,d_1)]\right).$$ The theorem follows easily.$``$ ###### Theorem 4.2.2. Let $`V=V^+V^{}=𝒪(k)𝒪(l)`$ on $`^s`$. Let $`\iota :X^s`$ be the zero locus of a generic section of $`V^+`$. Assume that $`X`$ is smooth and dim $`X>2`$. Let $`Y`$ be a smooth projective variety such that $`j:XY`$ with $`𝒩=𝒩_{X/Y}=\iota ^{}(V^{})`$. Assume that if $`CY`$ is a curve with $`[C]MX`$ then all the irreducible components $`C_i`$ of $`C`$ satisfy $`C_iX`$. Let $`j^{}`$ be the map constructed as in (62). Let $`J_Y`$ be the generator of the pure $`𝒟`$-module of $`Y`$ (). Then $$\iota _!(j^{}(J_Y))=\text{E}(V^+)J_V.$$ where $`\iota _!`$ is the Gysin map on cohomology. Proof. Since dim$`X>2`$ it follows that $`H^2X`$ is generated by $`\iota ^{}(H)`$. Let $`\beta _1`$ be the Poincaré dual to $`\iota ^{}(H)`$ and let $`D_1,D_2,\mathrm{},D_r`$ be a set of generators of $`H^2(Y,)`$. We may assume that $`j^{}(D_1)=\iota ^{}(H)`$ and $`j^{}(D_i)=0`$ for $`i>1`$. Let $`tD:=t_1D_1+\mathrm{}+t_rD_r`$. Now (64) $$J_Y=\text{exp}(\frac{t_0+tD}{\mathrm{}})\underset{\beta H^2(Y,)}{}q^\beta PD^1e_{1}^{}{}_{}{}^{}\left(\frac{[\overline{M}_{0,1}(Y,\beta )]^{\text{virt}}}{\mathrm{}(\mathrm{}c)}\right)$$ Consider the following diagram $$\begin{array}{ccccc}\overline{M}_{0,1}(Y,d_1\beta _1)& \stackrel{j_1}{}& \overline{M}_{0,1}(X,d_1\beta _1)& \stackrel{\iota _1}{}& \overline{M}_{0,1}(^s,d_1)\\ e_1& & e_1& & e_1\\ Y& \stackrel{j}{}& X& \stackrel{\iota }{}& ^s\end{array}$$ The square on the left is a fibre diagram. We repeatedly use the projection formula: $$\iota _{}(j^{}(J_Y))=\iota _!\left(\text{exp}\left(\frac{t_0+t_1\iota ^{}(H)}{\mathrm{}}\right)\underset{d_1=0}{\overset{\mathrm{}}{}}q_1^{d_1}PD^1j^{}e_{1}^{}{}_{}{}^{}\left(\frac{[\overline{M}_{0,1}(Y,d_1j_{}(\beta _1))]^{\text{virt}}}{\mathrm{}(\mathrm{}c)}\right)\right)$$ $$=\iota _!\left(\text{exp}\left(\frac{t_0+t_1\iota ^{}(H)}{\mathrm{}}\right)\underset{d_1=0}{\overset{\mathrm{}}{}}q_1^{d_1}\iota ^{}(𝔼^{})PD^1e_{1}^{}{}_{}{}^{}j_{1}^{}{}_{}{}^{}\left(\frac{[\overline{M}_{0,1}(Y,d_1j_{}(\beta _1))]^{\text{virt}}}{\mathrm{}(\mathrm{}c)}\right)\right)$$ $$=\text{exp}\left(\frac{t_0+t_1H}{\mathrm{}}\right)\left(𝔼^++\underset{d_1=1}{\overset{\mathrm{}}{}}q_1^{d_1}(𝔼^{})PD^1\iota _{}e_{1}^{}{}_{}{}^{}j_{1}^{}{}_{}{}^{}\left(\frac{[\overline{M}_{0,1}(Y,d_1j_{}(\beta _1))]^{\text{virt}}}{\mathrm{}(\mathrm{}c)}\right)\right)$$ (65) $$=\text{exp}\left(\frac{t_0+t_1H}{\mathrm{}}\right)\left(𝔼^++\underset{d_1=1}{\overset{\mathrm{}}{}}q_1^{d_1}(𝔼^{})PD^1e_{1}^{}{}_{}{}^{}\iota _{1}^{}{}_{}{}^{}j_{1}^{}{}_{}{}^{}\left(\frac{[\overline{M}_{0,1}(Y,d_1j_{}(\beta _1))]^{\text{virt}}}{\mathrm{}(\mathrm{}c)}\right)\right).$$ The equality in the second row follows from excess intersection theory in the left square. An argument similar to Lemma 4.2.1 implies that (66) $$\overline{M}_{0,1}(X,d_1\beta _1)=\overline{M}_{0,1}(Y,d_1j_{}(\beta _1)).$$ There are two obstruction theories in this moduli stack corresponding to the moduli problems of maps to $`X`$ and $`Y`$ respectively. They differ exactly by the bundle $`R^1\pi _{2}^{}{}_{}{}^{}e_2^{}(𝒩)`$ where $$\pi _2:\overline{M}_{0,2}(X,d_1\beta _1)\overline{M}_{0,1}(X,d_1\beta _1)$$ is the map that forgets the second marked point and $`𝒩=𝒩_{X/Y}`$. It follows that: $$j_{1}^{}{}_{}{}^{}([\overline{M}_{0,1}(Y,d_1j_{}(\beta _1))]^{\text{virt}})=\text{E}(R^1\pi _{2}^{}{}_{}{}^{}e_2^{}(𝒩))[\overline{M}_{0,1}(X,d_1\beta _1)]^{\text{virt}}.$$ Consider the following commutative diagram: $$\begin{array}{ccc}\overline{M}_{0,2}(X,d_1\beta _1)& \stackrel{e_2}{}& X\\ \iota _2& & \iota & & \\ \overline{M}_{0,2}(^s,d_1)& \stackrel{e_2}{}& ^s\end{array}$$ We compute: $$e_2^{}(𝒩)=e_2^{}(\iota ^{}(𝒪(l)))=\iota _{2}^{}{}_{}{}^{}e_{2}^{}{}_{}{}^{}(𝒪(l)).$$ There is the following fibre square: $$\begin{array}{ccc}\overline{M}_{0,2}(X,d_1\beta _1)& \stackrel{\iota _2}{}& \overline{M}_{0,2}(^s,d_1)\\ \pi _2& & \pi _2& & \\ \overline{M}_{0,1}(X,d_1\beta _1)& \stackrel{\iota _1}{}& \overline{M}_{0,1}(^s,d_1)\end{array}$$ We apply Proposition $`9.3`$ in to obtain: $$R^1\pi _{2}^{}{}_{}{}^{}e_2^{}(𝒩)=R^1\pi _{2}^{}{}_{}{}^{}\iota _{2}^{}{}_{}{}^{}\stackrel{~}{e_2}^{}(𝒪(l))=\iota _{1}^{}{}_{}{}^{}(R^1\pi _{2}^{}{}_{}{}^{}\stackrel{~}{e_2}^{}(𝒪(l)))=\iota _{1}^{}{}_{}{}^{}(V_{d_1}^{}).$$ Therefore: $`j_{1}^{}{}_{}{}^{}([\overline{M}_{0,1}(Y,d_1j_{}(\beta _1))]^{\text{virt}})=\text{E}(R^1\pi _{2}^{}{}_{}{}^{}e_2^{}(𝒩))[\overline{M}_{0,1}(X,d_1\beta _1)]^{\text{virt}}=`$ (67) $`\iota _1^{}(E_d^{})[\overline{M}_{0,1}(X,d_1\beta _1)]^{\text{virt}}.`$ On the other hand, Proposition $`\mathrm{11.2.3}`$ of says that : (68) $`\iota _{1}^{}{}_{}{}^{}[\overline{M}_{0,1}(X,d_1\beta _1)]^{\text{virt}}=𝔼_d^+[\overline{M}_{0,1}(^s,d_1)].`$ Substituting (67) and (68) in (65) we obtain $`\iota _{}(j^{}(J_Y))=\text{exp}\left({\displaystyle \frac{t_0+t_1H}{\mathrm{}}}\right)`$ (69) $`\left(𝔼^++{\displaystyle \underset{d=1}{\overset{\mathrm{}}{}}}q_1^dPD^1e_{1}^{}{}_{}{}^{}\left({\displaystyle \frac{𝔼_d}{\mathrm{}(\mathrm{}c)}}[\overline{M}_{0,1}(^s,d)]\right)(𝔼^{})\right).`$ Recall that on $`H^{}(\overline{M}_{0,1}(^s,d))`$ we have $`𝔼_d=𝔼_d^{}𝔼_d^{}e_1^{}(𝔼^+)`$. Substituting this in (69) and using the projection formula we obtain $`\iota _!(j^{}(J_Y))=\text{exp}\left({\displaystyle \frac{t_0+t_1H}{\mathrm{}}}\right)𝔼^+`$ (70) $`\left(1+{\displaystyle \underset{d=1}{\overset{\mathrm{}}{}}}q_1^dPD^1e_{1}^{}{}_{}{}^{}\left({\displaystyle \frac{𝔼_d^{}𝔼_d^{}}{\mathrm{}(\mathrm{}c)}}[\overline{M}_{0,1}(^s,d)]\right)(𝔼^{})\right).`$ The theorem is proven.$``$ ###### Remark 4.2.1. This naturally leads to local mirror symmetry. For example, let $`Y`$ be a Calabi-Yau threefold that contains $`X=^2`$. By adjunction formula, the normal bundle of $`^2`$ in $`X`$ is $`K_^2=𝒪_^2(3)`$. The last theorem asserts that the restriction of $`J_Y`$ in $`X`$ depends only on $`V=𝒪_^2(3)`$ i.e. in a neighborhood of $`X`$ in $`Y`$. Hence $`J_V`$ encodes Gromov-Witten correlators of the total space of $`𝒪_^2(3)`$ which is a local Calabi-Yau. In the next section we will see that mirror symmetry can be applied to $`J_V`$ establishing that mirror symmetry is local at least on the A-side. Interesting calculations in this direction can be found in . ## 5. Mirror Theorem In this section we will formulate and prove The Mirror Theorem which computes the generator $`J_V`$. Recall that $`V=(_{iI}𝒪(k_i))(_{jJ}𝒪(l_j))=V^+V^{}`$ with $`k_i,l_j>0`$ for all $`iI`$ and $`jJ`$. Consider the $`H^{}^s`$-valued hypergeometric series (71) $$I_V(t_0,t_1):=\text{exp}\left(\frac{t_0+t_1H}{\mathrm{}}\right)\underset{d=0}{\overset{\mathrm{}}{}}q^d\frac{_{iI}_{m=1}^{k_id}(kH+m\mathrm{})_{jJ}_{m=0}^{l_jd1}(l_jHm\mathrm{})}{_{m=1}^d(H+m\mathrm{})^{s+1}}.$$ ###### Theorem 5.0.1. (Mirror theorem). Assume that $`_{iI}k_i+_{jJ}l_js+1`$ and that $`J`$ is nonempty. If $`\|J\|>1`$ or $`_{iI}k_i+_{jJ}l_j<s+1`$ then $`J_V=I_V`$. Otherwise, there exists a power series $`I_1`$ of $`q`$ such that $`J_V(t_0,t_1+I_1)=I_V(t_0,t_1)`$ as power series of $`q`$. ###### Remark 5.0.1. The case in which $`J`$ is empty has been treated in , , , . It was suggested by Givental that his techniques should apply in the case in which $`J`$ is nonempty. ### 5.1. The Equivariant Mirror Theorem. We use Givental’s approach for complete intersections in projective spaces to prove an equivariant version of the theorem. For the remainder of this paper we will use the standard diagonal action of $`T=(^{})^{s+1}`$ on $`^s`$ with weights $`(\lambda _0,\mathrm{},\lambda _s)`$. Recall from section $`2.2`$ that $`𝒫=H_T^{}(^s)=[\lambda ][p]/_{i=0}^s(p\lambda _i)`$ and $`=(\lambda )[p]/_{i=0}^s(p\lambda _i)`$. Denote $`J_V^{eq}:=\text{exp}\left({\displaystyle \frac{t_0+pt_1}{\mathrm{}}}\right){\displaystyle \underset{d=0}{\overset{\mathrm{}}{}}}q^de_{1}^{}{}_{}{}^{}\left({\displaystyle \frac{E_d^{}E_d^{}}{\mathrm{}(\mathrm{}c)}}\right)({\displaystyle \underset{jJ}{}}l_jp)=`$ (72) $`\text{exp}\left({\displaystyle \frac{t_0+t_1p}{\mathrm{}}}\right)S(q,\mathrm{}).`$ and $`I_V^{eq}:=\text{exp}\left({\displaystyle \frac{t_0+t_1p}{\mathrm{}}}\right){\displaystyle \underset{d=0}{\overset{\mathrm{}}{}}}q^d{\displaystyle \frac{_{iI}_{m=1}^{k_id}(k_ip+m\mathrm{})_{jJ}_{m=0}^{l_jd1}(l_jpm\mathrm{})}{_{m=1}^d_{i=0}^s(p\lambda _i+m\mathrm{})}}=`$ (73) $`\text{exp}\left({\displaystyle \frac{t_0+t_1p}{\mathrm{}}}\right)S^{}(q,\mathrm{}).`$ Obviously the nonequivariant limits of $`J_V^{eq}`$ and $`I_V^{eq}`$ are respectively $`J_V`$ and $`I_V`$. The mirror theorem will follow as a nonequivariant limit of the following ###### Theorem 5.1.1. (The Equivariant Mirror Theorem). The same change of variables from theorem $`\mathrm{5.0.3}`$ transforms $`I_V^{eq}`$ into $`J_V^{eq}`$. ###### Remark 5.1.1. As the reader will see, the central part of the proof of the Mirror Theorem (up to section $`5.5`$) involves lengthy formulaes and algebraic manipulations. To simplify the presentation, we will assume during this part that $`V=𝒪(k)𝒪(l)`$. The general case is similar. We will return to the general case $`V=V^+V^{}`$ in section $`5.5`$. Recall that the equivariant Thom classes $`\varphi _i`$ of the fixed points $`p_i`$ form a basis of $``$ as a $`(\lambda )`$-vector space. Let $`S_i`$ and $`S_i^{}`$ be the restrictions of $`S`$ and $`S^{}`$ at the fixed point $`p_i`$. By the localization theorem in $`^s`$ they determine $`S`$ and $`S^{}`$. By the projection formula (74) $`S_i={\displaystyle _{_T^s}}S\varphi _i=1+{\displaystyle \underset{d=1}{\overset{\mathrm{}}{}}}q^d{\displaystyle _{\overline{M}_{0,1}(^s,d)_T}}{\displaystyle \frac{e_1^{}(lp\varphi _i)}{\mathrm{}(\mathrm{}c)}}E_d^{}E_d^{}.`$ The proof of the equivariant mirror theorem is based on exhibiting similar properties of the correlators $`S_i`$ and $`S_i^{}`$. The extra property $`S_i=1+o(\mathrm{}^2)`$ determines $`S_i`$ uniquely. After the change of variables that property is satisfied by $`S_i^{}`$ as well which implies $`S_i=S_i^{}`$. We now proceed with displaying properties of the correlators $`S_i`$ and $`S_i^{}`$. ### 5.2. Linear recursion relations The first property is given by this ###### Lemma 5.2.1. The correlators $`S_i`$ satisfy the following linear recursion relations: (75) $$S_i=1+\underset{d=1}{\overset{\mathrm{}}{}}q^dR_{id}+\underset{d=1}{\overset{\mathrm{}}{}}\underset{ji}{}q^dC_{ijd}S_j(q,\frac{\lambda _j\lambda _i}{d})$$ where $`R_{id}(\lambda )[\mathrm{}^1]`$ are polynomials in $`\mathrm{}^1`$ and (76) $$C_{ijd}=\frac{(\lambda _j\lambda _i)_{m=1}^{kd}(k\lambda _i+m\frac{\lambda _i\lambda _j}{d})_{m=0}^{ld1}(l\lambda _i+m\frac{\lambda _i\lambda _j}{d})}{d\mathrm{}(d\mathrm{}+\lambda _i\lambda _j)_{m=1}^d_{k=0,(k,m)(j,d)}^s(\lambda _i\lambda _k+m\frac{\lambda _j\lambda _i}{d})}.$$ Proof. We will see during the proof that $`S_j`$ is regular at $`\mathrm{}={\displaystyle \frac{\lambda _j\lambda _i}{d}}`$. The integrals that appear in the formula for $`S_i`$ can be evaluated using localization theorem (77) $$_{\overline{M}_{0,1}(^s,d)_T}\frac{e_1^{}(lp\varphi _i)}{\mathrm{}(\mathrm{}c)}E_d^{}E_d^{}=\underset{\mathrm{\Gamma }}{}_{(\overline{M}_\mathrm{\Gamma })_T}\frac{1}{a_\mathrm{\Gamma }\text{E}_T(N_\mathrm{\Gamma })}\left(\frac{e_1^{}(lp\varphi _i)}{\mathrm{}(\mathrm{}c)}E_d^{}E_d^{}\right)_\mathrm{\Gamma }$$ There are three types of fixed point components $`M_\mathrm{\Gamma }`$ of $`\overline{M}_{0,1}^T(^s,d)`$. The first one consists of those $`M_\mathrm{\Gamma }`$ where the component of the curve that contains the marked point is collapsed to $`p_i`$. We denote the set of these components by $`_{\mathrm{𝟏},𝐝}^𝐢`$. Let $`_{\mathrm{𝟐},𝐝}^𝐢`$ be the set of those $`M_\mathrm{\Gamma }`$ in which the marked point is mapped at $`p_i`$ and its incident component is a multiple cover of the line $`\overline{p_i,p_j}`$ for some $`ji`$. Finally let $`_{\mathrm{𝟎},𝐝}^𝐢`$ be the rest of the fixed point components. Notice first that: (78) $$\underset{\mathrm{\Gamma }_{0,d}^i}{}_{(\overline{M}_\mathrm{\Gamma })_T}\frac{1}{a_\mathrm{\Gamma }\text{E}_T(N_\mathrm{\Gamma })}\left(\frac{e_1^{}(lp\varphi _i)}{\mathrm{}(\mathrm{}c)}E_d^{}E_d^{}\right)_\mathrm{\Gamma }=0.$$ Indeed, let $`\mathrm{\Gamma }_j_{\mathrm{𝟎},𝐝}^𝐢`$ represent a fixed point component with the marked point mapped to the fixed point $`p_j`$ for some $`ji`$. Since $`(e_1^{}(\varphi _i))_{\mathrm{\Gamma }_j}=0`$ we are done. Next, in each fixed point component that belongs to $`_{\mathrm{𝟏},𝐝}^𝐢`$ the class $`c`$ is nilpotent. Indeed, if $`\mathrm{\Gamma }`$ is the decorated graph that represents such a fixed point component, let $`\overline{M}_{0,k}`$ correspond to the vertex of $`\mathrm{\Gamma }`$ that contains the marked point. Then $`kd+1`$. There is a morphism: $$\phi :M_\mathrm{\Gamma }\overline{M}_{0,k}$$ such that $`\phi ^{}(c)=c_\mathrm{\Gamma }`$. For dimension reasons $`c^{d1}=0`$ on $`\overline{M}_{0,k}`$ therefore $`{\displaystyle \frac{1}{\mathrm{}(\mathrm{}c)}}`$ is a polynomial of $`c`$ in $`\overline{M}_\mathrm{\Gamma }`$. Hence (79) $$\underset{\mathrm{\Gamma }_{1,d}^i}{}_{(\overline{M}_\mathrm{\Gamma })_T}\frac{1}{a_\mathrm{\Gamma }\text{E}_T(N_\mathrm{\Gamma })}\left(\frac{e_1^{}(lp\varphi _i)}{\mathrm{}(\mathrm{}c)}E_d^{}E_d^{}\right)_\mathrm{\Gamma }=R_{id}$$ is a polynomial in $`\mathrm{}^1`$. We now consider the fixed point components in $`_{\mathrm{𝟐},𝐝}^𝐢`$. Again let $`\mathrm{\Gamma }`$ represent such a component. For a stable map $`(C,x_1,f)`$ in $`\mathrm{\Gamma }`$ let $`C^{}`$ be the component of $`C`$ containing $`x_1`$, $`C^{\prime \prime }`$ the rest of the curve, $`x=C^{}C^{\prime \prime }`$ and $`f(x)=p_j`$ for some $`ji`$. Let $`d^{}`$ be the degree of the map f on the component $`C^{}`$ and $`d^{\prime \prime }=dd^{}`$. Then $`(C^{\prime \prime },x,f\|_{C^{\prime \prime }})`$ is a fixed point in $`\overline{M}_{0,1}(^s,d^{\prime \prime })`$. Denote its decorated graph by $`\mathrm{\Gamma }^{\prime \prime }`$. Choose the coordinates on $`C^{}`$ such that the restriction of $`f`$ on $`C^{}`$ is given by $`f(y_0,y_1)=(0,\mathrm{},z_i=y_0^d^{},\mathrm{},z_j=y_1^d^{},\mathrm{},0)`$. As $`\mathrm{\Gamma }`$ moves in $`_{\mathrm{𝟐},𝐝}^𝐢`$, the set of all such $`\mathrm{\Gamma }^{\prime \prime }`$ exhausts all the fixed points in $`\overline{M}_{0,1}(^s,d^{\prime \prime })`$ where the first marked point is not mapped to $`p_i`$. Since $`\text{Aut}(\mathrm{\Gamma })=\text{Aut}(\mathrm{\Gamma }^{\prime \prime })`$ it follows from (10) that: (80) $`a_\mathrm{\Gamma }=d^{}a_{\mathrm{\Gamma }^{\prime \prime }}.`$ The local coordinate at $`p_i`$ on the component $`C^{}`$ is $`z={\displaystyle \frac{y_1}{y_0}}`$. The weight of the $`T`$-action on $`y_l`$ is $`{\displaystyle \frac{\lambda _l}{d^{}}}`$ for $`l=0,1`$. It follows that the weight of the action on the coordinate $`z`$ and hence on $`T_{p_i}^{}C^{}`$ is $`{\displaystyle \frac{\lambda _j\lambda _i}{d^{}}}`$ therefore $`c_\mathrm{\Gamma }={\displaystyle \frac{\lambda _j\lambda _i}{d^{}}}`$. Now $`\text{E}_T(N_\mathrm{\Gamma })`$ can be split in three pieces: smoothing the node $`x`$ and deforming the maps $`f\|_{C^{\prime \prime }}`$ and $`f\|_C^{}`$. It follows $`\text{E}_T(N_\mathrm{\Gamma })=({\displaystyle \frac{\lambda _j\lambda _i}{d^{}}}c_\mathrm{\Gamma }^{\prime \prime })\text{E}_T(N_{\mathrm{\Gamma }^{\prime \prime }})`$ (81) $`{\displaystyle \underset{m=0}{\overset{d^{}1}{}}}{\displaystyle \underset{k=0,(m,k)(0,i)}{\overset{s}{}}}(\lambda _i\lambda _k+m{\displaystyle \frac{\lambda _j\lambda _i}{d^{}}}).`$ Next, we find the localization of $`E_d^{}`$ and $`E_d^{}`$ on the fixed point component $`M_{d_1d_2}^i`$. Consider the normalization sequence at the node $`x`$ (82) $$0𝒪_C𝒪_C^{}𝒪_{C^{\prime \prime }}𝒪_x0.$$ Twisting it by $`f^{}(V^+)`$ and $`f^{}(V^{})`$ respectively and taking the cohomology sequence yields $$(E_d^{})_\mathrm{\Gamma }=(l\lambda _j)(E_{d^{\prime \prime }}^{})_{\mathrm{\Gamma }^{\prime \prime }}(E_d^{}^{})_\mathrm{\Gamma }^{}$$ and (83) $$(E_d^+)_\mathrm{\Gamma }=\frac{(E_d^{}^+)_\mathrm{\Gamma }^{}(E_{d^{\prime \prime }}^+)_{\mathrm{\Gamma }^{\prime \prime }}}{k\lambda _j}$$ An explicit basis for $`H^1(C^{},f^{}(V^{}))=H^1(𝒪_^1(ld^{})`$ consists of (84) $`{\displaystyle \frac{y_0^sy_1^{ld^{}2s}}{(y_0y_1)^{ld^{}1}}}={\displaystyle \frac{1}{y_0^{ld^{}s1}y_1^{1+s}}}:s=0,1,\mathrm{},ld^{}2.`$ It allows us to compute: $$(E_d^{}^{})_\mathrm{\Gamma }^{}=\underset{s=0}{\overset{ld^{}2}{}}\left(\frac{1+sld^{}}{d^{}}\lambda _i\frac{1+s}{d^{}}\lambda _j\right)=\underset{s=1}{\overset{ld^{}1}{}}\left(l\lambda _i+s\frac{\lambda _i\lambda _j}{d^{}}\right).$$ Therefore we have: (85) $$(E_d^{})_\mathrm{\Gamma }=(l\lambda _j)\underset{s=1}{\overset{ld^{}1}{}}\left(l\lambda _i+s\frac{\lambda _i\lambda _j}{d^{}}\right)(E_{d^{\prime \prime }}^{})_{\mathrm{\Gamma }^{\prime \prime }}.$$ A basis for $`H^0(C^{},f^{}(V^+))=H^0(𝒪_^1(kd^{}))`$ consists of monomials $`y_0^sy_1^{kd^{}s}`$ for $`s=0,\mathrm{}kd^{}`$. It can be used to calculate $`(E_d^{}^+)_\mathrm{\Gamma }^{}`$ similarly to $`(E_d^{}^{})_\mathrm{\Gamma }^{}`$ above. Recall from (55) that $`E_d^+=e_1^{}(\text{E}(V^+))E_d^{}`$. The line bundle $`e_1^{}(V^+)`$ is trivial on $`M_\mathrm{\Gamma }`$ but the torus acts on it with weight $`k\lambda _i`$. Hence $`(E_d^+)_\mathrm{\Gamma }=k\lambda _i(E_d^{})_\mathrm{\Gamma }`$. Substituting in (83) yields (86) $$(E_d^{})_\mathrm{\Gamma }=\underset{r=1}{\overset{kd^{}}{}}\left(k\lambda _i+r\frac{\lambda _j\lambda _i}{d^{}}\right)(E_{d^{\prime \prime }}^{})_{\mathrm{\Gamma }^{\prime \prime }}.$$ We pause here to show that $`S_i`$ is regular at $`\mathrm{}={\displaystyle \frac{\lambda _j\lambda _i}{d}}`$ for any $`ji`$ and any $`d>0`$. It follows from (78) and (79) that: (87) $$S_i=1+\underset{d=1}{\overset{\mathrm{}}{}}q^dR_{id}+\underset{\mathrm{\Gamma }_{2,d}^i}{}q^d_{(\overline{M}_\mathrm{\Gamma })_T}\frac{(l\lambda _i)_{ki}(\lambda _i\lambda _k)(E_d^{}E_d^{})_\mathrm{\Gamma }}{\mathrm{}(\mathrm{}c_\mathrm{\Gamma })a_\mathrm{\Gamma }\text{E}_T(N_\mathrm{\Gamma })}.$$ From this representation of $`S_i`$ it is clear that the coefficients of the power series $`S_i=_{d=0}^{\mathrm{}}S_{id}q^d`$ belong to $`(\lambda ,\mathrm{})`$. But $`c_\mathrm{\Gamma }={\displaystyle \frac{\lambda _j\lambda _i}{d^{}}}`$ for some $`d^{}d`$ and $`R_{id}`$ has poles only at $`\mathrm{}=0`$ therefore $`S_i`$ is regular at $`\mathrm{}={\displaystyle \frac{\lambda _i\lambda _j}{d}}`$. We use the equations (81), (86) and (85) to compute $`{\displaystyle \underset{\mathrm{\Gamma }_{2,d}^i}{}}q^d{\displaystyle _{(\overline{M}_\mathrm{\Gamma })_T}}{\displaystyle \frac{(l\lambda _j)_{ki}(\lambda _i\lambda _k)(E_d^{}E_d^{})_\mathrm{\Gamma }}{\mathrm{}(\mathrm{}c_\mathrm{\Gamma })a_\mathrm{\Gamma }\text{E}_T(N_\mathrm{\Gamma })}}=`$ $`{\displaystyle \underset{d^{}=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{ji}{}}q^d^{}C_{ijd^{}}{\displaystyle \underset{\mathrm{\Gamma }^{\prime \prime }}{}}q^{d^{\prime \prime }}{\displaystyle _{(\overline{M}_{\mathrm{\Gamma }^{\prime \prime }})_T}}{\displaystyle \frac{l\lambda _j_{kj}(\lambda _j\lambda _k)(E_{d^{\prime \prime }}^{}E_{d^{\prime \prime }}^{})_{\mathrm{\Gamma }^{\prime \prime }}}{(\frac{\lambda _j\lambda _i}{d^{}})(\frac{\lambda _j\lambda _i}{d^{}}c_\mathrm{\Gamma }^{\prime \prime })a_{\mathrm{\Gamma }^{\prime \prime }}\text{E}_T(N_{\mathrm{\Gamma }^{\prime \prime }})}}=`$ (88) $`{\displaystyle \underset{d^{}=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{ji}{}}q^d^{}C_{ijd^{}}S_j(q,{\displaystyle \frac{\lambda _j\lambda _i}{d^{}}}).`$ The lemma follows by substituting the above identity into (87)$``$. ###### Lemma 5.2.2. The correlators $`S_i^{}`$ satisfy the same linear recursion relations as $`S_i`$. Proof. We know that $`S_i^{}=_{d=0}^{\mathrm{}}q^dS_{id}^{}`$ with (89) $$S_{id}^{}=\frac{_{m=1}^{kd}(k\lambda _i+m\mathrm{})_{m=0}^{ld1}(l\lambda _im\mathrm{})}{d\mathrm{}_{m=1}^d_{j=0,(j,m)(i,d)}^s(\lambda _i\lambda _j+m\mathrm{})}$$ Note that $`S_{id}^{}(\lambda ,\mathrm{})`$ is a proper rational expression of $`\mathrm{}`$. It has multiple poles at $`\mathrm{}=0`$ and simple poles at $`\mathrm{}={\displaystyle \frac{\lambda _r\lambda _i}{m}}`$ for any $`ri`$ and any $`1md`$. Applying calculus of residues in the $`\mathrm{}`$-variable yields: $`S_{id}^{}=R_{id}+{\displaystyle \underset{m=1}{\overset{d}{}}}{\displaystyle \underset{ri}{}}{\displaystyle \frac{1}{d\mathrm{}(\lambda _i\lambda _r+m\mathrm{})}}`$ (90) $`\times {\displaystyle \frac{_{n=1}^{kd}(k\lambda _i+n\frac{\lambda _r\lambda _i}{m})_{n=0}^{ld1}(l\lambda _in\frac{\lambda _r\lambda _i}{m})}{_{n=1,(j,n)(r,m)}^d_{j=0,(j,n)(i,d)}^s(\lambda _i\lambda _j+n\frac{\lambda _r\lambda _i}{m})}}`$ for some polynomials $`R_{id}(\lambda )[\mathrm{}^1]`$ such that $`R_{id}(0)=0`$. Substitute equation (90) in (89) to obtain: $`S_i^{}=1+{\displaystyle \underset{d=1}{\overset{\mathrm{}}{}}}q^dR_{id}+{\displaystyle \underset{d=1}{\overset{\mathrm{}}{}}}q^d{\displaystyle \underset{ri}{}}{\displaystyle \underset{m=1}{\overset{d}{}}}{\displaystyle \frac{1}{d\mathrm{}(\lambda _i\lambda _r+m\mathrm{})}}`$ (91) $`\times {\displaystyle \frac{_{n=1}^{kd}(k\lambda _i+n\frac{\lambda _r\lambda _i}{m})_{n=0}^{ld1}(l\lambda _in\frac{\lambda _r\lambda _i}{m})}{_{n=1,(j,n)(r,m)}^d_{j=0,(j,n)(i,d)}^s(\lambda _i\lambda _j+n\frac{\lambda _r\lambda _i}{m})}}.`$ Changing the order of summation in the last equation yields: $`S_i^{}1{\displaystyle \underset{d=1}{\overset{\mathrm{}}{}}}q^dR_{id}={\displaystyle \underset{ri}{}}{\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}q^m{\displaystyle \frac{1}{\mathrm{}(\lambda _i\lambda _r+m\mathrm{})}}`$ (92) $`\times {\displaystyle \underset{d=m}{\overset{\mathrm{}}{}}}q^{dm}{\displaystyle \frac{_{n=1}^{kd}(k\lambda _i+n\frac{\lambda _r\lambda _i}{m})_{n=0}^{ld1}(l\lambda _in\frac{\lambda _r\lambda _i}{m})}{d_{n=1,(j,n)(r,m)}^d_{j=0,(j,n)(i,d)}^s(\lambda _i\lambda _j+n\frac{\lambda _r\lambda _i}{m})}}.`$ The lemma follows from the identity $`{\displaystyle \underset{d=m}{\overset{\mathrm{}}{}}}q^{dm}{\displaystyle \frac{_{n=1}^{kd}(k\lambda _i+n\frac{\lambda _r\lambda _i}{m})_{n=0}^{ld1}(l\lambda _in\frac{\lambda _r\lambda _i}{m})}{d_{n=1,(j,n)(r,m)}^d_{j=0,(j,n)(i,d)}^s(\lambda _i\lambda _j+n\frac{\lambda _r\lambda _i}{m})}}`$ $`={\displaystyle \frac{(\frac{\lambda _r\lambda _i}{m})_{n=1}^{km}(k\lambda _i+n\frac{\lambda _r\lambda _i}{m})_{n=0}^{lm1}(l\lambda _in\frac{\lambda _r\lambda _i}{m})}{_{n=1,(j,n)(r,m)}^m_{j=0}^s(\lambda _i\lambda _j+n\frac{\lambda _r\lambda _i}{m})}}`$ (93) $`\times {\displaystyle \underset{u=0}{\overset{\mathrm{}}{}}}q^u{\displaystyle \frac{_{n=1}^{ku}(k\lambda _r+n\frac{\lambda _r\lambda _i}{m})_{n=0}^{lu1}(l\lambda _rn\frac{\lambda _r\lambda _i}{m})}{u(\frac{\lambda _r\lambda _i}{m})_{n=1,(j,n)(r,u)}^u_{j=0}^s\left(\lambda _r\lambda _j+n\frac{\lambda _r\lambda _i}{m}\right)}}`$ ### 5.3. Double polynomiality Recall from Section $`3.1`$ that $`V`$ induces a modified equivariant integral $`\omega _V:(\lambda )`$ defined as follows: (94) $$\omega _V(a):=_{_𝐓^s}a\frac{E^+}{E^{}}.$$ As one can see, in the case $`V=𝒪(k)𝒪(l)`$, this modified equivariant integral simplifies via $`{\displaystyle \frac{E^+}{E^{}}}={\displaystyle \frac{kp}{lp}}={\displaystyle \frac{k}{l}}`$. We have chosen not to simplify this integral in the proof of the following Lemma so that it is easier to see how to proceed in the general case. ###### Lemma 5.3.1. If $`z`$ is a variable, the expression: $$P(z,\mathrm{})=\omega _V\left(e^{pz}S(qe^z\mathrm{},\mathrm{})S(q,\mathrm{})\right)$$ belongs to $`(\lambda )[\mathrm{}][[q,z]].`$ Proof. In Section $`2.3`$ we introduced the action of $`T^{}=T\times ^{}`$ on $`^s\times ^1`$ with weights $`(\lambda _0,\mathrm{},\lambda _s)`$ on the first factor and $`(\mathrm{},0)`$ in the second factor. Consider the following $`T^{}`$-equivariant diagram $$\begin{array}{ccc}\overline{M}_{0,1}(^s\times ^1,(d,1))& \stackrel{e_1}{}& ^s\times ^1\\ \pi & & \\ M_d=\overline{M}_{0,0}(^s\times ^1,(d,1))\end{array}$$ Define: (95) $$W_d=W_d^+W_d^{}:=\pi _{}((e_1)^{}(𝒪_^s(k)𝒪_^1))R^1\pi _{}((e_1)^{}(𝒪_^s(l)𝒪_^1))$$ The lemma will follow from the identity: (96) $$P(z,\mathrm{})=\underset{d=0}{\overset{\mathrm{}}{}}q^d_{(M_d)_T^{}}e^{z\psi ^{}\kappa }\text{E}_T^{}(W_d)$$ where $`\psi `$ and $`\kappa `$ were defined in Section $`2.3`$. The localization formula for the diagonal action of $`T`$ on $`^s`$ applied to the left side gives (97) $$P(z,\mathrm{})=\underset{i=0}{\overset{s}{}}\frac{S_i(qe^z\mathrm{},\mathrm{})e^{z\lambda _i}S_i(q,\mathrm{})}{_{ki}(\lambda _i\lambda _k)}\left(\frac{k\lambda _i}{l\lambda _i}\right).$$ We recall from identity (74) (98) $$S_i=1+\underset{d=1}{\overset{\mathrm{}}{}}q^d_{\overline{M}_{0,1}(^s,d)_T}\frac{e_1^{}(lp\varphi _i)}{\mathrm{}(\mathrm{}c)}E_d^{}E_d^{}.$$ To compute the integrals on the right side of (96) we will use localization for the action of $`T^{}`$ on $`M_d`$. In Section $`2.3`$ we found that the components of the fixed point loci have the form $`M_{d_1d_2}^i=\overline{M}_{\mathrm{\Gamma }_{d_1}^i}\times \overline{M}_{\mathrm{\Gamma }_{d_2}^i}`$ for some $`i=0,1,\mathrm{},s`$ and a splitting $`d=d_1+d_2`$. We first compute the restriction of $`\text{E}_T^{}(W_d)`$ in such a component. Consider the following normalization sequence (99) $$0𝒪_C𝒪_{C_0}𝒪_{C_1}𝒪_{C_2}𝒪_{x_1}𝒪_{x_2}0$$ Twist (99) by $`f^{}(𝒪(l)𝒪_^1)`$ and take the corresponding long exact cohomology sequence. We obtain $$0𝒪_{x_1}(l)𝒪_{x_2}(l)W_d^{}W_{d_1}^{}W_{d_2}^{}0.$$ The first piece is trivial since it comes from the isomorphism $$(𝒪_^s(l)\times 𝒪_^1)\|_{C_0}𝒪_{C_0}.$$ The left hand side is generated by $`{\displaystyle \frac{1}{z_{i}^{}{}_{}{}^{l}}}`$ therefore the weight of that piece is $`l\lambda _i`$. It follows that $$\text{E}_T^{}(W_d^{})=(l\lambda _i)E_{d_1}^{}E_{d_2}^{}.$$ Similarly, twisting the normalization sequence (99) by $`f^{}(𝒪(k)𝒪_^1)`$ and taking the corresponding cohomology sequence we obtain: $$\text{E}_T^{}(W_d^+)=(k\lambda _i)E_{d_1}^{}E_{d_2}^{}.$$ We now use the localization theorem to calculate the integrals on the right side of (96). The equivariant Euler class of the normal bundle of the fixed point component $`M_{d_1d_2}^i`$ has been calculated at the end of section $`2.3`$. (100) $`{\displaystyle _{(M_d)_T^{}}}e^{z\psi ^{}\kappa }\text{E}_T^{}(W_d)={\displaystyle \underset{\mathrm{\Gamma }_{d_1}^i,\mathrm{\Gamma }_{d_2}^i}{}}(k\lambda _i)(l\lambda _i){\displaystyle \frac{e^{z(\lambda _i+d_2\mathrm{})}}{_{ki}(\lambda _i\lambda _k)}}`$ $`\times {\displaystyle _{(\overline{M}_{\mathrm{\Gamma }_{d_1}^i})_T}}{\displaystyle \frac{1}{E_T(N_{\mathrm{\Gamma }_{d_1}^i})}}\left({\displaystyle \frac{e_1^{}(\varphi _i)E_{d_1}^{}E_{d_1}^{}}{\mathrm{}(\mathrm{}c_1)}}\right)_{\mathrm{\Gamma }_{d_1}^i}{\displaystyle _{(\overline{M}_{\mathrm{\Gamma }_{d_2}^i})_T}}{\displaystyle \frac{1}{E_T(N_{\mathrm{\Gamma }_{d_2}^i})}}\left({\displaystyle \frac{e_1^{}(\varphi _i)E_{d_2}^{}E_{d_2}^{}}{\mathrm{}(\mathrm{}c_2)}}\right)_{\mathrm{\Gamma }_{d_2}^i}=`$ $`{\displaystyle \underset{\mathrm{\Gamma }_{d_1}^i,\mathrm{\Gamma }_{d_2}^i}{}}{\displaystyle \frac{k\lambda _i}{l\lambda _i}}{\displaystyle \frac{e^{z(\lambda _i+d_2\mathrm{})}}{_{ki}(\lambda _i\lambda _k)}}{\displaystyle _{(\overline{M}_{\mathrm{\Gamma }_{d_1}^i})_T}}{\displaystyle \frac{1}{E_T(N_{\mathrm{\Gamma }_{d_1}^i})}}\left({\displaystyle \frac{e_1^{}(l\lambda _i\varphi _i)E_{d_1}}{\mathrm{}(\mathrm{}c_1)}}\right)_{\mathrm{\Gamma }_{d_1}^i}\times `$ $`\times {\displaystyle _{(\overline{M}_{\mathrm{\Gamma }_{d_2}^i})_T}}{\displaystyle \frac{1}{E_T(N_{\mathrm{\Gamma }_{d_2}^i})}}\left({\displaystyle \frac{e_1^{}(l\lambda _i\varphi _i)E_{d_2}}{\mathrm{}(\mathrm{}c_2)}}\right)_{\mathrm{\Gamma }_{d_2}^i}.`$ If we use localization to compute $`S_i`$ in (98) and then substitute in (97) we obtain the right side of the last equation.$``$ ###### Lemma 5.3.2. If $`z`$ is a variable, the expression: (101) $`P^{}(z,h)=\omega _V\left(S^{}(qe^z\mathrm{},\mathrm{})e^{pz}S^{}(q,\mathrm{})\right)`$ belongs to $`(\lambda )[\mathrm{}][[q,z]].`$ Proof. The lemma will follow from the identity (102) $$P^{}(z,h)=\underset{d=0}{\overset{\mathrm{}}{}}q^d_{(N_d)_T^{}}e^{z\kappa }\underset{m=0}{\overset{kd}{}}(k\kappa m\mathrm{})\underset{m=1}{\overset{ld1}{}}(l\kappa +m\mathrm{}).$$ For $`d=0`$ the convention $$_{(N_d)_T^{}}e^{z\kappa }\underset{m=0}{\overset{kd}{}}(k\kappa +m\mathrm{})\underset{m=1}{\overset{ld1}{}}(l\kappa +m\mathrm{})=_{_𝐓^s}e^{pz}\left(\frac{_{iI}k_ip}{_{jJ}l_jp}\right)$$ is taken. Apply the localization formula to the integrals (101). $`P^{}(z,\mathrm{})={\displaystyle \underset{i=0}{\overset{s}{}}}{\displaystyle \frac{k\lambda _ie^{\lambda _iz}}{(l\lambda _i)_{ji}(\lambda _i\lambda _j)}}`$ $`\times {\displaystyle \underset{d_1=0}{\overset{\mathrm{}}{}}}(qe^z\mathrm{})^{d_1}{\displaystyle \frac{_{m=1}^{kd_1}(k\lambda _i+m\mathrm{})_{m=0}^{ld_11}(l\lambda _im\mathrm{})}{_{m=1}^{d_1}_{j=0}^s(\lambda _i\lambda _j+m\mathrm{})}}`$ $`\times {\displaystyle \underset{d_2=0}{\overset{\mathrm{}}{}}}q^{d_2}{\displaystyle \frac{_{m=1}^{kd_2}(k\lambda _im\mathrm{})_{m=0}^{ld_21}(l\lambda _i+m\mathrm{})}{_{m=1}^{d_2}_{j=0}^s(\lambda _i\lambda _jm\mathrm{})}}=`$ $`{\displaystyle \underset{i=0}{\overset{s}{}}}{\displaystyle \frac{1}{_{ji}(\lambda _i\lambda _j)}}{\displaystyle \underset{d_1=0}{\overset{\mathrm{}}{}}}q^{d_1z}e^{(\lambda _i+d_1\mathrm{})z}{\displaystyle \frac{_{m=0}^{kd_1}(k\lambda _i+m\mathrm{})_{m=1}^{ld_11}(l\lambda _im\mathrm{})}{_{m=1}^{d_1}_{j=0}^s(\lambda _i\lambda _j+m\mathrm{})}}`$ (103) $`\times {\displaystyle \underset{d_2=0}{\overset{\mathrm{}}{}}}q^{d_2}{\displaystyle \frac{_{m=1}^{kd_2}(k\lambda _im\mathrm{})_{m=0}^{ld_21}(l\lambda _i+m\mathrm{})}{_{m=1}^{d_2}_{j=0}^s(\lambda _i\lambda _jm\mathrm{})}}.`$ But for $`d_1,d_2>0`$ $`{\displaystyle \frac{_{m=0}^{kd_1}(k\lambda _i+m\mathrm{})_{m=1}^{ld_11}(l\lambda _im\mathrm{})_{m=1}^{kd_2}(k\lambda _im\mathrm{})_{m=0}^{ld_21}(l\lambda _i+m\mathrm{})}{_{ji}(\lambda _i\lambda _j)_{m=1}^{d_1}_{j=0}^s(\lambda _i\lambda _j+m\mathrm{})_{m=1}^{d_2}_{j=0}^s(\lambda _i\lambda _jm\mathrm{})}}=`$ $`{\displaystyle \frac{_{m=0}^{k(d_1+d_2)}(k(\lambda _i+d_1\mathrm{})m\mathrm{})_{m=1}^{l(d_1+d_2)1}(l(\lambda _i+d_1\mathrm{})+m\mathrm{})}{_{j=0}^s_{m=0,(j,m)(i,d_1)}^{d_1+d_2}(\lambda _i+d_1\mathrm{}\lambda _jm\mathrm{})}}.`$ Therefore $`P^{}(z,\mathrm{})={\displaystyle \underset{d=0}{\overset{\mathrm{}}{}}}q^d{\displaystyle \underset{d_1=0}{\overset{d}{}}}{\displaystyle \underset{i=0}{\overset{s}{}}}e^{(\lambda _i+d_1\mathrm{})z}`$ (104) $`\times {\displaystyle \frac{_{m=0}^{kd}(k(\lambda _i+d_1\mathrm{})m\mathrm{})_{m=1}^{ld1}(l(\lambda _i+d_1\mathrm{})+m\mathrm{})}{_{j=0}^s_{m=0,(j,m)(i,d_1)}^d(\lambda _i+d_1\mathrm{}\lambda _jm\mathrm{})}}.`$ By the localization formula in $`N_d`$ we can see that $$P^{}(z,\mathrm{})=\underset{d=0}{\overset{\mathrm{}}{}}q^d_{(N_d)_T^{}}e^{z\kappa }\underset{m=0}{\overset{kd}{}}(k\kappa m\mathrm{})\underset{m=1}{\overset{ld1}{}}(l\kappa +m\mathrm{}).$$ The lemma is proven.$``$ ### 5.4. Mirror transformation and uniqueness The following two theorems carry over from . The first lemma deals with uniqueness. ###### Lemma 5.4.1. Let $`S=_{d=0}^{\mathrm{}}S_dq^d`$ and $`S^{}=_{d=0}^{\mathrm{}}S_d^{}q^d`$ be two power series with coefficients in $`[[\mathrm{}^1]]`$ that satisfy the following conditions: 1. $`S_0=S_0^{}=1`$ 2. They both satisfy the recursion relations of Lemma 5.2.1. 3. They both have the double polynomiality property of Lemma 5.3.2. 4. For any d, $`S_dS_d^{}`$ mod $`(\mathrm{}^2)`$. Then $`S=S^{}`$. The second lemma describes a transformation which preserves the properties of lemma $`\mathrm{5.4.1}`$. ###### Lemma 5.4.2. Let $`I_1`$ be a power series in $`q`$ whose first term is zero. Then $`\text{exp}(\frac{I_1p}{\mathrm{}})S(qe^{I_1},\mathrm{})`$ satisfies conditions $`1,2,3`$ of Lemma $`\mathrm{5.4.1}`$. ### 5.5. The conclusion of the proof of Mirror Theorem Recall that $$I_V^{eq}=\text{exp}\left(\frac{t_0+t_1p}{\mathrm{}}\right)\left(1+\underset{d=1}{\overset{\mathrm{}}{}}q^d\frac{_{iI}_{m=1}^{k_id}(k_ip+m\mathrm{})_{jJ}_{m=0}^{l_jd1}(l_jpm\mathrm{})}{_{m=1}^d_{i=0}^s(p\lambda _i+m\mathrm{})}\right).$$ We are assuming that there is at least one negative line bundle. We expand the second factor of $`I_V^{eq}`$ as a polynomial of $`\mathrm{}^1`$. Each negative line bundle produces a factor of $`{\displaystyle \frac{p}{\mathrm{}}}`$. For example, in the case $`V=𝒪(k)𝒪(l)`$ the expansion yields: (105) $$I_V^{eq}=\text{exp}\left(\frac{t_0+pt}{\mathrm{}}\right)\left(1+\frac{p}{\mathrm{}}\underset{d=1}{\overset{\mathrm{}}{}}q^d\frac{(1)^{ld}(ld1)!(kd)!}{(d!)^{s+1}}\frac{1}{\mathrm{}^{d(s+1kl)}}+o(\frac{1}{\mathrm{}^2})\right)$$ If $`V`$ contains two or more negative line bundles it follows that $$I_V^{eq}=\text{exp}\left(\frac{t_0+pt}{\mathrm{}}\right)\left(1+o(\frac{1}{\mathrm{}^2})\right)$$ Lemma $`(\mathrm{5.4.1})`$ and $`(\mathrm{5.4.2})`$ imply that $`J_V^{eq}=I_V^{eq}`$. If $`_{iI}k_i+_{jJ}l_j<s+1`$ the presence of $`{\displaystyle \frac{1}{\mathrm{}^{d(s+1kl)}}}`$ in the above expansion shows that again $$I_V^{eq}=\text{exp}\left(\frac{t_0+pt}{\mathrm{}}\right)\left(1+o(\frac{1}{\mathrm{}^2})\right)$$ hence $`J^{eq}=I^{eq}`$. We may assume that $`_{iI}k_i+_{jJ}l_j=s+1`$ and $`\|J\|=1`$. In this case $$I_V^{eq}=\text{exp}\left(\frac{t_0+pt}{\mathrm{}}\right)\left(1+I_1\frac{p}{\mathrm{}}+o(\frac{1}{\mathrm{}^2})\right)$$ where $`I_1`$ is a power series of $`q`$ whose first term is zero. For example if $`V=𝒪(k)𝒪(l)`$ the power series $`I_1`$ is $$I_1=\underset{d=1}{\overset{\mathrm{}}{}}q^d\frac{(1)^{ld}(ld1)!(kd)!}{(d!)^{s+1}}$$ Recall that $`S=1+o(\mathrm{}^2)`$. Therefore (106) $$\text{exp}(\frac{I_1p}{\mathrm{}})S(qe^{I_1},\mathrm{})=1+I_1\frac{p}{\mathrm{}}+o(\mathrm{}^2).$$ Lemma $`\mathrm{5.4.2}`$ implies that both $`\text{exp}\left({\displaystyle \frac{I_1p}{\mathrm{}}}\right)S(qe^{I_1},\mathrm{})`$ and $`S^{}(q,\mathrm{})`$ satisfy the conditions of the Lemma $`\mathrm{5.4.1}`$. It follows that $$\text{exp}\left(\frac{I_1p}{\mathrm{}}\right)S(qe^{I_1},\mathrm{})=S^{}(q,\mathrm{})$$ Multiplying both sides of this identity by $`\text{exp}\left({\displaystyle \frac{t_0+pt}{\mathrm{}}}\right)`$ yields (107) $$J_V^{eq}(t_0,t+I_1)=I_V^{eq}(t_0,t).$$ This completes the proof.$``$ ###### Corollary 5.5.1. Let $`V=(_{iI}𝒪(k_i))(_{jJ}𝒪(l_j))`$. For $`\|J\|>1`$ or $`k+l<s+1`$ $$e_{1}^{}{}_{}{}^{}\left(\frac{𝔼_d^{}𝔼_d^{}}{\mathrm{}(\mathrm{}c)}\right)=\frac{_{iI}_{m=1}^{k_id}(k_iH+m\mathrm{})_{jJ}_{m=1}^{l_jd1}(l_jHm\mathrm{})}{_{m=1}^d(H+m\mathrm{})^{s+1}}.$$ Proof. As mentioned above in this case we have $`J_V^{eq}=I_V^{eq}`$. Recall that $$J_V^{eq}=\text{exp}\left(\frac{t_0+pt_1}{\mathrm{}}\right)\left(1+\underset{d>0}{}q^de_{1}^{}{}_{}{}^{}\left(\frac{E_d^{}E_d^{}}{\mathrm{}(\mathrm{}c)}\right)\underset{jJ}{}(l_jp)\right).$$ We obtain the equivariant identity: $$e_{1}^{}{}_{}{}^{}\left(\frac{E_d^{}E_d^{}}{\mathrm{}(\mathrm{}c)}\right)\underset{jJ}{}(l_jp)=\frac{_{iI}_{m=1}^{k_id}(k_ip+m\mathrm{})_{jJ}_{m=0}^{l_jd1}(l_jpm\mathrm{})}{_{m=1}^d_{k=0}^s(p\lambda _k+m\mathrm{})}.$$ The restriction of $`p`$ to any fixed point $`p_i`$ is nonzero. This implies that $`p`$ is invertible. Therefore we obtain $$e_{1}^{}{}_{}{}^{}\left(\frac{E_d^{}E_d^{}}{\mathrm{}(\mathrm{}c)}\right)=\frac{_{iI}_{m=1}^{k_id}(k_ip+m\mathrm{})_{jJ}_{m=1}^{l_jd1}(l_jpm\mathrm{})}{_{m=1}^d_{k=1}^{s+1}(p\lambda _k+m\mathrm{})}.$$ The nonequivariant limit of this identity reads $$e_{1}^{}{}_{}{}^{}\left(\frac{𝔼_d^{}𝔼_d^{}}{\mathrm{}(\mathrm{}c)}\right)=\frac{_{iI}_{m=1}^{k_id}(k_iH+m\mathrm{})_{jJ}_{m=1}^{l_jd1}(l_jHm\mathrm{})}{_{m=1}^d(H+m\mathrm{})^{s+1}}.$$ The theorem is proven.$``$ This corollary is particularly useful when $`\text{Euler}(V^{})=0`$ in $`^s`$. In this case $`J_V=I_V=\text{exp}\left({\displaystyle \frac{t_0+Ht_1}{\mathrm{}}}\right)`$ hence the mirror theorem is true trivially. An example of such a situation is $`V=𝒪_^1(1)𝒪_^1(1)`$ which is treated in the next section. ## 6. Examples ### 6.1. Multiple covers Let $`C`$ be a smooth rational curve in a Calabi-Yau threefold $`X`$ with normal bundle $`N=𝒪(1)𝒪(1)`$ and $`\beta =[C]H_2(X,)`$. Since $`K_X=𝒪_X`$ the expected dimension of the moduli space $`\overline{M}_{0,0}(X,d\beta )`$ is zero. However, this moduli space contains a component of positive dimension, namely $`\overline{M}_{0,0}(^1,d)`$. Indeed, let $`f:^1C`$ be an isomorphism, and $`g:^1^1`$ a degree $`d`$ multiple cover. Then $`fg`$ is a stable map that belongs to $`\overline{M}_{0,0}(X,d\beta )`$. For a proof of the fact that $`\overline{M}_{0,0}(^1,d)`$ is a component of $`\overline{M}_{0,0}(X,d\beta )`$ see section $`\mathrm{7.4.4}`$ in . Let $`N_d`$ be the degree of $`[\overline{M}_{0,0}(X,d\beta )]^{\text{virt}}`$. We want to compute the contribution $`n_d`$ of $`\overline{M}_{0,0}(^1,d)`$ to $`N_d`$. Kontsevich asserted in and Behrend proved in that the restriction of $`[\overline{M}_{0,0}(X,d\beta )]^{\text{virt}}`$ to $`\overline{M}_{0,0}(^1,d)`$ is precisely $`𝔼_d`$ for $`V=𝒪(1)𝒪(1)`$. Therefore: $$n_d=_{\overline{M}_{0,0}(^1,d)}𝔼_d.$$ Note that dim $`\overline{M}_{0,0}(^1,d)=2d2`$ and the rank of the bundle $`V_d`$ is also $`2d2`$. We use the mirror theorem to compute numbers $`n_d`$. Since $`V`$ contains two negative line bundles we can apply Corollary $`\mathrm{5.5.1}`$ $$e_{1}^{}{}_{}{}^{}\left(\frac{𝔼_d}{\mathrm{}(\mathrm{}c)}\right)=\frac{_{m=1}^{d1}(Hm\mathrm{})^2}{_{m=1}^d(H+m\mathrm{})^2}=\frac{1}{(H+d\mathrm{})^2}.$$ An expansion of the left hand side using the divisor property for the modified gravitational descendants yields $$e_{1}^{}{}_{}{}^{}\left(\frac{𝔼_d}{\mathrm{}(\mathrm{}c)}\right)=\frac{dn_d}{\mathrm{}^2}+\frac{H}{\mathrm{}^3}_{\overline{M}_{0,1}(^1,d)}c𝔼_d.$$ where $`c`$ is the chern class of the cotangent line bundle at the marked point. On the other hand: $$\frac{1}{(H+d\mathrm{})^2}=\frac{1}{d^2\mathrm{}^2}\frac{2H}{d^3\mathrm{}^3}.$$ We obtain the Aspinwall-Morrison formula $$n_d=\frac{1}{d^3},$$ which has been proved by several different methods ,,. We also obtain $$_{\overline{M}_{0,1}(^1,d)}c𝔼_d=\frac{2}{d^3}.$$ ### 6.2. Virtual numbers of plane curves Let $`X`$ be a Calabi-Yau threefold containing a $`^2`$. As we saw in Remark $`\mathrm{4.2.1}`$, the normal bundle of $`^2`$ in $`X`$ is $`K_^2=𝒪(3)`$. Let $`C`$ be a rational curve of degree $`d`$ in $`^2`$. Since $`K_X=𝒪_X`$, the expected dimension of the moduli space $`\overline{M}_{0,0}(X,[C])`$ is zero. Lemma 4.2.1 says that $`\overline{M}_{0,0}(^2,d)=\overline{M}_{0,0}(X,[C])`$, hence the dimension of this moduli stack is $`3d1`$. Recall the diagram $$\begin{array}{ccc}\overline{M}_{0,1}(^2,d)& \stackrel{e_1}{}& ^2\\ \pi _1& & \\ \overline{M}_{0,0}(^2,d)\end{array}$$ From Lemma 4.2.2, the virtual fundamental class of $`\overline{M}_{0,0}(X,[C])`$ is the refined top Chern class of the bundle $`V_d=R^1\pi _{1}^{}{}_{}{}^{}(e_1^{}(K_^2))`$ over $`\overline{M}_{0,0}(^2,d)`$. The zero pointed Gromov-Witten invariant: $$N_d:=\text{deg}[\overline{M}_{0,0}(X,[C])]^{\text{virt}}=_{\overline{M}_{0,0}(^2,d)}𝔼_d$$ is called the virtual number of degree $`d`$ rational curves in $`X`$. As promised in Remark $`\mathrm{3.2.1}`$, we will show that the modified equivariant quantum product in this case has a nonequivariant limit. We will also use The Mirror Theorem to calculate these numbers $`N_d`$. The modified pairing on $`_{^{}}^2`$ corresponding to $`V=𝒪_^2(3)`$ is $$a,b:=_{(^2)_{^{}}}ab\left(\frac{1}{3p\lambda }\right).$$ Recall that $`p`$ denotes the equivariant hyperplane class in $`^2`$. Then $`1,p,p^2`$ is a basis for $``$ as a $`(\lambda )`$-module. A simple calculation shows that $`\lambda p^2,3p^2\lambda H,3p\lambda `$ is its dual basis with respect to the above pairing. Since both bases and $`E_d`$ are polynomials in $`\lambda `$, we can restrict $`\stackrel{~}{I}_d`$ and $`_V`$ in $`𝒫=H^{}(^2,[\lambda ])`$ and take the nonequivariant limit. Denote by $`H`$ the nonequivariant limit of $`p`$. We obtain the following nonequivariant quantum product on $`H^{}^2[[q]]`$ $$a_Vb:=ab+\underset{d=1}{\overset{\mathrm{}}{}}q^dT^kI_d(a,b,3HT_k)$$ where $`T^0=1,T^1=H,T^2=H^2`$ and for $`\gamma _1,\gamma _2,\mathrm{},\gamma _nH^{}^2`$ (108) $$I_d(\gamma _1,\gamma _2,\mathrm{},\gamma _n)=_{\overline{M}_{0,n}(^2,d)}e_1^{}\gamma _1e_2^{}\gamma _2\mathrm{}e_n^{}\gamma _n𝔼_d.$$ For example, using the divisor axiom we obtain $$H_VH=H^2(13\underset{d>0}{}q^dd^3N_d).$$ Theorem 3.2.1 implies the following: ###### Theorem 6.2.1. $`(H^{}^2[[q]],_V)`$ is an associative, commutative and unital ring with unity $`1=[^2]`$. Denote by $`i`$ the embedding $`i:^2X.`$ Since the normal bundle of $`^2`$ in $`X`$ is $`𝒪_^2(3)`$, it follows that $`i^{}({\displaystyle \frac{1}{3}}[^2])=T^1`$ and $`i^{}({\displaystyle \frac{1}{3}}[l])=T^2.`$ Therefore the map $`i^{}:(H^{}X,)(H^{}^2,)`$ is surjective. Consider the small quantum cohomology rings $`SQH^{}X=(H^{}X[[\beta ]],)`$ and $`SQH_V^{}^2:=(H^{}^2[[q]],_V)`$ where the products are given by three point correlators. Recall from section $`4.2`$ the extension of $`i^{}:H^{}(X,)H^{}(^2,)`$ to $`\stackrel{~}{i^{}}:SQH^{}XSQH_V^{}^2.`$ There is a natural relation between the modified quantum product in $`^2`$ and the pure product in $`X`$. ###### Theorem 6.2.2. The map $`\stackrel{~}{i^{}}`$ is a ring homomorphism. Proof. Complete $`\tau ^0=[X],\tau ^1=\frac{1}{3}[^2],\tau ^2=\frac{1}{3}H`$ into a basis of $`(H^{}X,)`$ by adding elements from $`\text{Ker}(i^{})`$. Let $`\tau _0=[pt],\tau _1=H,\tau _2=[^2],\mathrm{}`$ be the dual basis. Let $`a,bH^{}X`$. We want to show $$\stackrel{~}{i}^{}(ab)=i^{}(a)_Vi^{}(b).$$ But $$ab=\underset{\beta H_2(X,)}{}\underset{r}{}q^\beta \tau ^r_{[\overline{M}_{0,3}(X,\beta )]^{\text{virt}}}e_1^{}ae_2^{}be_3^{}\tau _r.$$ Note that this formula is true for a $``$-basis, but due to the uniqueness of the quantum product, it is true for any $``$-basis as well. Now, $$\stackrel{~}{i^{}}(ab)=\underset{d0}{}\underset{r}{}q^di^{}\tau ^r_{\overline{M}_{0,3}(^2,d)}e_1^{}(i^{}a)e_2^{}(i^{}b)e_3^{}(i^{}\tau _r)𝔼_d.$$ But $`i^{}(\tau ^r)=T^k`$ for $`r=0,1,2`$ and $`i^{}(\tau ^r)=0`$ for $`r2`$. The theorem follows from the readily checked fact: $`i^{}(\tau _k)=3HT_k`$ for $`k=0,1,2`$.$``$ Using the divisor and fundamental class properties of the modified gravitational descendants it is easy to show that: $$J_V=\text{exp}\left(\frac{t_0+t_1H}{\mathrm{}}\right)\left(13\frac{H^2}{\mathrm{}^2}\underset{d=1}{\overset{\mathrm{}}{}}q^ddN_d\right).$$ The hypergeometric series corresponding to the total space of $`V=𝒪_^2(3)`$ is: $$I_V:=\text{exp}\left(\frac{t_0+t_1H}{\mathrm{}}\right)\underset{d=0}{\overset{\mathrm{}}{}}q^d\frac{_{m=0}^{3d1}(3Hm\mathrm{})}{_{m=1}^d(H+m\mathrm{})^3}.$$ We expand this function $$I_V=\text{exp}\left(\frac{t_0+t_1H}{\mathrm{}}\right)\left(1+I_1\frac{H}{\mathrm{}}+o(\frac{1}{\mathrm{}})\right)$$ where $$I_1=3\underset{d=1}{\overset{\mathrm{}}{}}q^d(1)^d\frac{(3d1)!}{(d!)^3}.$$ The mirror theorem for this case says that $`J(t_0,t_1+I_1)=I_V(t_0,t_1)`$. This theorem allows us to compute the virtual number of rational plane curves in the Calabi-Yau $`X`$. The first few numbers are $`3,{\displaystyle \frac{45}{8}},{\displaystyle \frac{244}{9}}`$. Department of Mathematics and Statistics American University 4400 Massachusetts Ave Washington, Dc 20016 aelezi@american.edu
warning/0004/hep-ph0004151.html	ar5iv	text	# The Super - Kamiokande Day - Night Effect Data and the MSW Solutions of the Solar Neutrino Problem preprint: Ref. SISSA 32/2000/EP March, 2000 hep – ph/0004151 The current Super-Kamiokande data on the D-N asymmetry between the the day event rate and the Night (Mantle) and Core event rates, produced by solar neutrinos which respectively cross the Earth along any trajectory (cross the Earth mantle but do not cross the core), and cross the Earth core before reaching the detector, imply rather stringent constraints on the MSW small mixing angle (SMA) $`\nu _e\nu _{\mu (\tau )}`$ solution of the solar neutrino problem. A simplified analysis shows, in particular, that a substantial subregion of the SMA solution region is disfavored by these data. The Core D-N asymmetry data alone allow to rule out at 99.7% C.L. a part of this subregion. The constraints on the MSW large mixing angle and LOW $`\nu _e\nu _{\mu (\tau )}`$ solutions as well as on the MSW $`\nu _e\nu _s`$ solution, following from the data on the Mantle, Night and Core D-N asymmetries are also discussed. I. Introduction: The MSW Solutions of the Solar Neutrino Problem and the Day-Night Effect The hypothesis of MSW transitions of solar neutrinos continues to provide a viable solution of the solar neutrino problem . The current (mean event rate) solar neutrino data admit three types of MSW $`\nu _e\nu _{\mu (\tau )}`$ transition solutions: the well-known small mixing angle (SMA) non-adiabatic and large mixing angle (LMA) adiabatic (see, e.g., ) and the so-called “LOW” solution (very recent analyses can be found in, e.g., ). While the SMA and LMA solutions have been shown to be rather stable with respect to variations in the values of the various physical quantities which enter into the calculations (the fluxes of <sup>8</sup>B and <sup>7</sup>Be neutrinos, nuclear reaction cross-sections, etc.), and of the data utilized in the analyses, the LOW solution is of the “borderline” type: its existence even at 99% C.L. is not stable with respect to relatively small changes in the data and/or in the relevant theoretical predictions (see, e.g., and the references quoted therein.). To the three solutions there correspond (at a given C.L.) three distinct regions in the plane of values of the two parameters, $`\mathrm{\Delta }m^2`$ and $`\mathrm{sin}^22\theta `$, characterizing the transitions. One finds using the standard solar model predictions for the solar neutrino fluxes (<sup>8</sup>B, <sup>7</sup>Be, $`pp`$, etc.) that at 99% C.L. the SMA MSW solution requires values in the intervals $`4.0\times 10^6\mathrm{eV}^2<\mathrm{\Delta }m^2<10.0\times 10^6\mathrm{eV}^2`$, $`1.3\times 10^3<\mathrm{sin}^22\theta <1.0\times 10^2`$, the LMA solutions is realized for $`\mathrm{\Delta }m^2`$ and $`\mathrm{sin}^22\theta `$ from the region $`7.0\times 10^6\mathrm{eV}^2<\mathrm{\Delta }m^2<2.0\times 10^4\mathrm{eV}^2`$, $`0.50<\mathrm{sin}^22\theta <1.0`$, and the LOW solution lies approximately in the region $`0.4\times 10^7\mathrm{eV}^2<\mathrm{\Delta }m^2<1.5\times 10^7\mathrm{eV}^2`$, $`0.80<\mathrm{sin}^22\theta <1.0`$. The SMA and LMA solution regions expand in the direction of smaller values of $`\mathrm{sin}^22\theta `$ up to $`0.6\times 10^3`$ and to $`0.3`$, respectively, if one adopts a more conservative approach in analyzing the data in terms of the MSW effect and treats the <sup>8</sup>B neutrino flux as a free parameter in the analysis (see, e.g., ). A unique testable prediction of the MSW solutions of the solar neutrino problem is the day-night (D-N) effect - a difference between the solar neutrino event rates during the day and during the night, caused by the additional transitions of the solar neutrinos taking place at night while the neutrinos cross the Earth on the way to the detector (see, e.g., and the references quoted therein). The experimental observation of a non-zero D-N asymmetry $$A_{DN}^N\frac{R_NR_D}{(R_N+R_D)/2},$$ (1) where $`R_N`$ and $`R_D`$ are, e.g., the one year averaged event rates in a given detector respectively during the night and the day, would be a very strong evidence in favor (if not a proof) of the MSW solution of the solar neutrino problem. Extensive predictions for the magnitude of the D-N effect for the Super-Kamiokande detector have been obtained in . Earlier results have been derived in . High precision calculations of the D-N asymmetry in the one year averaged recoil-e<sup>-</sup> spectrum and in the energy-integrated event rates were performed for three event samples, Night, Core and Mantle, in . The night fractions of these event samples are due to neutrinos which respectively cross the Earth along any trajectory, cross the Earth core, and cross only the Earth mantle (but not the core), on the way to the detector. The measurement of the D-N asymmetry in the Core sample was found to be of particular importance because of the strong enhancement of the asymmetry, caused by a constructive interference between the amplitudes of the neutrino transitions in the Earth mantle and in the Earth core . The effect differs from the MSW one . The mantle-core enhancement effect is caused by the existence (for a given neutrino trajectory through the Earth core) of points of resonance-like total neutrino conversion in the corresponding space of neutrino oscillation parameters . The location of these points determines the regions where the relevant probability of transitions in the Earth of the Earth-core-crossing solar neutrinos is large Being a constructive interference effect between the amplitudes of neutrino transitions in the mantle and in the core, this is not just “core enhancement” effect, but rather mantle-core enhancement effect. . At small mixing angles and in the case of $`\nu _e\nu _{\mu (\tau )}`$ transitions the predicted D-N asymmetry in the Core sample of the Super-Kamiokande event rate data was shown to be much bigger due to the mantle-core enhancement effect The term “neutrino oscillation length resonance” (NOLR) was used in , in particular, to denote the enhancement in this case. \- by a factor of up to $`6`$, than the asymmetry in the Night sample. The asymmetry in the Mantle sample was found to be smaller than the asymmetry in the Night sample. On the basis of these results it was concluded in that it can be possible to test a substantial part of the MSW $`\nu _e\nu _{\mu (\tau )}`$ SMA solution region in the $`\mathrm{\Delta }m^2\mathrm{sin}^22\theta `$ plane by performing selective, i.e., Core and Night (or Mantle) D-N asymmetry measurements. The current Super-Kamiokande data shows a D-N asymmetry in the Night sample, which is different from zero at 1.9 s.d. level: $$A_{DN}^N=0.065\pm 0.031(stat.)\pm 0.013(syst.).$$ (2) These data allow to probe only a relatively small subregion of the SMA “conservative” solution region: the predicted asymmetry is too small (see, e.g., ). However, the Super-Kamiokande night data is given in 5 bins and 80% of the events in the bin N5 are due to Earth-core-crossing solar neutrinos , while the remaining 20% are produce by neutrinos which cross only the Earth mantle. Since the predicted D-N asymmetry in the Mantle sample is practically negligible in the case of the MSW SMA solution of interest , we have for the D-N asymmetry measured using the night N5 bin data: $`A_{DN}^{N5}0.8A_{DN}^C`$, $`A_{DN}^C`$ being the asymmetry in the Core sample. The data on $`A_{DN}^{N5}`$ permitted to exclude a part of the MSW SMA solution region located in the area $`\mathrm{sin}^22\theta (0.0070.01)`$, $`\mathrm{\Delta }m^2(0.51.0)\times 10^5\mathrm{eV}^2`$. It should be obvious from the above discussion that the measurement of the Core asymmetry $`A_{DN}^C`$, as suggested in , will provide a more effective test of the the MSW SMA solution than the measurement of $`A_{DN}^{N5}`$. Recently the Super-Kamiokande collaboration for the first time published data on the Core D-N asymmetry $`A_{DN}^C`$ : $$A_{DN}^C=0.0175\pm 0.0622(stat.)\pm 0.013(syst.).$$ (3) The experimental value of the Mantle asymmetry can also be deduced from the data : $$A_{DN}^M=0.0769\pm 0.034(stat.)\pm 0.013(syst.).$$ (4) In the present article we use the Super-Kamiokande results on the D - N effect, eqs. (2), (3) and (4), to derive constraints on the MSW $`\nu _e\nu _{\mu (\tau )}`$ transition solutions of the solar neutrino problem. We show below, in particular, that, as has been suggested in , the data on the Core and Night or Mantle D - N asymmetries allow to perform a very effective test of the MSW SMA solution. We also obtain constraints on the MSW LMA and LOW $`\nu _e\nu _{\mu (\tau )}`$ solutions as well as on the MSW $`\nu _e\nu _s`$ solution, following from the Super-Kamiokande data on the Mantle, Night and Core D-N asymmetries. II. Constraints on the MSW Solutions from the Super-Kamiokande Data on the Core and Mantle or Night Day-Night Asymmetries We use the high precision methods of calculation of the energy-integrated one year average Core, Night and Mantle D-N asymmetries developed for our earlier studies of the D-N effect for the Super-Kamiokande detector, which are described in detail in . The cross section of the $`\nu _ee^{}`$ elastic scattering reaction was taken from . We used in our calculations the <sup>8</sup>B neutrino spectrum derived in . The probability of survival of the solar $`\nu _e`$ when they travel in the Sun and further to the surface of the Earth, $`\overline{P}_{}(\nu _e\nu _e)`$, was computed on the basis of the analytic expression obtained in and using the method developed in . In the calculation of $`\overline{P}_{}(\nu _e\nu _e)`$ we have utilized (as in ) the density profile of the Sun and the <sup>8</sup>B neutrino production distribution in the Sun, predicted in . These predictions were updated in , but a detailed test study showed (see also ) that using the results of instead of those in leads to a change of the survival probability by less than 1% for any set of values of the parameters $`\mathrm{\Delta }m^2`$ and $`\mathrm{sin}^22\theta `$, relevant for the calculation of the D-N effect. As in , the Earth matter effects were calculated using the Stacey model from 1977 for the Earth density distribution. The latter practically coincides with that predicted by the more recent Earth model . Our results are obtained for the standard values of the electron fraction number in the Earth mantle and the Earth core, $`Y_e^{man}=0.49`$ and $`Y_e^c=0.467`$, which reflect the chemical composition of the two major Earth structures (see, e.g, ). In our simplified analysis we use the Super-Kamiokande data on the Core and Mantle or Night D-N asymmetries, treating the Core and Mantle, and the Core and Night asymmetries as two pairs of independent observables. While this is certainly justified in the case of the Core and Mantle asymmetries, the treatment of the Core and Night asymmetries neglects the possible (weak) correlation between the values of the two asymmetries. Note, however, that i) since such a correlation does not exist in the case of the Core and Mantle asymmetries and ii) the statistics in the Night sample of Super-Kamiokande events is dominated by that in the Mantle sample, we expect that the results obtained using the data on Core and Night asymmetries without taking into account the correlation between the two will not differ substantially from those derived by accounting for the correlation. The indicated results should not differ much also from those found by using the data on the Core and Mantle asymmetries. Actually, the restrictions following from the latter set of data turned out to be somewhat stronger that those derived by us on the basis of Core and Night asymmetry data. The results of our study are presented graphically in Figs. 1a - 1c and 2a - 2c. In Figs. 1a - 1c (2a - 2c) the “conservative” MSW SMA (LMA and LOW) solution region in the $`\mathrm{\Delta }m^2\mathrm{sin}^22\theta `$ plane of oscillation parameters is shown (grey area) together with the regions allowed by the Super-Kamiokande data on the Core, Night and Mantle D-N asymmetries at 1.0 s.d. (a), 1.5 s.d. (b) and 2.0 s.d. (c). The MSW solution regions were obtained using the mean event rate solar neutrino data . Contours of constant Core D-N asymmetry in the plane of the two parameters are also shown (thin solid lines), with the value of the asymmetry corresponding to a given contour indicated on the contour. The regions allowed by the Core asymmetry data in Figs. 1a - 1c are located to the left of the thick (black) solid line, while in Figs. 2a - 2c these regions are situated above the upper and below the lower thick (black) solid lines. The regions allowed by the Night asymmetry data in Figs. 1a - 1c are between the two dashed lines, in Figs. 2a - 2b they are between the two upper-most and between the two lower-most dashed lines, and in Fig. 2c they are above the upper and below the lower dashed lines. Finally, the regions allowed by the Mantle asymmetry data in Fig. 1a are between the two dash-dotted lines, in Figs. 1b - 1c they are located to the right of the dash-dotted line; in Figs. 2a - 2c these regions are situated between the two upper-most and between the two lower-most dash-dotted lines. In Figs. 1c and 2c we have shown also the contour corresponding to the maximal allowed value of $`A_{DN}^C`$ at 3 s.d. (the double-thick (black) solid lines). In obtaining the allowed intervals of values of $`A_{DN}^N`$, $`A_{DN}^T`$ and $`A_{DN}^C`$ at a given C.L. we have added the errors in eqs. (2), (3) and (4) in quadratures. For the Core asymmetry this procedure gives the following 1 s.d., 2 s.d. and 3 s.d. allowed intervals of values: $`A_{DN}^C=[0.081,0.046];[0.145,0.110];[0.208,0.173]`$. The corresponding intervals of allowed values for the Night and Mantle asymmetries are: $`A_{DN}^N=[0.031,0.099];[0.002,0.132];[0.036,0.166]`$, and $`A_{DN}^M=[0.041,0.113];[0.005,0.149];[0.031,0.185]`$. Since the minimal value of the predicted Core D-N asymmetry in the case of the SMA solution is $`(0.03)`$ and the lower bound on $`A_{DN}^C`$ following from the current Super-Kamiokande data at 1 s.d. is $`min(A_{DN}^C)(0.081)`$, it is clear that the latter will not play a role in constraining the SMA solution region. The current Super-Kamiokande upper bound on $`A_{DN}^C`$, however, can be used to constrain the MSW SMA solution region. As Fig. 1a demonstrates, the MSW SMA “conservative” solution region is incompatible with Super-Kamiokande data on the asymmetries $`A_{DN}^N`$ or $`A_{DN}^T`$ and $`A_{DN}^C`$, eqs. (2) or (4) and (3), including 1 s.d. uncertainties. The same conclusion is practically valid at 1.5 s.d. as well, as Fig. 1b shows: the region compatible with the data on the Night and Core asymmetries in this case is a tiny “triangle” around the point $`\mathrm{\Delta }m^24\times 10^6\mathrm{eV}^2`$ and $`\mathrm{sin}^22\theta 0.0085`$, while the Mantle and Core asymmetry data leave allowed only a small point-like region at $`\mathrm{\Delta }m^23.9\times 10^6\mathrm{eV}^2`$ and $`\mathrm{sin}^22\theta 0.0091`$. At 2 s.d. a large subregion of the SMA “conservative” solution region is ruled out by the Super-Kamiokande data on the Core and Night D-N asymmetries: the allowed region in Fig. 1c is located between the thick dashed line corresponding to $`min(A_{DN}^N)=(0.002)`$ and the dash-dotted line corresponding to $`max(A_{DN}^C)=0.110`$. The allowed region is even smaller if one uses the data on the Core and Mantle asymmetries: this is the triangular shaped region to the right of the dash-dotted line and to the left of the thick (black) solid line. At 3 s.d. level only the data on the Core asymmetry constraints the SMA solution region (Fig. 1c): the subregion located to the right of the contour corresponding to $`max(A_{DN}^C)=0.173`$ (double-thick (black) solid line) is ruled out by these data. Most of the indicated subregion (not to mention the results described above) could not be probed at the indicated C.L. by the earlier Super-Kamiokande data on the asymmetry $`A_{DN}^{N5}`$ associated with the night N5 bin data (see, e.g., ). The combined use of the Super-Kamiokande data on the Core and Mantle (Night) D-N asymmetries is less effective in constraining the MSW LMA solution because the mantle-core enhancement effect leads to a modest amplification of the Core asymmetry with respect to the Night (Mantle) one in the LMA solution region. The same conclusion is valid for the LOW solution region. The results of our analysis for the LMA and LOW solutions are depicted in Figs. 2a - 2c. It is interesting that, as Fig. 2a shows, large parts of both the LMA and the LOW solution regions are incompatible at 1 s.d. with the Super-Kamiokande data on the asymmetries $`A_{DN}^C`$ and $`A_{DN}^N`$: the only allowed regions lie in the narrow bands corresponding to $`\mathrm{\Delta }m^2(4.05.0)\times 10^5\mathrm{eV}^2`$ and to $`\mathrm{\Delta }m^2(1.02.0)\times 10^7\mathrm{eV}^2`$. Even these narrow bands are barely compatible with the allowed values of $`A_{DN}^M`$ at 1 s.d. (Fig. 2a). Similar conclusion is valid if one uses the D-N asymmetry data, eqs. (2) - (4), including the 1.5 s.d. uncertainties, as Fig. 2b shows: the allowed values of $`A_{DN}^C`$, $`A_{DN}^M`$ and $`A_{DN}^N`$ are compatible with the subregions located between the upper-most dash-dotted and the thick solid lines, and between the lower-most dash-dotted and the thick solid lines. The allowed values of $`\mathrm{\Delta }m^2`$ in the case of the LMA solution are confined to the interval $`\mathrm{\Delta }m^2(2.26.5)\times 10^5\mathrm{eV}^2`$. The region of the LOW solution corresponding to $`\mathrm{\Delta }m^29.5\times 10^8\mathrm{eV}^2`$ is incompatible at the indicated C.L. by the current Super-Kamiokande data on the asymmetry $`A_{DN}^M`$. These data are considerably less restrictive at 2 s.d. (Fig. 2c): the regions between the upper-most dash-dotted and thick solid lines, and between the lower-most dash-dotted and thick solid lines are compatible with the data. The data on the Core asymmetry alone including the 3 s.d. uncertainties (see Fig. 2c) are incompatible with the subregion of the LMA “conservative” solution region located at $`\mathrm{\Delta }m^210^5\mathrm{eV}^2`$. The solar neutrino data can also be explained assuming that the solar neutrinos undergo MSW transitions into sterile neutrino in the Sun: $`\nu _e\nu _s`$ (see, e.g, ). In this case only a SMA solution is compatible with the data. The corresponding solution region obtained (at a given C.L.) using the mean event rate solar neutrino data and the predictions of ref. for the different solar neutrino flux components practically coincides in magnitude and shape with the SMA $`\nu _e\nu _{\mu (\tau )}`$ solution region, but is shifted by a factor of $`1.2`$ along the $`\mathrm{\Delta }m^2`$ axis to smaller values of $`\mathrm{\Delta }m^2`$. The “conservative” $`\nu _e\nu _s`$ transition solution region extends both in the direction of smaller and larger values of $`\mathrm{sin}^22\theta `$ down to $`0.7\times 10^3`$ and up to 0.4 . The main contribution to the energy-integrated D-N asymmetries in the three Super-Kamiokande event samples of interest comes from the “high” energy tail of the <sup>8</sup>B neutrino spectrum . This is related to the fact the relevant neutrino effective potential difference in the Earth matter in the case of the $`\nu _e\nu _s`$ transitions is approximately by a factor of 2 smaller than in the case of $`\nu _e\nu _{\mu (\tau )}`$ transitions . As a consequence, the predicted Core and Night (Mantle) D-N asymmetries for the Super-Kamiokande detector are considerably smaller in the case of the MSW $`\nu _e\nu _s`$ transition solution than in the case of the MSW SMA $`\nu _e\nu _{\mu (\tau )}`$ solution. Nevertheless, the current Super-Kamiokande data on the Core asymmetry $`A_{DN}^C`$ including the 2 s.d. uncertainties, as can be shown, excludes the subregion of the MSW $`\nu _e\nu _s`$ “conservative” solution region, located, depending on $`\mathrm{\Delta }m^2`$, approximately at $`\mathrm{sin}^22\theta (0.0160.020)`$ (see Fig. 6 in ). III. Conclusions We have studied the implications of the current Super-Kamiokande data on the Core, Night and Mantle D-N asymmetries, eqs. (2) and (3), for the MSW solutions of the solar neutrino problem. The Super-Kamiokande collaboration published recently for the first time data on the Core asymmetry . Performing a very simplified analysis we have found that practically the whole “conservative” region of the MSW SMA $`\nu _e\nu _{\mu (\tau )}`$ solution in the $`\mathrm{\Delta }m^2\mathrm{sin}^22\theta `$ plane is incompatible with the indicated data if one includes the 1 s.d. and 1.5 s.d. uncertainties, the only exception in the latter case being a point-like region at $`\mathrm{\Delta }m^23.9\times 10^6\mathrm{eV}^2`$ and $`\mathrm{sin}^22\theta 0.0091`$. At 2 s.d. a large subregion of the MSW SMA $`\nu _e\nu _{\mu (\tau )}`$ solution region is incompatible with the D-N effect data, while at 3 s.d. the data on the Core D-N asymmetry alone excludes a non-negligible subregion of the indicated solution region (Fig. 1c). The constraints on the LMA and LOW solution regions from the Super-Kamiokande data on the Core and Night or Mantle D-N asymmetries are somewhat weaker. Nevertheless, both the LMA and the LOW solutions are barely compatible with the Core and Night asymmetry data at 1 s.d.: the subregions allowed by the indicated data are contained in the narrow bands determined by $`\mathrm{\Delta }m^2(4.05.0)\times 10^5\mathrm{eV}^2`$ and $`\mathrm{\Delta }m^2(1.02.0)\times 10^7\mathrm{eV}^2`$ (Fig. 2a). At 1.5 s.d. these data are incompatible with substantial subregions of both solution regions (Fig. 2b), while at 2 s.d. the LMA solution region at $`\mathrm{\Delta }m^22.0\times 10^5\mathrm{eV}^2`$ and the whole LOW solution region are allowed by the Super-Kamiokande Core and Night or Mantle D-N asymmetry data (Fig. 2c). The current Super-Kamiokande data on the asymmetries $`A_{DN}^C`$ and $`A_{DN}^N`$ allows to constrain the MSW $`\nu _e\nu _s`$ “conservative” solution region as well: depending on $`\mathrm{\Delta }m^2`$ the subregion located approximately at $`\mathrm{sin}^22\theta (0.0160.020)`$ is ruled out at 2 s.d. by the data. The simplified analysis we have performed demonstrates, in particular, the remarkable potential which the data on the Night or Mantle and especially on the Core D-N asymmetry have for testing the MSW $`\nu _e\nu _{\mu (\tau )}`$ solutions of the solar neutrino problem. As was suggested in , these data are particularly effective in testing the MSW SMA $`\nu _e\nu _{\mu (\tau )}`$ solution. The Core D-N asymmetry is strongly enhanced in the case of MSW SMA $`\nu _e\nu _{\mu (\tau )}`$ solution by the mantle-core enhancement effect . With the increase of the statistics of the Super-Kamiokande experiment the data on the two indicated D-N asymmetries will allow to probe larger and larger subregions of the SMA solution region. Our results indicate that the data on the Core and Night or Mantle asymmetries can be used to perform rather effective tests of the LMA and LOW $`\nu _e\nu _{\mu (\tau )}`$ solutions as well. Using these data one can probe and constrain also the MSW $`\nu _e\nu _s`$ “conservative” solution region. Similarly, one can use the data on the Core, Night and Mantle D-N asymmetries in the charged current event rate in the SNO detector to perform equally effective tests the MSW solutions of the solar neutrino problem . Acknowledgements. We would like to thank Y. Suzuki and M. Smy for very useful and clarifying correspondence concerning the Super-Kamiokande data on the D-N effect. The work of (S.T.P.) was supported in part by the EC grant ERBFMRX CT96 0090. FIGURE CAPTIONS Figs. 1a - 1c. Constraints on the MSW SMA $`\nu _e\nu _{\mu (\tau )}`$ solution from the Super-Kamiokande data on the Core, Night and Mantle D-N asymmetries, $`A_{DN}^C`$, $`A_{DN}^M`$ and $`A_{DN}^N`$, eqs. (3), (4) and (2). The “conservative” SMA solution region is shown in grey color. The thick solid line (the (left) dash-dotted line) corresponds to $`max(A_{DN}^C)`$ ($`min(A_{DN}^M)`$), the dashed lines correspond to the $`min(A_{DN}^N)`$ and $`max(A_{DN}^N)`$, allowed by the Super-Kamiokande data at 1 s.d. (a), 1.5 s.d. (b) and 2 s.d. (c). The double-thick solid line in Fig. 1c corresponds to the 3 s.d. maximal value of $`A_{DN}^C`$ allowed by the data. Contours of given constant Core D-N asymmetry are also shown (thin solid lines). The region to the left of the thick solid line (double-thick solid line in Fig. 1c) is allowed by the data on $`A_{DN}^C`$ (at 3 s.d.), the region between the two dashed lines is allowed by the data on $`A_{DN}^N`$, while the region between the two dash-dotted lines (a) or to the left of the dash-dotted line (b,c) is allowed by the data on $`A_{DN}^M`$. Figs. 2a - 2c. The same as in figures 1a - 1c for the MSW LMA and LOW $`\nu _e\nu _{\mu (\tau )}`$ transition solutions. The regions allowed by the Core asymmetry data (at 3 s.d.) are located above the upper and below the lower thick solid lines (double-thick solid lines in Fig. 2c). The regions allowed by the Mantle asymmetry data are in Figs. 2a - 2c between the two upper and between the two lower dash-dotted lines, while those allowed by the Night asymmetry data are in Figs. 2a - 2b between the two upper and between the two lower dashed lines, and in Fig. 2c they are above the upper and below the lower dashed lines.
warning/0004/hep-ph0004214.html	ar5iv	text	# The model and the weak angle ## The model and the weak angle The gauge group closest to the $`SU(3)\times SU(2)\times U(1)`$ of the Standard Model one can hope to derive from type I/I string theory in the above context is $`U(3)\times U(2)\times U(1)`$. The first factor arises from three coincident D-branes (“color” branes). An open string with one end on them is a triplet under $`SU(3)`$ and carries the same $`U(1)`$ charge for all three components. Thus, the $`U(1)`$ factor of $`U(3)`$ has to be identified with gauged baryon number. Similarly, $`U(2)`$ arises from two coincident “weak” D-branes and the corresponding abelian factor is identified with gauged weak-doublet number. A priori, one might expect that $`U(3)\times U(2)`$ would be the minimal choice. However, this is not good enough because the hypercharge cannot be expressed as a linear combination of baryon and weak-doublet numbers <sup>2</sup><sup>2</sup>2See nevertheless the comments at the end of this section for a string embedding of the Standard Model based on $`U(3)\times U(2)`$, where the two $`U(1)`$’s are not the baryon and weak-doublet numbers. The model is though unsatisfactory for phenomenological reasons.. Therefore, at least one additional $`U(1)`$ factor corresponding to an extra D-brane (“$`U(1)`$” brane) is necessary in order to accommodate the Standard Model. In principle this $`U(1)`$ brane can be chosen to be independent of the other two collections with its own gauge coupling. To improve the predictability of the model, here we choose to put it on top of either the color or the weak D-branes. In either case, the model has two independent gauge couplings $`g_3`$ and $`g_2`$ corresponding, respectively, to the gauge groups $`U(3)`$ and $`U(2)`$. The $`U(1)`$ gauge coupling $`g_1`$ is equal to either $`g_3`$ or $`g_2`$. Let us denote by $`Q_3`$, $`Q_2`$ and $`Q_1`$ the three $`U(1)`$ charges of $`U(3)\times U(2)\times U(1)`$, in a self explanatory notation. Under $`SU(3)\times SU(2)\times U(1)_3\times U(1)_2\times U(1)_1`$, the members of a family of quarks and leptons have the following quantum numbers: $`Q`$ $`(\mathrm{𝟑},\mathrm{𝟐};1,w,0)_{1/6}`$ $`u^c`$ $`(\overline{\mathrm{𝟑}},\mathrm{𝟏};1,0,x)_{2/3}`$ $`d^c`$ $`(\overline{\mathrm{𝟑}},\mathrm{𝟏};1,0,y)_{1/3}`$ $`L`$ $`(\mathrm{𝟏},\mathrm{𝟐};0,1,z)_{1/2}`$ $`l^c`$ $`(\mathrm{𝟏},\mathrm{𝟏};0,0,1)_1`$ Here, we normalize all $`U(N)`$ generators according to $`\mathrm{Tr}T^aT^b=\delta ^{ab}/2`$, and measure the corresponding $`U(1)_N`$ charges with respect to the coupling $`g_N/\sqrt{2N}`$, with $`g_N`$ the $`SU(N)`$ coupling constant. Thus, the fundamental representation of $`SU(N)`$ has $`U(1)_N`$ charge unity. The values of the $`U(1)`$ charges $`x,y,z,w`$ will be fixed below so that they lead to the right hypercharges, shown for completeness as subscripts. The quark doublet $`Q`$ corresponds necessarily to a massless excitation of an open string with its two ends on the two different collections of branes. The $`Q_2`$ charge $`w`$ can be either $`+1`$ or $`1`$ depending on whether $`Q`$ transforms as a $`\mathrm{𝟐}`$ or a $`\overline{\mathrm{𝟐}}`$ under $`U(2)`$. The antiquark $`u^c`$ corresponds to fluctuations of an open string with one end on the color branes and the other on the $`U(1)`$ brane for $`x=\pm 1`$, or on other branes in the bulk for $`x=0`$. Ditto for $`d^c`$. Similarly, the lepton doublet $`L`$ arises from an open string with one end on the weak branes and the other on the $`U(1)`$ brane for $`z=\pm 1`$, or in the bulk for $`z=0`$. Finally, $`l^c`$ corresponds necessarily to an open string with one end on the $`U(1)`$ brane and the other in the bulk. We defined its $`Q_1=1`$. The weak hypercharge $`Y`$ is a linear combination of the three $`U(1)`$’s:<sup>3</sup><sup>3</sup>3A study of hypercharge embeddings in gauge groups obtained from M-branes was considered in Ref. . In the context of Type I groundstates such embeddings were considered in . $$Y=c_1Q_1+c_2Q_2+c_3Q_3.$$ (2) $`c_1=1`$ is fixed by the charges of $`l^c`$ in eq. (The model and the weak angle), while for the remaining two coefficients and the unknown charges $`x,y,z,w`$, we obtain four possibilities: $`c_2={\displaystyle \frac{1}{2}},c_3={\displaystyle \frac{1}{3}};`$ $`x=1,y=0,z=0,w=1`$ $`c_2={\displaystyle \frac{1}{2}},c_3={\displaystyle \frac{1}{3}};`$ $`x=1,y=0,z=1,w=1`$ $`c_2={\displaystyle \frac{1}{2}},c_3={\displaystyle \frac{2}{3}};`$ $`x=0,y=1,z=0,w=1`$ (3) $`c_2={\displaystyle \frac{1}{2}},c_3={\displaystyle \frac{2}{3}};`$ $`x=0,y=1,z=1,w=1`$ Orientifold models realizing the $`c_3=1/3`$ embedding in the supersymmetric case with intermediate string scale $`M_s10^{11}`$ GeV have been described in . To compute the weak angle $`\mathrm{sin}^2\theta _W`$, we use from eq. (2) that the hypercharge coupling $`g_Y`$ is given by <sup>4</sup><sup>4</sup>4The gauge couplings $`g_{2,3}`$ are determined at the tree-level by the string coupling and other moduli, like radii of longitudinal dimensions. In higher orders, they also receive string threshold corrections.: $$\frac{1}{g_Y^2}=\frac{2}{g_1^2}+\frac{4c_2^2}{g_2^2}+\frac{6c_3^2}{g_3^2},$$ (4) with $`g_1=g_2`$ or $`g_1=g_3`$ at the string scale. On the other hand, with the generator normalizations employed above, the weak $`SU(2)`$ gauge coupling is $`g_2`$. Thus, $$\mathrm{sin}^2\theta _W\frac{g_Y^2}{g_2^2+g_Y^2}=\frac{1}{1+4c_2^2+2g_2^2/g_1^2+6c_3^2g_2^2/g_3^2},$$ (5) which for $`g_1=g_2`$ reduces to: $$\mathrm{sin}^2\theta _W(M_s)=\frac{1}{4+6c_3^2g_2^2(M_s)/g_3^2(M_s)},$$ (6) while for $`g_1=g_3`$ it becomes: $$\mathrm{sin}^2\theta _W(M_s)=\frac{1}{2+2(1+3c_3^2)g_2^2(M_s)/g_3^2(M_s)}.$$ (7) We now show that the above predictions agree with the experimental value for $`\mathrm{sin}^2\theta _W`$ for a string scale in the region of a few TeV. For this comparison, we use the evolution of gauge couplings from the weak scale $`M_Z`$ as determined by the one-loop beta-functions of the Standard Model with three families of quarks and leptons and one Higgs doublet, $$\frac{1}{\alpha _i(M_s)}=\frac{1}{\alpha _i(M_Z)}\frac{b_i}{2\pi }\mathrm{ln}\frac{M_s}{M_Z};i=3,2,Y$$ (8) where $`\alpha _i=g_i^2/4\pi `$ and $`b_3=7`$, $`b_2=19/6`$, $`b_Y=41/6`$. We also use the measured values of the couplings at the $`Z`$ pole $`\alpha _3(M_Z)=0.118\pm 0.003`$, $`\alpha _2(M_Z)=0.0338`$, $`\alpha _Y(M_Z)=0.01014`$ (with the errors in $`\alpha _{2,Y}`$ less than 1%) . In order to compare the theoretical relations for the two cases (6) and (7) with the experimental value of $`\mathrm{sin}^2\theta _W=g_Y^2/(g_2^2+g_Y^2)`$ at $`M_s`$, we plot in Fig. 1 the corresponding curves as functions of $`M_s`$. The solid line is the experimental curve. The dashed line is the plot of the function (6) for $`c_3=1/3`$ while the dotted-dashed line corresponds to the function (7) for $`c_3=2/3`$. Thus, the second case, where the $`U(1)`$ brane is on top of the color branes, is compatible with low energy data for $`M_s68`$ TeV. This selects the last two possibilities of charge assignments in Eq. (3). Indeed, the curve corresponding to $`g_1=g_3`$ and $`c_3=1/3`$ intersects the experimental curve for $`\mathrm{sin}^2\theta _W`$ at a scale $`M_s`$ of the order of a few thousand TeV. Since this value appears to be too high to protect the hierarchy, it is less interesting and is not shown in Fig. 1. The other case, where the $`U(1)`$ brane is on top of the weak branes, is not interesting either. For $`c_3=2/3`$, the corresponding curve does not intersect the experimental one at all and is not shown in the Fig. 1, while the case of $`c_3=1/3`$ leads to $`M_s`$ of a few hundred GeV and is most likely excluded experimentally. In the sequel we shall restrict ourselves to the last two possibilities of Eq. (3). From the general solution (3) and the requirement that the Higgs doublet has hypercharge $`1/2`$, one finds the following possible assignments for it, in the notation of Eq. (The model and the weak angle): $`c_2={\displaystyle \frac{1}{2}}:`$ $`H(\mathrm{𝟏},\mathrm{𝟐};0,1,1)_{1/2}`$ $`H^{}(\mathrm{𝟏},\mathrm{𝟐};0,1,0)_{1/2}`$ (9) $`c_2={\displaystyle \frac{1}{2}}:`$ $`\stackrel{~}{H}(\mathrm{𝟏},\mathrm{𝟐};0,1,1)_{1/2}`$ $`\stackrel{~}{H}^{}(\mathrm{𝟏},\mathrm{𝟐};0,1,0)_{1/2}`$ (10) It is straightforward to check that the allowed (trilinear) Yukawa terms are: $`c_2={\displaystyle \frac{1}{2}}:H^{}Qu^c,H^{}Ll^c,H^{}Qd^c`$ (11) $`c_2={\displaystyle \frac{1}{2}}:\stackrel{~}{H}^{}Qu^c,\stackrel{~}{H}^{}Ll^c,\stackrel{~}{H}^{}Qd^c`$ (12) Thus, two Higgs doublets are in each case necessary and sufficient to give masses to all quarks and leptons. Let us point out that the presence of the second Higgs doublet changes very little the curves of Fig. 1 and consequently our previous conclusions about $`M_s`$ and $`\mathrm{sin}^2\theta _W`$. A few important comments are now in order: (i) The spectrum we assumed in Eq. (The model and the weak angle) does not contain right-handed neutrinos on the branes. They could in principle arise from open strings in the bulk. Their interactions with the particles on the branes would then be suppressed by the large volume of the transverse space . More specifically, conservation of the three U(1) charges allow for the following Yukawa couplings involving the right-handed neutrino $`\nu _R`$: $$c_2=\frac{1}{2}:H^{}L\nu _L;c_2=\frac{1}{2}:\stackrel{~}{H}L\nu _R$$ (13) These couplings lead to Dirac type neutrino masses between $`\nu _L`$ from $`L`$ and the zero mode of $`\nu _R`$, which is naturally supressed by the volume of the bulk. (ii) Implicit in the above was our assumption of three generations (The model and the weak angle) of quarks and lepton in the light spectrum. They can arise, for example, from an orbifold action along the lines of the model described in Ref. . (iii) From Eq. (7) and Fig. 1, we find the ratio of the $`SU(2)`$ and $`SU(3)`$ gauge couplings at the string scale to be $`\alpha _2/\alpha _30.4`$. This ratio can be arranged by an appropriate choice of the relevant moduli. For instance, one may choose the weak branes to extend along one extra dimension transverse to the color branes, with size around twice the string length. Another possibility would be to move slightly off the orientifold point which may be necessary also for other reasons (see discussion towards the end of the paper). (iv) Finally, it should be stressed that the charge assignments (The model and the weak angle) were based on the assumption that the anti-quarks $`u^c`$ and $`d^c`$ arise as excitations of open strings with only one end on the color D-branes. This is not however the only possibility. The fact that the $`\overline{\mathrm{𝟑}}`$ of $`SU(3)`$ can also be obtained as the antisymmetric product of two $`\mathrm{𝟑}`$’s implies that $`u^c`$ and $`d^c`$ may also arise as fluctuations of open strings with both ends on the color branes. Similarly, the anti-lepton $`l^c`$ which is $`SU(2)`$ singlet can be obtained as the antisymmetric product of two doublets and consequently it may arise as a fluctuation of an open string with both ends on the weak branes. In these cases, the quantum numbers of the corresponding particles will be: $$u^c:(\overline{\mathrm{𝟑}},\mathrm{𝟏};2,0,0)_{2/3}d^c:(\overline{\mathrm{𝟑}},\mathrm{𝟏};2,0,0)_{1/3}l^c:(\mathrm{𝟏},\mathrm{𝟏};0,2,0)_1$$ (14) One should then repeat the previous analysis from the beginning, with any combination of the particles $`u^c`$, $`d^c`$ and $`l^c`$ in Eq. (The model and the weak angle) replaced by the corresponding $`u^c`$, $`d^c`$ and $`l^c`$. However, as we argue next, the only physically viable alternative scenario to the one discussed above is just to replace $`l^c`$ by $`l^c`$. First observe that $`d^c`$ cannot be replaced by $`d^c`$. Indeed, this would fix $`c_3=1/6`$ and the $`Q`$ hypercharge would set $`c_2=0`$. It is then easy to see that one cannot satisfy the hypercharge assignments of leptons for either choice of $`l^c`$ or $`l^c`$. Next, let us replace $`u^c`$ by $`u^c`$. This fixes $`c_3=1/3`$. If we keep $`l^c`$, then $`c_1=1`$ and one is left with the first two lines of Eq. (3) as the two possible solutions for $`y,z,w`$ ($`x`$ is absent in this case). From our previous analysis of $`\mathrm{sin}^2\theta _W`$, these solutions are not very interesting since the string scale comes out to be either too low or too high. On the other hand, if we substitute also $`l^c`$ by $`l^c`$, the solutions for the remaining parameters are: (a) the second line of Eq. (3) with $`c_1=1`$ as before, which is uninteresting, and (b) the first line of Eq. (3) with $`c_1`$ undetermined. In this case, the $`Q_1`$ charges of all particles vanish, the corresponding gauge field decouples, and the gauge group becomes effectively $`U(3)\times U(2)`$ with $`Y=Q_2/2Q_3/3`$. At first sight, this seems to be a more economical embedding of the Standard Model than the one based on $`U(3)\times U(2)\times U(1)`$. In this case, the $`g_1`$ term drops from Eq. (4) and the weak angle is given by $`1/\mathrm{sin}^2\theta _W(M_s)=2+2g_2^2(M_s)/3g_3^2(M_s)`$. Unfortunately, comparison with the experimental value of $`\mathrm{sin}^2\theta _W`$ at $`M_Z`$ requires a string scale of the order of $`10^{14}`$ GeV. An additional unsatisfactory feature of the models obtained by replacing $`u^c`$ with $`u^c`$ is the absence of appropriate Yukawa couplings to give masses to the up-quarks. The last case to be examined is the substitution of $`l^c`$ alone by $`l^c`$. This leads to the same four solutions (3) as with $`l^c`$, and thus, to the same conclusions for $`\mathrm{sin}^2\theta _W`$ and $`M_s`$. However, the case with $`c_2=1/2`$ is problematic because the charge assignments do not allow tree-level Yukawa interactions to give masses to the leptons. In the case with $`c_2=1/2`$, the Yukawa couplings of the leptons (12) are slightly modified to $$c_2=\frac{1}{2}:H^{}Ll^c,$$ (15) implying that they acquire masses from the Higgs which gives masses also to the up-quarks. ## Extra $`U(1)`$’s, anomalies and proton stability The model under discussion has three $`U(1)`$ gauge interactions corresponding to the generators $`Q_1`$, $`Q_2`$, $`Q_3`$. From the previous analysis, the hypercharge was shown to be either one of the two linear combinations: $$Y=Q_1\frac{1}{2}Q_2+\frac{2}{3}Q_3.$$ (16) It is easy to see that the remaining two $`U(1)`$ combinations orthogonal to $`Y`$ are anomalous. Indeed, the generic two-parameter generator orthogonal to $`Y`$ is $$\stackrel{~}{Q}=(\pm \frac{\beta }{2}\frac{2\gamma }{3})Q_1+\beta Q_2+\gamma Q_3,$$ (17) which satisfies $`\mathrm{Tr}\stackrel{~}{Q}T_{SU(2)}^2=\pm 5\beta /4\gamma /3`$ and $`\mathrm{Tr}\stackrel{~}{Q}T_{SU(3)}^2=2\beta +3\gamma /2`$ (for $`c_2=1/2`$), or $`3\beta /4+11\gamma /6`$ (for $`c_2=1/2`$); they can both vanish only for $`\beta =\gamma =0`$. We have assumed throughout that this model can be derived as a consistent type I/I string vacuum without additional light states charged under $`U(3)\times U(2)\times U(1)`$. In such a vacuum, the anomalies should be canceled by appropriate shifts of Ramond-Ramond axions in the bulk . Since we have two independent anomalous $`U(1)`$ currents, we need two axions $`a,a^{}`$ that couple to the non-abelian Pontryagin densities with coefficients fixed by the anomalies. The relevant part of the low-energy effective lagrangian can be written as: $`_{eff}^{(1)}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(a+\lambda M_sA)^2+{\displaystyle \frac{1}{32\pi ^2}}{\displaystyle \frac{a}{\lambda M_s}}(k_2\mathrm{Tr}F_2F_2+k_3\mathrm{Tr}F_3F_3)`$ $`+`$ $`{\displaystyle \frac{1}{2}}(a^{}+\lambda M_sA^{})^2+{\displaystyle \frac{1}{32\pi ^2}}{\displaystyle \frac{a^{}}{\lambda M_s}}(k_2^{}\mathrm{Tr}F_2F_2+k_3^{}\mathrm{Tr}F_3F_3),`$ where $`F_2`$ and $`F_3`$ are the $`SU(2)`$ and $`SU(3)`$ gauge field strengths, and $`A`$, $`A^{}`$ the gauge fields corresponding to two independent and mutually orthogonal anomalous abelian charges $`Q_A`$, $`Q_A^{}`$ of the form (17). $`k_2`$, $`k_2^{}`$ ($`k_3`$, $`k_3^{}`$) are their respective mixed anomalies with $`SU(2)`$ ($`SU(3)`$) given by $$k_2=\mathrm{Tr}Q_AT_{SU(2)}^2,k_2^{}=\mathrm{Tr}Q_A^{}T_{SU(2)}^2,k_3=\mathrm{Tr}Q_AT_{SU(3)}^2,k_3^{}=\mathrm{Tr}Q_A^{}T_{SU(3)}^2,$$ (19) while $`\lambda `$ is a calculable parameter in every particular string vacuum. The theory is invariant under the gauge transformation $`AA+\mathrm{\Lambda }/g_A`$, $`aa\lambda M_s\mathrm{\Lambda }`$, together with appropriate transformations of the fermion fields. Indeed, under this transformation the action (Extra $`U(1)`$’s, anomalies and proton stability) changes by exactly the amount necessary to cancel the phase of the fermionic determinant. Ditto for $`A^{}`$. According to Eq. (Extra $`U(1)`$’s, anomalies and proton stability), the gauge fields $`A`$ and $`A^{}`$ become massive with masses $`\lambda g_AM_s`$ and $`\lambda g_A^{}M_s`$, respectively, with $`g_A`$ and $`g_A^{}`$ the corresponding gauge couplings. The axions $`a`$ and $`a^{}`$ become their longitudinal components. Note that we have chosen $`A`$ and $`A^{}`$ so that their mass matrix is diagonal. Gravitational anomalies proportional to the trace of a single charge are also canceled by similar axionic couplings to $`RR`$. This mechanism can be generalized to show the cancellation of the mixed $`U(1)`$ anomalies. These are of two types. First, the ones associated to the non-vanishing traces $`\mathrm{Tr}Q_AY^2k_Y`$ and $`\mathrm{Tr}Q_A^{}Y^2k_Y^{}`$ can be canceled by introducing in $`_{eff}^{(1)}`$ the additional terms $$_{eff}^{(1)}_{eff}^{(1)}+\frac{1}{32\pi ^2}\frac{1}{\lambda M_s}(k_Ya+k_Y^{}a^{})F_YF_Y,$$ (20) which are needed to cancel the $`F_YF_Y`$ contribution to the divergence of the two anomalous currents. The coefficients $`k_Y`$ and $`k_Y^{}`$ can be deduced from the anomaly of the generic current (17). In the case of $`l^c`$ we obtain $`\mathrm{Tr}\stackrel{~}{Q}Y^2=4\beta /343\gamma /18`$ (for $`c_2=1/2`$), or $`\beta /1237\gamma /18`$ (for $`c_2=1/2`$), while for $`l^c`$ (and $`c_2`$ necessarily $`1/2`$) $`\mathrm{Tr}\stackrel{~}{Q}Y^2=7\beta /631\gamma /18`$. Second, there are mixed anomalies related to the non-vanishing trace $`\mathrm{Tr}Y\stackrel{~}{Q}^2=\beta ^2/2+20\gamma ^2/9\beta \gamma /3`$ (for $`c_2=1/2`$), or $`3\beta ^2/4+16\gamma ^2/95\beta \gamma /3`$ (for $`c_2=1/2`$), or $`17\beta ^2/4+16\gamma ^2/9+\beta \gamma /3`$ (in the case of $`l^c`$ for $`c_2=1/2`$). Using this general formula, we can uniquely determine the two orthogonal combinations $`Q_A`$ and $`Q_A^{}`$ in such a way that the triple mixed trace vanishes. We thus have: $$\mathrm{Tr}YQ_A^2\xi _A,\mathrm{Tr}YQ_A^{}^2\xi _A^{},\mathrm{Tr}YQ_AQ_A^{}=0.$$ (21) These mixed anomalies seem to violate the hypercharge gauge invariance of the Standard Model. However, in the context of a consistent string theory, they should also be eliminated. This can be achieved without giving mass to the hypercharge gauge field $`A_Y`$ if the low-energy effective lagrangian contains Chern-Simons terms of the form $`A_Y\omega _A`$ and $`A_Y\omega _A^{}`$. Finally, the violation of the $`U(1)_A`$ and $`U(1)_A^{}`$ gauge invariances introduced by these new terms can be remedied by adding non-diagonal axionic couplings proportional to $`aF_YF_A`$ and $`a^{}F_YF_A^{}`$. To summarize, one may cancel all anomalies of our model by modifying the relevant to the anomaly part of the effective lagrangian $`_{eff}^{(1)}`$ in Eq. (Extra $`U(1)`$’s, anomalies and proton stability) to: $`_{eff}^{\mathrm{anom}}=_{eff}^{(1)}`$ $`+`$ $`{\displaystyle \frac{1}{32\pi ^2}}{\displaystyle \frac{1}{\lambda M_s}}(k_Ya+k_Y^{}a^{})F_YF_Y`$ $``$ $`{\displaystyle \frac{1}{32\pi ^2}}A_Y(\xi _A\omega _A+\xi _A^{}\omega _A^{})+{\displaystyle \frac{1}{32\pi ^2}}{\displaystyle \frac{1}{\lambda M_s}}F_Y(\xi _AaF_A+\xi _A^{}a^{}F_A^{}).`$ For completeness, we give the linear combinations $`Q_A`$ and $`Q_A^{}`$ that satisfy Eq. (21): $`Q_A`$ $``$ $`{\displaystyle \frac{3}{2}}Q_1\pm {\displaystyle \frac{13}{3}}Q_2+Q_3+t({\displaystyle \frac{2}{3}}Q_1+Q_3)`$ $`Q_A^{}`$ $``$ $`t({\displaystyle \frac{3}{2}}Q_1\pm {\displaystyle \frac{13}{3}}Q_2+Q_3)+{\displaystyle \frac{61}{4}}({\displaystyle \frac{2}{3}}Q_1+Q_3)`$ (23) where the $`\pm `$ signs correspond to $`c_2=1/2`$, respectively. For $`c_2=1/2`$ the value of $`t`$ is $`t=(1159\pm 13\sqrt{21533})/388`$. For $`c_2=1/2`$, $`t=(427\pm 13\sqrt{1342})/54`$ for $`l^c`$, and $`t=(671\pm 91\sqrt{61})/60`$ for $`l^c`$. An important property of the above Green-Schwarz anomaly cancellation mechanism is that the two $`U(1)`$ gauge bosons $`A`$ and $`A^{}`$ acquire masses leaving behind the corresponding global symmetries (23) . This is in contrast to what would had happened in the case of an ordinary Higgs mechanism. These global symmetries remain exact to all orders in type I string perturbation theory around the orientifold vacuum. This follows from the topological nature of Chan-Paton charges in all string amplitudes. On the other hand, one expects non-perturbative violation of global symmetries and consequently exponentially small in the string coupling, as long as the vacuum stays at the orientifold point. Once we move sufficiently far away from it, we expect the violation to become of order unity. This can be justified in a supersymmetric theory as follows. Every Ramond-Ramond axion $`a`$ is part of a chiral superfield $`a+im/g_s`$ with $`g_s`$ the string coupling. Its scalar component $`m`$ is a NS-NS (Neuveu-Scwharz) closed string modulus, whose vacuum expectation value (VEV) blows up the orbifold singularities moving away from the orientifold point. Using the shift of the axion under gauge transformations, one can form the complex field $`e^{(iam/g_s)/M_s}`$ that transforms covariantly with charge $`\lambda `$. A matter interaction term with charge $`n\lambda `$ (with integer $`n`$), multiplied by the $`n`$-th power of this field forms a neutral operator which can appear in the effective action. For $`<m>0`$, one thus obtains charge violating non-perturbative interaction terms with strength $`𝒪(e^{<m>/g_sM_s})`$. A small $`<m>`$ of order $`g_sM_s`$ leads therefore to charge violations of order unity. So, as long as we stay at the orientifold point, all three charges $`Q_1`$, $`Q_2`$, $`Q_3`$ are conserved and since $`Q_3`$ is the baryon number, proton stability is guaranteed. To break the electroweak symmetry, the Higgs doublets in Eq. (9) or (10) should acquire non-zero VEV’s. Since the model is non-supersymmetric, this may be achieved radiatively . From Eqs. (11) and (12), to generate masses for all quarks and leptons, it is necessary for both Higgses to get non-zero VEV’s. The baryon number conservation remains intact because both Higgses have vanishing $`Q_3`$. However, the linear combination $`(t61/4)Q_A+(t+1)Q_A^{}`$ which does not contain $`Q_3`$, will be broken spontaneously, as follows from their quantum numbers in Eqs. (9) and (10). This leads to an unwanted massless Goldstone boson of the Peccei-Quinn type. A possible way out is to break this global symmetry explicitly, by moving away from the orientifold point along the direction of non-vanishing $`(t61/4)m+(t+1)m^{}`$, so that baryon number remains conserved. In conclusion, we presented a particular embedding of the Standard Model in a non-supersymmetric D-brane configuration of type I/I string theory. The strong and electroweak couplings are not unified because strong and weak interactions live on different branes. Nevertheless, $`\mathrm{sin}^2\theta _W`$ is naturally predicted to have the right value for a string scale of the order of a few TeV. The model contains two Higgs doublets needed to give masses to all quarks and leptons, and preserves baryon number as a (perturbatively) exact global symmetry. The model satisfies the main phenomenological requirements for a viable low energy theory and its explicit derivation from string theory deserves further study. ## Acknowledgments This work was partly supported by the EU under TMR contract ERBFMRX-CT96-0090 and INTAS contract N 99 0590. I.A. would like to thank the Department of Physics of the University of Crete, while T.N.T. the CPHT of the Ecole Polytechnique, as well as the L.P.T. of the Ecole Normale Supérieure, for their hospitality. E.K. is grateful to C. Bachas, L. Ibáñez, H.P. Nilles, A. Sagnotti and A. Uranga for discussions.
warning/0004/hep-th0004121.html	ar5iv	text	# 1 Introduction ## 1 Introduction Affine Toda field theory is an integrable quantum field theory in two-dimensional Minkowski space-time which possesses remarkable properties and rich algebraic structure (for a review see ). This theory becomes more interesting when it is restricted to a half-line. However, for most of the Toda theories corresponding to affine simply-laced algebras, the boundary conditions are limited to a finite number which preserve integrability. Corrigan et.al have classified the boundary conditions which preserve classical integrability. However, there still remains much to be studied in relation to quantum integrability on the half-line. For models based on $`a_n^{(1)}`$ much is now known . The simplest affine Toda theory, the sinh-Gordon model has been studied much more than other models in the context of integrable boundaries. This model is the only theory in the $`ade`$ series of affine Toda field theory for which continuous boundary parameters are possible. In recent years there has been considerable interest in perturbative affine Toda field theory. The motivation behind this fact is that the boundary S-matrices of the models are largely unknown. The most progress has been made for $`a_1^{(1)}`$ affine Toda field theory for which the general form of the boundary S-matrix has been found by Ghoshal . In fact, the boundary bootstrap equations yield the boundary S-matrices up to some unknown parameters. The perturbation method not only provides an additional check of the results which come from the bootstrap technique, but also it could make a connection between the unknown parameters of the boundary S-matrices and the boundary parameters which are involved in the Lagrangian formulation of the theories. Firstly, Ghoshal and Zamolodchikov obtained the soliton reflection factors in the sine-Gordon model with a boundary consistent with integrability. Then, Ghoshal using these results calculated the reflection factors of the soliton-anti-soliton bound states (the breathers) of the model. One of the interesting problem in the boundary sine(sinh)-Gordon model is to find the relation between the free parameters appearing in Ghoshal’s formula and the boundary data appearing in the Lagrangian of the model. Corrigan was the first to notice that the lightest breather reflection factor of the sine-Gordon model is identical to the reflection factor of the sinh-Gordon model after an analytic continuation in the coupling constant. In a recent paper Corrigan and Delius studied the boundary breather states of the sinh-Gordon model on a half-line. They calculated the energy spectrum of the boundary states in two ways, by using the bootstrap equations then by using a WKB approximation. By comparing the results obtained by the two methods, they provided strong evidence for a conjectured relationship between the boundary parameters, the bulk coupling constant and the parameters appearing in the quantum reflection factor calculated by Ghoshal. They carried out the calculations in the special case when the boundary parameters are equal and the boundary condition preserve the $`\varphi \varphi `$ symmetry of the bulk theory. In the quantum corrections up to $`O(\beta ^2)`$ to the classical reflection factor of the sinh-Gordon model were found when the boundary parameters are equal. In this case, the static background configuration is $`\varphi =0`$. If the boundary data are different then, the lowest energy solution will not be a trivial background. The corresponding perturbation theory involves complicated coupling constants and two-point Green function as well. Recently the quantum reflection factor has been calculated in one loop order up to the first order in the difference of the two boundary parameters. The result provide a further verification of Ghoshal’s formula. This paper extends the results of , by calculating the quantum reflection factor for any value of the boundary parameters. It is found that most part of the related quantum corrections to the classical reflection factor may be expressed in terms of hypergeometric functions. This result and how it could relate to Ghoshal’s formula is discussed in the conclusions. ## 2 Boundary sinh-Gordon model The sinh-Gordon theory on the half-line is a massive scalar quantum field theory in 1+1 dimension whose corresponding untwisted affine Kac-Moody algebra is $`a_1^{(1)}`$. The Lagrangian density of the theory is: $$\overline{}=\theta (x)\delta (x)$$ (2.1) Here, $``$ is the bulk Lagrangian density of the model which is given by $``$ $`=`$ $`{\displaystyle \frac{1}{2}}_\mu \varphi ^\mu \varphi V(\varphi )`$ (2.2) $`=`$ $`{\displaystyle \frac{1}{2}}_\mu \varphi ^\mu \varphi {\displaystyle \frac{2m^2}{\beta ^2}}\mathrm{cosh}(\beta \alpha \varphi ),`$ where m and $`\beta `$ are a mass scale and a coupling constant of the theory. Moreover, the boundary potential $``$ has the generic form $$=\frac{m}{\beta ^2}\left(\sigma _0e^{\frac{\beta }{\sqrt{2}}\varphi }+\sigma _1e^{\frac{\beta }{\sqrt{2}}\varphi }\right).$$ (2.3) In the above relation, the two real coefficients $`\sigma _0`$ and $`\sigma _1`$ are arbitrary and indicate the degrees of freedom allowed at the boundary. In fact, Bowcock et.al obtained some results about the form of the the boundary term via a generalized Lax pair when there is a boundary. For further discussion on the boundary parameters see . The sinh-Gordon model is integrable classically which means there are infinitely many independent conserved quantities $`Q_{\pm s}`$ where s is an arbitrary odd integer. On the other hand, the model is integrable after quantizing which implies the S-matrix describing the n-particles scattering factorises into a product of two-particles scattering amplitudes. The S-matrix describing the elastic scattering of two sinh-Gordon particles of relative rapidity $`\theta `$ is conjectured to have the form $$S(\theta )=\frac{1}{(B)(2B)}$$ (2.4) where we use the hyperbolic building blocks $$(x)=\frac{\mathrm{sinh}(\theta /2+\frac{i\pi x}{4})}{\mathrm{sinh}(\theta /2\frac{i\pi x}{4})},$$ (2.5) and the quantity B is related to the coupling constant $`\beta `$ by $`B=\frac{2\beta ^2}{4\pi +\beta ^2}`$. In order to maintain the integrability on the half-line, the boundary potential must satisfy the following equation $$\frac{\varphi }{x}=\frac{}{\varphi }\text{at}x=0,$$ (2.6) or $$\frac{\varphi }{x}=\frac{\sqrt{2}m}{\beta }\left(\sigma _1e^{\beta \varphi /\sqrt{2}}\sigma _0e^{\beta \varphi /\sqrt{2}}\right)\text{at}x=0,$$ (2.7) where we use the normalization condition $`\alpha ^2=2`$ which is customary in affine Toda field theory. In what follows the dimensional mass parameter $`m`$ will be taken to unity. It is also convenient to use $`\sigma _i=\mathrm{cos}a_i\pi `$. For the boundary sinh-Gordon model, besides to the two-particle S-matrix it is necessary to know the boundary S-matrix or reflection factor describing one particle reflection off the boundary. Firstly, Ghoshal and Zamolodchikov calculated the soliton reflection factors for the sine-Gordon model by solving the boundary Yang-Baxter equation. Then, Ghoshal calculated the soliton-antisoliton bound state reflection factor. He used the boundary bootstrap equations along with the result of reference . The general form of the quantum reflection factor in sinh-Gordon model may be derived by regarding the lightest breather reflection factor of the sine-Gordon model , calculated by Ghoshal, and performing analytic continuation in the coupling constant to find $$K_q(\theta )=\frac{(1)(2B/2)(1+B/2)}{(1E(\sigma _0,\sigma _1,\beta ))(1+E(\sigma _0,\sigma _1,\beta ))(1F(\sigma _0,\sigma _1,\beta ))(1+F(\sigma _0,\sigma _1,\beta ))}.$$ (2.8) Note the bulk reflection symmetry leads to $`F=0`$ or $`E=0`$ when $`\sigma _0=\sigma _1`$ (note only one of them vanishes). In fact, the exact form of the $`E`$ and $`F`$ is an open and hard problem. Recently Corrigan and Delius obtained the function E in the special case when $`\sigma _0=\sigma _1=\mathrm{cos}a\pi `$ and $`\frac{1}{2}<a<1`$ as $$E=2a(1B/2).$$ (2.9) They found the above formula by equating the results of the WKB approximation method and the bootstrap technique. ## 3 Low order perturbation theory For affine Toda field theory the perturbative calculation is performed around the static background field configuration, so standard Feynman Rules may be used. By expanding the bulk and boundary potentials in terms of the coupling constant $`\beta `$, the three and four point couplings can be deduced. We find for the sinh-Gordon theory $$C_{bulk}^{(3)}=\frac{2\sqrt{2}}{3}\beta \mathrm{sinh}(\sqrt{2}\beta \varphi _0)$$ (3.1) $$C_{bulk}^{(4)}=\frac{1}{3}\beta ^2\mathrm{cosh}(\sqrt{2}\beta \varphi _0),$$ (3.2) where $`\varphi _0`$ represent the background solution to the equation of motion of the model and similarly $$C_{boundary}^{(3)}=\frac{\sqrt{2}\beta }{12}\left(\sigma _1e^{\beta \varphi _0/\sqrt{2}}\sigma _0e^{\beta \varphi _0/\sqrt{2}}\right)$$ (3.3) $$C_{boundary}^{(4)}=\frac{\beta ^2}{48}\left(\sigma _1e^{\beta \varphi _0/\sqrt{2}}+\sigma _0e^{\beta \varphi _0/\sqrt{2}}\right).$$ (3.4) On the other hand, the static background field can be found through linear perturbation of the equation of motion and the boundary condition of the model to obtain $$e^{\beta \varphi _0/\sqrt{2}}=\frac{1+e^{2(xx_0)}}{1e^{2(xx_0)}},$$ (3.5) where the parameter $`x_0`$ is related to the boundary parameters by $$\mathrm{coth}x_0=\sqrt{\frac{1+\sigma _0}{1+\sigma _1}}.$$ (3.6) So, the three and four point couplings corresponding to the bulk potential take the forms $$C_{bulk}^{(3)}=\frac{4\sqrt{2}}{3}\beta \mathrm{cosh}2(xx_0)\left(\mathrm{coth}^22(xx_0)1\right),$$ (3.7) $$C_{bulk}^{(4)}=\frac{1}{3}\beta ^2\left(2\mathrm{coth}^22(xx_0)1\right).$$ (3.8) In the same manner the three and four point couplings of the boundary are given by $$C_{boundary}^{(3)}=\frac{\sqrt{2}\beta }{12}\left(\sigma _1\mathrm{coth}x_0\sigma _0\mathrm{tanh}x_0\right),$$ (3.9) $$C_{boundary}^{(4)}=\frac{\beta ^2}{48}\left(\sigma _1\mathrm{coth}x_0+\sigma _0\mathrm{tanh}x_0\right).$$ (3.10) The next step is to find the propagator for the theory. It has been shown that the two-point Green function for the sinh-Gordon model on a half-line is $`G(x,t;x^{},t^{})`$ $`=`$ $`{\displaystyle }{\displaystyle }{\displaystyle \frac{d\omega }{2\pi }}{\displaystyle \frac{dk}{2\pi }}{\displaystyle \frac{ie^{i\omega (tt^{})}}{\omega ^2k^24+i\rho }}(f(k,x)f(k,x^{})e^{ik(xx^{})}`$ (3.11) $`+K_cf(k,x)f(k,x^{})e^{ik(x+x^{})}),`$ where $$f(k,x)=\frac{ik2\mathrm{coth}2(xx_0)}{ik+2}$$ (3.12) and $`K_c`$ is the classical reflection factor of the model which is equal to $$K_c=\left(\frac{(ik)^2+2ik\sqrt{1+\sigma _0}\sqrt{1+\sigma _1}+2(\sigma _0+\sigma _1)}{(ik)^22ik\sqrt{1+\sigma _0}\sqrt{1+\sigma _1}+2(\sigma _0+\sigma _1)}\right)\left(\frac{ik2}{ik+2}\right).$$ (3.13) The classical reflection factor (3.13) can be derived from the quantum reflection factor (2.8) by considering the classical limit i.e. when $`\beta 0`$. Because, in this limit $`E=a_0+a_1`$ and $`F=a_0a_1`$. Now following the idea introduced by Kim and developed by Corrigan , we may calculate the one loop quantum corrections to the classical reflection factor after perturbation calculation of the two-point function and then by finding the coefficient of $`e^{ik(x+x^{})}`$ in the residue of the on-shell pole in the asymptotic region $`x,x^{}\mathrm{}`$. In order to calculate the one loop $`\left(O(\beta ^2)\right)`$ quantum corrections to the classical reflection factor, we use the standard perturbation theory which is generalized to the affine Toda field theory on the half-line. In general, at $`\left(O(\beta ^2)\right)`$ there are three basic kinds of Feynman diagrams contribute to the two-point propagator of affine Toda field theory . These are shown in figure 1. Moreover, by inspection of the forms of the three point and four point couplings which we have found, it is clear that all types of these diagrams are involved in our problem. I II III Figuer 1: Three basic Feynman diagrams in one loop order. In fact, when the boundary parameters are not equal then, the calculations corresponding to the one loop order in the sinh-Gordon model are lengthy and intricate. In the following sections we try to compute the contributions of types I and III diagrams to the reflection factor. The remaining diagrams will be treated elsewhere. Meanwhile, it is instructive to start with type III. ## 4 Type III diagram (boundary-boundary) In this section we shall calculate the contribution of the type III diagram to the reflection factor when both vertices are located at the boundary $`x=0`$. We are led to the following integral $`{\displaystyle \frac{\beta ^2}{4}}(\sigma _1\mathrm{coth}x_0\sigma _0\mathrm{tanh}x_0)^2`$ $`\times {\displaystyle }{\displaystyle }dtdt^{}G(x_1,t_1;0,t)G(0,t;0,t^{})G(0,t^{};0,t^{})G(0,t;x_2,t_2)`$ (4.1) in which the combinatorial factor has been taken into account. Let us start by looking at the loop propagator $`G(0,t^{};0,t^{})`$ which is equal to $`G(0,t^{};0,t^{})`$ $`=`$ $`i{\displaystyle \frac{d\omega ^{}}{2\pi }\frac{dk^{}}{2\pi }\frac{1}{\omega ^2k^24+i\rho }f(k^{},0)}`$ (4.2) $`\times \left(f(k^{},0)+K^{}(k^{})f(k^{},0)\right),`$ where $$f(k^{},0)=\frac{ik^{}+2\mathrm{coth}2x_0}{ik^{}+2}$$ (4.3) and $`K^{}(k^{})`$ is the classical reflection factor (3.13). After some manipulation, we obtain $`G(0,t^{};0,t^{})`$ $`=`$ $`i{\displaystyle \frac{d\omega ^{}}{2\pi }\frac{dk^{}}{2\pi }\frac{1}{\omega ^2k^24+i\rho }}`$ (4.4) $`\times {\displaystyle \frac{2ik^{}\left(ik^{}2\mathrm{coth}2x_0\right)}{(ik^{})^22ik^{}\sqrt{1+\sigma _0}\sqrt{1+\sigma _1}+2(\sigma _0+\sigma _1)}}.`$ The above integral is clearly divergent however, the divergence can be removed by the infinite renormalization of the boundary term. In other words, considering the following relation $`{\displaystyle \frac{2ik^{}\left(ik^{}2\mathrm{coth}2x_0\right)}{(ik^{})^22ik^{}\sqrt{1+\sigma _0}\sqrt{1+\sigma _1}+2(\sigma _0+\sigma _1)}}`$ $`=2+4{\displaystyle \frac{ik^{}\left(\sqrt{1+\sigma _0}\sqrt{1+\sigma _1}\mathrm{coth}2x_0\right)(\sigma _0+\sigma _1)}{(ik^{})^22ik^{}\sqrt{1+\sigma _0}\sqrt{1+\sigma _1}+2(\sigma _0+\sigma _1)}},`$ (4.5) it is seen that a minimal subtraction of the divergent part can be made by adding an appropriate counter term to the boundary, replace the logarithmically divergent integral by the finite part. Hence, $`G(0,t^{};0,t^{})`$ $`=`$ $`4i{\displaystyle \frac{d\omega ^{}}{2\pi }\frac{dk^{}}{2\pi }\frac{1}{\omega ^2k^24+i\rho }}`$ (4.6) $`\times {\displaystyle \frac{ik^{}\left(\sqrt{1+\sigma _0}\sqrt{1+\sigma _1}\mathrm{coth}2x_0\right)(\sigma _0+\sigma _1)}{(ik^{})^22ik^{}\sqrt{1+\sigma _0}\sqrt{1+\sigma _1}+2(\sigma _0+\sigma _1)}}.`$ The integration over $`\omega ^{}`$ may be performed by closing the contour into the upper half-plane and collecting a pole at $`\omega ^{}=\sqrt{k^2+4}`$ so that $$G(0,t^{};0,t^{})=2\frac{dk^{}}{2\pi }\frac{1}{\sqrt{k^2+4}}\frac{ik^{}\left(\sqrt{1+\sigma _0}\sqrt{1+\sigma _1}\mathrm{coth}2x_0\right)(\sigma _0+\sigma _1)}{(ik^{})^22ik^{}\sqrt{1+\sigma _0}\sqrt{1+\sigma _1}+2(\sigma _0+\sigma _1)}.$$ (4.7) In order to integrate over $`k^{}`$, as before, one chooses a contour in the upper half-plane, however due to the branch cut the contour has to run around the cut line. Moreover we assume that the roots of the denominator of the integrand i.e. $`2\mathrm{cos}\frac{(a_0\pm a_1)\pi }{2}`$ are positive, otherwise we may close the contour in the lower half-plane. Therefore, (4.7) converts to $$4_2^{\mathrm{}}\frac{dy}{2\pi }\frac{1}{\sqrt{y^24}}\frac{y\left(\sqrt{1+\sigma _0}\sqrt{1+\sigma _1}\mathrm{coth}2x_0\right)+(\sigma _0+\sigma _1)}{y^2+2y\sqrt{1+\sigma _0}\sqrt{1+\sigma _1}+2(\sigma _0+\sigma _1)}$$ (4.8) and the above integral gives the following result $$G(0,t^{};0,t^{})=\frac{a_0}{2}\mathrm{cot}a_0\pi \frac{a_1}{2}\mathrm{cot}a_1\pi .$$ (4.9) Now it is convenient to calculate the time integral of the other middle propagator in (4) which is equal to $`{\displaystyle 𝑑t^{}G(0,t;0,t^{})}`$ $`=`$ $`{\displaystyle 𝑑t^{}\frac{d\omega }{2\pi }\frac{dk}{2\pi }e^{i\omega (tt^{})}\frac{i}{\omega ^2k^24+i\rho }}`$ (4.10) $`\times {\displaystyle \frac{2ik\left(ik2\mathrm{coth}2x_0\right)}{(ik)^22ik\sqrt{1+\sigma _0}\sqrt{1+\sigma _1}+2(\sigma _0+\sigma _1)}}.`$ Clearly, in the boundary-boundary contribution (4), it is seen that the $`t^{}`$ dependence is involved only in the propagator $`G(0,t;0,t^{})`$ so, the integration over $`t^{}`$ gives us a Dirac delta function which means we must substitute zero for $`\omega `$ and hence $$𝑑t^{}G(0,t;0,t^{})=\frac{dk}{2\pi }\left(\frac{i}{k^24}\right)\left(\frac{2ik\left(ik2\mathrm{coth}2x_0\right)}{(ik)^22ik\sqrt{1+\sigma _0}\sqrt{1+\sigma _1}+2(\sigma _0+\sigma _1)}\right).$$ (4.11) As we mentioned before, throughout this paper we assume that the roots of $`P(k)=(ik)^22ik\sqrt{1+\sigma _0}\sqrt{1+\sigma _1}+2(\sigma _0+\sigma _1)`$ which are equal to $`2\mathrm{cos}\frac{(a_0\pm a_1)\pi }{2}`$ are positive, so the $`P(k)`$ has no pole in the upper half-plane. Obviously, if the roots are negative then we can choose the contour in the lower half-plane in which no pole is inserted. Therefore, $$𝑑t^{}G(0,t;0,t^{})=\frac{i(1+\mathrm{coth}2x_0)}{2+2\sqrt{1+\sigma _0}\sqrt{1+\sigma _1}+(\sigma _0+\sigma _1)}$$ (4.12) and by substituting $`\sigma _0=\mathrm{cos}a_0\pi `$ and $`\sigma _1=\mathrm{cos}a_1\pi `$, we obtain $$𝑑t^{}G(0,t;0,t^{})=\frac{i}{4\mathrm{cos}\frac{a_0\pi }{2}\mathrm{cos}\frac{a_1\pi }{2}}.$$ (4.13) Up to now, the boundary-boundary contribution has the form $`{\displaystyle \frac{i\beta ^2}{32}}{\displaystyle \frac{(\sigma _1\mathrm{coth}x_0\sigma _0\mathrm{tanh}x_0)^2(a_0\mathrm{cot}a_0\pi +a_1\mathrm{cot}a_1\pi )}{\mathrm{cos}\frac{a_0\pi }{2}\mathrm{cos}\frac{a_1\pi }{2}}}`$ $`\times {\displaystyle }dt{\displaystyle }{\displaystyle }{\displaystyle \frac{d\omega _1}{2\pi }}{\displaystyle \frac{dk_1}{2\pi }}{\displaystyle \frac{ie^{i\omega _1(t_1t)}}{\omega _1^2k_1^24+i\rho }}(f(k_1,x_1)f(k_1,0)e^{ik_1x_1}.`$ $`+.K_1(k_1)f(k_1,x_1)f(k_1,0)e^{ik_1x_1})`$ $`\times {\displaystyle }{\displaystyle }{\displaystyle \frac{d\omega _2}{2\pi }}{\displaystyle \frac{dk_2}{2\pi }}{\displaystyle \frac{ie^{i\omega _2(tt_2)}}{\omega _2^2k_2^24+i\rho }}(f(k_2,x_2)f(k_2,0)e^{ik_2x_2}.`$ $`.+K_2(k_2)f(k_2,x_2)f(k_2,0)e^{ik_2x_2}).`$ (4.14) First of all, it is necessary to perform the transformation $`k_1k_1`$ in the first term of the first propagator. Secondly, integration over $`t`$ ensures energy conservation at the interaction vertex and generates a Dirac delta function because of which we can set $`\omega _1=\omega _2`$. Moreover, it is better to define a new function as $$A(k,x)=f(k,x)f(k,0)+K(k)f(k,x)f(k,0)$$ (4.15) or, in an expanded form, $`A(k,x)`$ $`=`$ $`{\displaystyle \frac{ik+2\mathrm{coth}2x_0}{ik+2}}{\displaystyle \frac{ik+2\mathrm{coth}2(xx_0)}{ik2}}`$ (4.16) $`+{\displaystyle \frac{(ik+2\mathrm{cos}\frac{(a_0+a_1)\pi }{2})(ik+2\mathrm{cos}\frac{(a_0a_1)\pi }{2})}{(ik2\mathrm{cos}\frac{(a_0+a_1)\pi }{2})(ik2\mathrm{cos}\frac{(a_0a_1)\pi }{2})}}`$ $`\times {\displaystyle \frac{ik2\mathrm{coth}2x_0}{ik2}}{\displaystyle \frac{ik+2\mathrm{coth}2(xx_0)}{ik+2}},`$ then, the expression (4) reduces to $`{\displaystyle \frac{i\beta ^2}{32}}{\displaystyle \frac{(\sigma _1\mathrm{coth}x_0\sigma _0\mathrm{tanh}x_0)^2(a_0\mathrm{cot}a_0\pi +a_1\mathrm{cot}a_1\pi )}{\mathrm{cos}\frac{a_0\pi }{2}\mathrm{cos}\frac{a_1\pi }{2}}}`$ $`\times {\displaystyle }{\displaystyle }{\displaystyle }{\displaystyle \frac{d\omega _1}{2\pi }}{\displaystyle \frac{dk_1}{2\pi }}{\displaystyle \frac{dk_2}{2\pi }}e^{i\omega _1(t_1t_2)}{\displaystyle \frac{i}{\omega _1^2k_1^24+i\rho }}{\displaystyle \frac{i}{\omega _1^2k_2^24+i\rho }}e^{ik_1x_1}e^{ik_2x_2}`$ $`\times A(k_1,x_1)A(k_2,x_2).`$ (4.17) Obviously, what we need to do next is to integrate over the momenta $`k_1`$ and $`k_2`$ and this may be achieved by closing the contours in the upper half-plane and considering the poles at $`\widehat{k}_1=k_1=k_2=\sqrt{\omega _1^24}`$. Note, the additional poles due to functions $`A(k_1,x_1)`$ and $`A(k_2,x_2)`$ are not important because their contributions will be exponentially damped as $`x_1,x_2\mathrm{}`$. Therefore, the boundary-boundary contribution is $`{\displaystyle \frac{i\beta ^2}{32}}{\displaystyle \frac{(\sigma _1\mathrm{coth}x_0\sigma _0\mathrm{tanh}x_0)^2(a_0\mathrm{cot}a_0\pi +a_1\mathrm{cot}a_1\pi )}{\mathrm{cos}\frac{a_0\pi }{2}\mathrm{cos}\frac{a_1\pi }{2}}}`$ $`\times {\displaystyle }{\displaystyle \frac{d\omega _1}{2\pi }}e^{i\omega _1(t_1t_2)}e^{i\widehat{k}_1(x_1+x_2)}{\displaystyle \frac{1}{(2\widehat{k}_1)^2}}A(\widehat{k}_1,x_1)A(\widehat{k}_2,x_2).`$ (4.18) Now recall the definition of the quantum reflection factor as the coefficient of $`e^{ik(x+x^{})}`$ in the two-point Green function in the residue of the on-shell pole when $`x,x^{}\mathrm{}`$. Thus, the correction to the reflection factor from the type III (boundary- boundary) piece is $`{\displaystyle \frac{i\beta ^2}{32}}{\displaystyle \frac{(\sigma _1\mathrm{coth}x_0\sigma _0\mathrm{tanh}x_0)^2(a_0\mathrm{cot}a_0\pi +a_1\mathrm{cot}a_1\pi )}{\mathrm{cos}\frac{a_0\pi }{2}\mathrm{cos}\frac{a_1\pi }{2}}}`$ $`\times {\displaystyle \frac{1}{2\widehat{k}_1}}({\displaystyle \frac{(i\widehat{k}_1+2\mathrm{coth}2x_0)^2}{(i\widehat{k}_1+2)^2}}+2K(\widehat{k}_1){\displaystyle \frac{(i\widehat{k}_1+2\mathrm{coth}2x_0)(i\widehat{k}_12\mathrm{coth}2x_0)}{(i\widehat{k}_1+2)^2}}`$ $`+K^2(\widehat{k}_1){\displaystyle \frac{(i\widehat{k}_12\mathrm{coth}2x_0)^2}{(i\widehat{k}_1+2)^2}}).`$ (4.19) ## 5 Type III (boundary-bulk) This section deals with the determination of the contribution of the type III Feynman diagram to the classical reflection factor when one of the vertices corresponding to the loop is situated at the boundary and the other one is inside the bulk region. It is evident that in this case we have to take into account the bulk three point coupling $`C_{bulk}^{(3)}`$ in the corresponding vertex as well as the boundary three point coupling $`C_{boundary}^{(3)}`$ in the other vertex. Meanwhile, the combinatorial factor associated with the related Feynman diagram must be considered as a coefficient factor. Therefore, the contribution of the type III (boundary-bulk) to the reflection factor may be written as $`2\beta ^2(\sigma _1\mathrm{coth}x_0\sigma _0\mathrm{tanh}x_0){\displaystyle 𝑑t𝑑t^{}𝑑xG(x_1,t_1;x,t)G(x,t;0,t^{})}`$ $`\times G(0,t^{};0,t^{})G(x,t;x_2,t_2)\mathrm{sinh}(\sqrt{2}\beta \varphi _0).`$ (5.1) The propagator $`G(0,t^{};0,t^{})`$ corresponding to the loop has been found in the previous section and is given by $$G(0,t^{};0,t^{})=\frac{a_0}{2}\mathrm{cot}a_0\pi \frac{a_1}{2}\mathrm{cot}a_1\pi .$$ (5.2) The calculation of the other middle propagator i.e. $`G(x,t;0,t^{})`$ is the next step and clearly, the $`t^{}`$ dependence in (5) is included only in this propagator. Hence, it is convenient to compute the following relation $`{\displaystyle 𝑑t^{}G(x,t;0,t^{})}`$ $`=`$ $`{\displaystyle }dt^{}{\displaystyle }{\displaystyle }{\displaystyle \frac{d\omega }{2\pi }}{\displaystyle \frac{dk}{2\pi }}e^{i\omega (tt^{})}{\displaystyle \frac{i}{\omega ^2k^24+i\rho }}(f(k,x)f(k,0)e^{ikx}.`$ (5.3) $`.+K(k)f(k,x)f(k,0)e^{ikx}).`$ Integrating over $`t^{}`$ generates a Dirac delta function and so, $`{\displaystyle 𝑑t^{}G(x,t;0,t^{})}`$ $`=`$ $`i{\displaystyle }{\displaystyle \frac{dk}{2\pi }}{\displaystyle \frac{1}{(k^24+i\rho )}}(f(k,x)f(k,0)e^{ikx}.`$ (5.4) $`.+K(k)f(k,x)f(k,0)e^{ikx})`$ and the residue theorem gives $$𝑑t^{}G(x,t;0,t^{})=\frac{ie^{2x}}{8}\left(c_0+c_1\mathrm{coth}2(xx_0)\right),$$ (5.5) where $$c_0=c_1=\left(1+\mathrm{tan}^2\frac{(a_0+a_1)\pi }{4}\right)\left(1+\mathrm{tan}^2\frac{(a_0a_1)\pi }{4}\right).$$ (5.6) In order to check the above result, if we set $`x=0`$ in (5.5) then it will be equal to (4.13). Up to now the type III (boundary-bulk) contribution take the following form, of course, after integrating over $`t`$: $`\beta ^2c_0(\sigma _1\mathrm{coth}x_0\sigma _0\mathrm{tanh}x_0)(a_0\mathrm{cot}a_0\pi +a_1\mathrm{cot}a_1\pi )`$ $`\times {\displaystyle _{\mathrm{}}^0}dx{\displaystyle }{\displaystyle }{\displaystyle \frac{d\omega _1}{2\pi }}{\displaystyle \frac{dk_1}{2\pi }}e^{i\omega _1(t_1t_2)}{\displaystyle \frac{i}{\omega _1^2k_1^24+i\rho }}(f(k_1,x_1)f(k_1,x)e^{ik_1(x_1x)}.`$ $`.+K_1(k_1)f(k_1,x_1)f(k_1,x)e^{ik_1(x+x_1)})`$ $`\times \left\{{\displaystyle \frac{ie^{2x}}{8}}\left(1\mathrm{coth}2(xx_0)\right)\mathrm{sinh}(\sqrt{2}\beta \varphi _0)\right\}`$ $`\times {\displaystyle }{\displaystyle \frac{dk_2}{2\pi }}{\displaystyle \frac{i}{\omega _1^2k_2^24+i\rho }}(f(k_2,x)f(k_2,x_2)e^{ik_2(xx_2)}.`$ $`.+K_2(k_2)f(k_2,x)f(k_2,x_2)e^{ik_2(x+x_2)}).`$ (5.7) By multiplying the two propagators in (5) by each other, it is clear that one obtains four pole pieces and, as far as the integration over $`x`$ is concerned, if we do the integration over $`x`$ on one of them then, obviously the other three pole pieces could be done in the same manner. Hence, in what follows it is sufficient to treat only one of them and, meanwhile, keeping those terms which are functions of $`x`$, we are led to the following complicated integral $$_{\mathrm{}}^0𝑑x\mathrm{exp}\left\{2+i(k_2k_1)x\right\}f(k_1,x)f(k_2,x)\mathrm{sinh}(\sqrt{2}\beta \varphi _0)\left(1\mathrm{coth}2(xx_0)\right).$$ (5.8) After some substitutions and collecting together powers of $`\mathrm{coth}2(xx_0)`$ we obtain $`{\displaystyle \frac{1}{(ik_12)(ik_2+2)}}{\displaystyle _{\mathrm{}}^0}dx\mathrm{exp}\{2+i(k_2k_1)x\}\mathrm{sinh}(\sqrt{2}\beta \varphi _0)(k_1k_2.`$ $`.+(2ik_22ik_1+k_1k_2)\mathrm{coth}2(xx_0)+(2ik_12ik_24)\mathrm{coth}^22(xx_0).`$ $`.+4\mathrm{coth}^32(xx_0)).`$ (5.9) It is clear that in order to solve the above integral, it is necessary to manipulate the following integrals $$_{\mathrm{}}^0𝑑x\mathrm{exp}\left\{2+i(k_2k_1)x\right\}\mathrm{sinh}(\sqrt{2}\beta \varphi _0)\mathrm{coth}^n2(xx_0),$$ (5.10) where, $`n=0,1,2,3`$. In fact in Appendix A, we have found the integrals (5.10) and the solutions of them are expressed in terms of hypergeometric functions. So, using the formulae in Appendix A and simplifying, we find that $``$ in (5.9) can be rewritten $`(k_1,k_2)`$ $`=`$ $`{\displaystyle \frac{3k_1k_2+4ik_24ik_1+8}{3}}{\displaystyle \frac{1}{\mathrm{sinh}2x_0}}{\displaystyle \frac{3k_1k_2+6ik_26ik_1+13}{6}}{\displaystyle \frac{\mathrm{cosh}2x_0}{\mathrm{sinh}^22x_0}}`$ (5.11) $`{\displaystyle \frac{ik_2ik_1+2}{3}}{\displaystyle \frac{\mathrm{cosh}^22x_0+1}{\mathrm{sinh}^32x_0}}{\displaystyle \frac{1}{6}}{\displaystyle \frac{\mathrm{cosh}^32x_0+5\mathrm{cosh}2x_0}{\mathrm{sinh}^42x_0}}`$ $`{\displaystyle \frac{12ik_1k_216k_2+16k_1+40i(k_2k_1)(9k_1k_2+16ik_216ik_1+34)}{3(k_2k_14i)}}`$ $`\times e^{2x_0}F(1,{\displaystyle \frac{i}{4}}(k_2k_1)+1,{\displaystyle \frac{i}{4}}(k_2k_1)+2,e^{4x_0})`$ $`{\displaystyle \frac{12ik_1k_248k_2+48k_1+136i(k_2k_1)(6k_1k_2+28ik_228ik_1+84)}{3(k_2k_18i)}}`$ $`\times e^{6x_0}F(2,{\displaystyle \frac{i}{4}}(k_2k_1)+2,{\displaystyle \frac{i}{4}}(k_2k_1)+3,e^{4x_0})`$ $`+{\displaystyle \frac{32k_232k_1192i+(k_2k_1)(16ik_216ik_1+104)}{3(k_2k_112i)}}`$ $`\times e^{10x_0}F(3,{\displaystyle \frac{i}{4}}(k_2k_1)+3,{\displaystyle \frac{i}{4}}(k_2k_1)+4,e^{4x_0})`$ $`+{\displaystyle \frac{16k_216k_132i}{(k_2k_116i)}}`$ $`\times e^{14x_0}F(4,{\displaystyle \frac{i}{4}}(k_2k_1)+4,{\displaystyle \frac{i}{4}}(k_2k_1)+5,e^{4x_0}).`$ Now regarding (5), after doing the transformation $`k_1k_1`$ in the first term of the first propagator, all that remains is to integrate over the momenta $`k_1`$ and $`k_2`$ and this can be achieved by closing the contours in the upper half-plane and considering poles at $`\widehat{k}_1=k_1=k_2=\sqrt{\omega _1^24}`$. The extra poles in the four functions $`(\pm k_1,\pm k_2)`$ are not important because their contributions will be discounted when $`x_1`$ and $`x_2`$ go to $`\mathrm{}`$. Let us write down the type III (boundary-bulk) contribution to the reflection factor $`{\displaystyle \frac{\beta ^2}{2}}{\displaystyle \frac{\mathrm{tan}\frac{(a_0+a_1)\pi }{4}\mathrm{tan}\frac{(a_0a_1)\pi }{4}}{\mathrm{cos}\frac{a_0\pi }{2}\mathrm{cos}\frac{a_1\pi }{2}}}(a_0\mathrm{cot}a_0\pi +a_1\mathrm{cot}a_1\pi )`$ $`\times {\displaystyle }{\displaystyle \frac{d\omega _1}{2\pi }}e^{i\omega _1(t_1t_2)}e^{i\widehat{k}_1(x_1+x_2)}{\displaystyle \frac{1}{(2\widehat{k}_1)^2}}{\displaystyle \frac{i\widehat{k}_1+2\mathrm{coth}2(x_1x_0)}{i\widehat{k}_12}}{\displaystyle \frac{i\widehat{k}_1+2\mathrm{coth}2(x_2x_0)}{i\widehat{k}_12}}`$ $`\times \{{\displaystyle \frac{i}{(i\widehat{k}_1+2)^2}}(\widehat{k}_1,\widehat{k}_1){\displaystyle \frac{i}{(i\widehat{k}_1+2)(i\widehat{k}_12)}}K_1(\widehat{k}_1)(\widehat{k}_1,\widehat{k}_1)`$ $`{\displaystyle \frac{i}{(i\widehat{k}_1+2)(i\widehat{k}_12)}}K_1(\widehat{k}_1)(\widehat{k}_1,\widehat{k}_1)+{\displaystyle \frac{i}{(i\widehat{k}_12)^2}}K_1^2(\widehat{k}_1)(\widehat{k}_1,\widehat{k}_1)\}`$ (5.12) Now looking at the function $`(k_1,k_2)`$ given by (5.11), let us show the detailed forms of $`(\widehat{k}_1,\widehat{k}_1)`$, $`(\widehat{k}_1,\widehat{k}_1)`$,$`(\widehat{k}_1,\widehat{k}_1)`$ and $`(\widehat{k}_1,\widehat{k}_1)`$. In fact, $`(\widehat{k}_1,\widehat{k}_1)`$ $`=`$ $`{\displaystyle \frac{3\widehat{k}_1^2+8}{3}}{\displaystyle \frac{1}{\mathrm{sinh}2x_0}}+{\displaystyle \frac{3\widehat{k}_1^2+13}{6}}{\displaystyle \frac{\mathrm{cosh}2x_0}{\mathrm{sinh}^22x_0}}`$ (5.13) $`{\displaystyle \frac{2}{3}}{\displaystyle \frac{\mathrm{cosh}^22x_0+1}{\mathrm{sinh}^32x_0}}{\displaystyle \frac{1}{6}}{\displaystyle \frac{\mathrm{cosh}^32x_0+5\mathrm{cosh}2x_0}{\mathrm{sinh}^42x_0}}`$ $`+{\displaystyle \frac{3\widehat{k}_1^2+10}{3}}e^{2x_0}F(1,1,2,e^{4x_0})`$ $`+{\displaystyle \frac{3\widehat{k}_1^2+34}{6}}e^{6x_0}F(2,2,3,e^{4x_0})`$ $`+{\displaystyle \frac{16}{3}}e^{10x_0}F(3,3,4,e^{4x_0})`$ $`+2e^{14x_0}F(4,4,5,e^{4x_0}).`$ It can be easily verified that $$(\widehat{k}_1,\widehat{k}_1)=(\widehat{k}_1,\widehat{k}_1),$$ (5.14) and $`(\widehat{k}_1,\widehat{k}_1)`$ $`=`$ $`{\displaystyle \frac{3\widehat{k}_1^28i\widehat{k}_18}{3}}{\displaystyle \frac{1}{\mathrm{sinh}2x_0}}+{\displaystyle \frac{3\widehat{k}_1^212i\widehat{k}_113}{6}}{\displaystyle \frac{\mathrm{cosh}2x_0}{\mathrm{sinh}^22x_0}}`$ (5.15) $`{\displaystyle \frac{2i\widehat{k}_1+2}{3}}{\displaystyle \frac{\mathrm{cosh}^22x_0}{\mathrm{sinh}^32x_0}}{\displaystyle \frac{1}{6}}{\displaystyle \frac{\mathrm{cosh}^32x_0+5\mathrm{cosh}2x_0}{\mathrm{sinh}^42x_0}}`$ $`{\displaystyle \frac{9\widehat{k}_1^338i\widehat{k}_1^250\widehat{k}_1+20i}{3(\widehat{k}_12i)}}e^{2x_0}F(1,{\displaystyle \frac{i}{2}}\widehat{k}_1+1,{\displaystyle \frac{i}{2}}\widehat{k}_1+2,e^{4x_0})`$ $`{\displaystyle \frac{6\widehat{k}_1^362i\widehat{k}_1^2132\widehat{k}_1+68i}{3(\widehat{k}_14i)}}e^{6x_0}F(2,{\displaystyle \frac{i}{2}}\widehat{k}_1+2,{\displaystyle \frac{i}{2}}\widehat{k}_1+3,e^{4x_0})`$ $`+{\displaystyle \frac{32i\widehat{k}_1^2+136\widehat{k}_196i}{3(\widehat{k}_16i)}}e^{10x_0}F(3,{\displaystyle \frac{i}{2}}\widehat{k}_1+3,{\displaystyle \frac{i}{2}}\widehat{k}_1+4,e^{4x_0})`$ $`+{\displaystyle \frac{16\widehat{k}_116i}{(\widehat{k}_18i)}}e^{14x_0}F(4,{\displaystyle \frac{i}{2}}\widehat{k}_1+4,{\displaystyle \frac{i}{2}}\widehat{k}_1+5,e^{4x_0}).`$ Finally $`(\widehat{k}_1,\widehat{k}_1)`$ can be obtained from $`(\widehat{k}_1,\widehat{k}_1)`$ after setting $`\widehat{k}_1\widehat{k}_1`$. ## 6 Type III(bulk-boundary) In this section we study the quantum correction to the classical reflection factor due to the contribution of the type III Feynman diagram, when the vertex associated with the loop is located at the bulk region and the other vertex coincides with the boundary. The associated contribution is given by $`𝒞=2\beta ^2(\sigma _1\mathrm{coth}x_0\sigma _0\mathrm{tanh}x_0){\displaystyle 𝑑t𝑑t^{}𝑑x^{}G(x_1,t_1;0,t)G(0,t;x^{},t^{})}`$ $`\times G(x^{},t^{};x^{},t^{})G(0,t;x_2,t_2)\mathrm{sinh}(\sqrt{2}\beta \varphi _0).`$ (6.1) The following relation which is some part of the contribution (6), can be derived independently from the remaining part $$𝒞_1=𝑑tG(x_1,t_1;0,t)G(0,t;x_2,t_2)$$ (6.2) or $`𝒞_1`$ $`=`$ $`{\displaystyle }dt{\displaystyle }{\displaystyle }{\displaystyle \frac{d\omega _1}{2\pi }}{\displaystyle \frac{dk_1}{2\pi }}e^{i\omega _1(t_1t)}{\displaystyle \frac{i}{\omega _1^2k_1^24+i\rho }}(f(k_1,x_1)f(k_1,0)e^{ik_1x_1}.`$ (6.3) $`.+K_1(k_1)f(k_1,x_1)f(k_1,0)e^{ik_1x_1})`$ $`\times {\displaystyle }{\displaystyle }{\displaystyle \frac{d\omega _2}{2\pi }}{\displaystyle \frac{dk_2}{2\pi }}e^{i\omega _2(tt_2)}{\displaystyle \frac{i}{\omega _2^2k_2^24+i\rho }}(f(k_2,0)f(k_2,x_2)e^{ik_2x_2}.`$ $`.+K_2(k_2)f(k_2,0)f(k_2,x_2)e^{ik_2x_2}).`$ First of all, it is necessary to set $`k_1k_1`$ in the first term of the first propagator. Secondly, integration over $`t`$ leads to the substitution of $`\omega _2=\omega _1`$. Finally integration over the momenta $`k_1`$ and $`k_2`$, as before, may be done immediately by closing the contour in the upper half-plane and looking at the poles at $`\widehat{k}_1=k_1=k_2=\sqrt{\omega _1^24}`$ and ignoring all the other poles as their contributions vanish rapidly as $`x_1,x_2\mathrm{}`$. Therefore, $`𝒞_1`$ $`=`$ $`{\displaystyle \frac{d\omega _1}{2\pi }e^{i\omega _1(t_1t_2)}e^{i\widehat{k}_1(x_1+x_2)}\frac{1}{(2\widehat{k}_1)^2}A(\widehat{k}_1,x_1)A(\widehat{k}_1,x_2)},`$ (6.4) where $$A(\widehat{k}_1,x_1)=f(\widehat{k}_1,x_1)f(\widehat{k}_1,0)+K(\widehat{k}_1)f(\widehat{k}_1,x_1)f(\widehat{k}_1,0).$$ (6.5) So, our next problem is to calculate the following integral which is the remaining part of the contribution. $$𝑑t^{}𝑑x^{}G(0,t;x^{},t^{})G(x^{},t^{};x^{},t^{})\mathrm{sinh}(\sqrt{2}\beta \varphi _0).$$ (6.6) Obviously, this part will be appeared as a constant and it must be multiplied by (6.4). Clearly, the time variable $`t^{}`$ appears only in one of the propagator i.e. in $`G(0,t;x^{},t^{})`$. On the other hand, this propagator along with integration over $`t^{}`$ has been obtained in the previous section. Hence, $$𝑑t^{}G(0,t;x^{},t^{})=\frac{i}{8}e^{2x^{}}c_0\left(1\mathrm{coth}2(x^{}x_0)\right),$$ (6.7) where $$c_0=\left(1+\mathrm{tan}^2\frac{(a_0+a_1)\pi }{4}\right)\left(1+\mathrm{tan}^2\frac{(a_0a_1)\pi }{4}\right).$$ (6.8) Therefore, (6.6) reduces to $$\frac{ic_0}{8}𝑑x^{}e^{2x^{}}\mathrm{sinh}(\sqrt{2}\beta \varphi _0)\left(1\mathrm{coth}2(x^{}x_0)\right)G(x^{},t^{};x^{},t^{}),$$ (6.9) where $`G(x^{},t^{};x^{},t^{})={\displaystyle }{\displaystyle }{\displaystyle \frac{d\omega ^{}}{2\pi }}{\displaystyle \frac{dk^{}}{2\pi }}{\displaystyle \frac{i}{\omega ^2k^24+i\rho }}(f(k^{},x^{})f(k^{},x^{}).`$ $`.+K(k^{})(f(k^{},x^{})(f(k^{},x^{})e^{2ik^{}x^{}})`$ (6.10) or after integration over $`\omega ^{}`$ $`G(x^{},t^{};x^{},t^{})={\displaystyle \frac{1}{2}}{\displaystyle }{\displaystyle \frac{dk^{}}{2\pi }}{\displaystyle \frac{1}{\sqrt{k^2+4}}}(f(k^{},x^{})f(k^{},x^{}).`$ $`.+K(k^{})(f(k^{},x^{})(f(k^{},x^{})e^{2ik^{}x^{}}).`$ (6.11) In fact, the above integrand has two parts, the first part can be easily manipulated but the other part which includes the exponential term is hard to calculate and we prefer to leave the computation of that part for later. Let us look at the first part of the loop propagator. The integral of this part is logarithmically divergent. Nevertheless, this divergence can be removed by an infinite renormalization of the mass parameter in the bulk potential. Then, doing the integration over $`k^{}`$, we obtain $$\frac{1}{2}\frac{dk^{}}{2\pi }\frac{1}{\sqrt{k^2+4}}f(k^{},x^{})f(k^{},x^{})=\frac{\left(1\mathrm{coth}^22(x^{}x_0)\right)}{2\pi }.$$ (6.12) To sum up, the integral (6.9) reduces to $`{\displaystyle \frac{ic_0}{16\pi }}{\displaystyle _{\mathrm{}}^0}𝑑x^{}e^{2x^{}}\left(1\mathrm{coth}^22(x^{}x_0)\right)\left(1\mathrm{coth}2(x^{}x_0)\right)\mathrm{sinh}(\sqrt{2}\beta \varphi _0)`$ $`+{\displaystyle \frac{ic_0}{16}}{\displaystyle _{\mathrm{}}^0}𝑑x^{}{\displaystyle \frac{dk^{}}{2\pi }\frac{1}{\sqrt{k^2+4}}\left(1\mathrm{coth}2(x^{}x_0)\right)\mathrm{sinh}(\sqrt{2}\beta \varphi _0)e^{2x^{}}}`$ $`\times \left\{{\displaystyle \frac{(ik^{})^2+2ik^{}\sqrt{1+\sigma _0}\sqrt{1+\sigma _1}+2(\sigma _0+\sigma _1)}{(ik^{})^22ik^{}\sqrt{1+\sigma _0}\sqrt{1+\sigma _1}+2(\sigma _0+\sigma _1)}}{\displaystyle \frac{\left(ik^{}+2\mathrm{coth}2(x^{}x_0)\right)^2}{(ik^{}+2)(ik^{}2)}}e^{2ik^{}x^{}}\right\}`$ (6.13) The above relation has two parts and the first part which is a single integral can be performed by means of the formulae in Appendix A and we write down only the solution of this part which is expressed in terms of hypergeometric functions, that is, $`𝒞_2`$ $`=`$ $`{\displaystyle \frac{ic_0}{16\pi }}{\displaystyle _{\mathrm{}}^0}𝑑x^{}e^{2x^{}}\left(1\mathrm{coth}^22(x^{}x_0)\right)\left(1\mathrm{coth}2(x^{}x_0)\right)\mathrm{sinh}(\sqrt{2}\beta \varphi _0)`$ (6.14) $`=`$ $`{\displaystyle \frac{i}{16\pi }}\left(1+\mathrm{tan}^2{\displaystyle \frac{(a_0+a_1)\pi }{4}}\right)\left(1+\mathrm{tan}^2{\displaystyle \frac{(a_0a_1)\pi }{4}}\right)`$ $`\times \{{\displaystyle \frac{1}{3\mathrm{sinh}2x_0}}{\displaystyle \frac{\mathrm{cosh}2x_0}{24\mathrm{sinh}^22x_0}}{\displaystyle \frac{\mathrm{cosh}^22x_0+1}{6\mathrm{sinh}^32x_0}}{\displaystyle \frac{\mathrm{cosh}^32x_0+5\mathrm{cosh}2x_0}{24\mathrm{sinh}^42x_0}}`$ $`{\displaystyle \frac{1}{6}}e^{2x_0}F(1,1,2,e^{4x_0})`$ $`+{\displaystyle \frac{11}{12}}e^{6x_0}F(2,2,3,e^{4x_0})`$ $`+{\displaystyle \frac{4}{3}}e^{10x_0}F(3,3,4,e^{4x_0})`$ $`+{\displaystyle \frac{1}{2}}e^{14x_0}F(4,4,5,e^{4x_0})\}.`$ So, in connection with the type III (bulk-boundary) contribution, the remaining integral is $`𝒞_3={\displaystyle \frac{ic_0}{16}}{\displaystyle _{\mathrm{}}^0}𝑑x^{}{\displaystyle \frac{dk^{}}{2\pi }\frac{1}{\sqrt{k^2+4}}\left(1\mathrm{coth}2(x^{}x_0)\right)\mathrm{sinh}(\sqrt{2}\beta \varphi _0)e^{2x^{}}}`$ $`\times {\displaystyle \frac{\left(ik^{}+2\mathrm{cos}\frac{(a_0+a_1)\pi }{2}\right)\left(ik^{}+2\mathrm{cos}\frac{(a_0a_1)\pi }{2}\right)}{\left(ik^{}2\mathrm{cos}\frac{(a_0+a_1)\pi }{2}\right)\left(ik^{}2\mathrm{cos}\frac{(a_0a_1)\pi }{2}\right)}}{\displaystyle \frac{\left(ik^{}+2\mathrm{coth}2(x^{}x_0)\right)^2}{(ik^{}+2)(ik^{}2)}}e^{2ik^{}x^{}}.`$ (6.15) As we mentioned before, it is more convenient to integrate over $`x^{}`$ then afterwards over $`k^{}`$. Since to integrate over $`k^{}`$ first is a difficult problem. Let us do partial fraction decomposition for the rational function in (6). Obviously we will have four elementary partial fraction including $$\frac{1}{\left(ik^{}2\mathrm{cos}\frac{(a_0+a_1)\pi }{2}\right)},\frac{1}{\left(ik^{}2\mathrm{cos}\frac{(a_0a_1)\pi }{2}\right)},\frac{1}{ik^{}+2},\frac{1}{ik^{}2}.$$ Now, in what follows we perform the calculations in detail for one of them, for example, $`\frac{1}{\left(ik^{}2\mathrm{cos}\frac{(a_0+a_1)\pi }{2}\right)}`$ due to the fact that for all the others the computations are similar except that $`\mathrm{cos}\frac{(a_0+a_1)\pi }{2}`$ is replaced by one of $`\mathrm{cos}\frac{(a_0a_1)\pi }{2}`$, -1, 1, respectively. What we need to do is to calculate the following : $`{\displaystyle \frac{\left(\mathrm{tan}^2\frac{(a_0+a_1)\pi }{4}\mathrm{cot}^2\frac{(a_0+a_1)\pi }{4}\right)}{\mathrm{cos}^2\frac{(a_0+a_1)\pi }{4}\mathrm{cos}^2\frac{(a_0a_1)\pi }{4}}}\mathrm{cot}{\displaystyle \frac{a_0\pi }{2}}\mathrm{cot}{\displaystyle \frac{a_1\pi }{2}}`$ $`\times {\displaystyle \frac{i}{16}}{\displaystyle _{\mathrm{}}^0}𝑑x^{}{\displaystyle \frac{dk^{}}{2\pi }\frac{1}{\sqrt{k^2+4}}\left(1\mathrm{coth}2(x^{}x_0)\right)\mathrm{sinh}(\sqrt{2}\beta \varphi _0)e^{(22ik^{})x^{}}}`$ $`\times \left(\mathrm{coth}2(x^{}x_0)+\mathrm{cos}{\displaystyle \frac{(a_0+a_1)\pi }{2}}\right)^2\left({\displaystyle \frac{1}{ik^{}2\mathrm{cos}\frac{(a_0+a_1)\pi }{2}}}\right).`$ (6.16) The integration over $`x^{}`$ may be done by using the formulae in Appendix A and gives $`{\displaystyle _{\mathrm{}}^0}𝑑x^{}e^{(22ik^{})x^{}}\left(1\mathrm{coth}2(x^{}x_0)\right)\mathrm{sinh}(\sqrt{2}\beta \varphi _0)`$ (6.17) $`\times \left(\mathrm{coth}2(x^{}x_0)+\mathrm{cos}{\displaystyle \frac{(a_0+a_1)\pi }{2}}\right)^2`$ $`=`$ $`L(a_0,a_1)`$ $`+`$ $`{\displaystyle \underset{n=1}{\overset{4}{}}}\left\{{\displaystyle \frac{(A_nk^{}+B_n)}{(k^{}+2ni)}}e^{(2+4(n1))x_0}F(n,{\displaystyle \frac{i}{2}}k^{}+n,{\displaystyle \frac{i}{2}}k^{}+n+1,e^{4x_0})\right\},`$ where $`L(a_0,a_1)`$ $`=`$ $`\left(\mathrm{cos}^2{\displaystyle \frac{(a_0+a_1)\pi }{2}}{\displaystyle \frac{4}{3}}\mathrm{cos}{\displaystyle \frac{(a_0+a_1)\pi }{2}}+{\displaystyle \frac{2}{3}}\right){\displaystyle \frac{1}{\mathrm{sinh}2x_0}}`$ (6.18) $`+`$ $`\left({\displaystyle \frac{1}{2}}\mathrm{cos}^2{\displaystyle \frac{(a_0+a_1)\pi }{2}}\mathrm{cos}{\displaystyle \frac{(a_0+a_1)\pi }{2}}+{\displaystyle \frac{13}{24}}\right){\displaystyle \frac{\mathrm{cosh}2x_0}{\mathrm{sinh}^22x_0}}`$ $``$ $`\left({\displaystyle \frac{1}{3}}\mathrm{cos}{\displaystyle \frac{(a_0+a_1)\pi }{2}}{\displaystyle \frac{1}{6}}\right){\displaystyle \frac{\mathrm{cosh}^22x_0+1}{\mathrm{sinh}^32x_0}}`$ $`+`$ $`{\displaystyle \frac{1}{24}}{\displaystyle \frac{\mathrm{cosh}^32x_0+5\mathrm{cosh}2x_0}{\mathrm{sinh}^42x_0}}`$ and the coefficients $`A_n`$, $`B_n`$, $`n=1,2,3,4`$ are constants which in fact only depend on $`\mathrm{cos}\frac{(a_0+a_1)\pi }{2}`$. Now the final calculation is to integrate over $`k^{}`$ and it is evident that in order to do this, we have to convert the hypergeometric functions to infinite series. Considering the equation (A.9) in Appendix A, we conclude that $$\frac{F(1,\frac{i}{2}k^{}+1,\frac{i}{2}k^{}+2,e^{4x_0})}{k^{}+2i}=\underset{n=0}{\overset{\mathrm{}}{}}\frac{e^{4nx_0}}{k^{}+i(2+2n)}.$$ (6.19) If we differentiate both sides of the above relation with respect to $`x_0`$, then the following identity may be derived $$\frac{F(2,\frac{i}{2}k^{}+2,\frac{i}{2}k^{}+3,e^{4x_0})}{k^{}+4i}=\underset{n=1}{\overset{\mathrm{}}{}}\frac{ne^{4(n1)x_0}}{k^{}+i(2+2n)}.$$ (6.20) In the same way, one obtains $$\frac{F(3,\frac{i}{2}k^{}+3,\frac{i}{2}k^{}+4,e^{4x_0})}{k^{}+6i}=\frac{1}{2!}\underset{n=2}{\overset{\mathrm{}}{}}\frac{n(n1)e^{4(n2)x_0}}{k^{}+i(2+2n)}$$ (6.21) and $$\frac{F(4,\frac{i}{2}k^{}+4,\frac{i}{2}k^{}+5,e^{4x_0})}{k^{}+8i}=\frac{1}{3!}\underset{n=3}{\overset{\mathrm{}}{}}\frac{n(n1)(n2)e^{4(n3)x_0}}{k^{}+i(2+2n)}.$$ (6.22) Now if we substitute (6.19), (6.20), (6.21) and (6.22) into (6.17), all that remains in connection with the contribution (6.16) is the integration over $`k^{}`$. Obviously we encounter integrals of the form $$_{\mathrm{}}^{\mathrm{}}\frac{dk^{}}{\sqrt{k^2+4}}\left(\frac{1}{ik^{}2\mathrm{cos}\frac{(a_0+a_1)\pi }{2}}\right)\left(\frac{Ak^{}+B}{k^{}+i(2+2n)}\right)$$ (6.23) and the $`k^{}`$ integration may be performed by closing the contour in the upper half-plane and onto the branch cut which stretches from $`k^{}=2i`$ to infinity along the imaginary axis. In fact, leaving the integrals along the branch cut to be evaluated later, we obtain the required formula $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{dk^{}}{\sqrt{k^2+4}}}\left({\displaystyle \frac{1}{ik^{}2\mathrm{cos}\frac{(a_0+a_1)\pi }{2}}}\right)\left({\displaystyle \frac{Ak^{}+B}{k^{}+i(2+2n)}}\right)`$ $`={\displaystyle \frac{\left(2\mathrm{cos}\frac{(a_0+a_1)\pi }{2}A+iB\right)}{\left(2n+22\mathrm{cos}\frac{(a_0+a_1)\pi }{2}\right)}}{\displaystyle \frac{\frac{(a_0+a_1)\pi }{2}}{\mathrm{sin}\frac{(a_0+a_1)\pi }{2}}}`$ $`{\displaystyle \frac{\left((2n+2)A+iB\right)}{\left(2n+22\mathrm{cos}\frac{(a_0+a_1)\pi }{2}\right)}}{\displaystyle \frac{1}{2\sqrt{n^2+2n}}}\mathrm{ln}\left\{{\displaystyle \frac{n+1\sqrt{n^2+2n}}{n+1+\sqrt{n^2+2n}}}\right\}.`$ (6.24) Note (6) is valid when $`n0`$, on the other hand if $`n=0`$ then one may find $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{dk^{}}{\sqrt{k^2+4}}}\left({\displaystyle \frac{1}{ik^{}2\mathrm{cos}\frac{(a_0+a_1)\pi }{2}}}\right)\left({\displaystyle \frac{Ak^{}+B}{k^{}+2i}}\right)`$ $`={\displaystyle \frac{\left(2\mathrm{cos}\frac{(a_0+a_1)\pi }{2}A+iB\right)}{\left(22\mathrm{cos}\frac{(a_0+a_1)\pi }{2}\right)}}{\displaystyle \frac{\frac{(a_0+a_1)\pi }{2}}{\mathrm{sin}\frac{(a_0+a_1)\pi }{2}}}{\displaystyle \frac{\left(2A+iB\right)}{\left(22\mathrm{cos}\frac{(a_0+a_1)\pi }{2}\right)}}.`$ (6.25) Now we are in a position to write down (6) or, in fact, the integral (6.15) $`𝒞_3`$ $`=`$ $`{\displaystyle \frac{i}{32\pi }}{\displaystyle \frac{\left(\mathrm{tan}^2\frac{(a_0+a_1)\pi }{4}\mathrm{cot}^2\frac{(a_0+a_1)\pi }{4}\right)}{\mathrm{cos}^2\frac{(a_0+a_1)\pi }{4}\mathrm{cos}^2\frac{(a_0a_1)\pi }{4}}}\mathrm{cot}{\displaystyle \frac{a_0\pi }{2}}\mathrm{cot}{\displaystyle \frac{a_1\pi }{2}}`$ (6.26) $`\times \{{\displaystyle \frac{\frac{(a_0+a_1)\pi }{2}}{\mathrm{sin}\frac{(a_0+a_1)\pi }{2}}}L(a_0,a_1)`$ $`+{\displaystyle \frac{1}{12}}e^{2x_0}\left({\displaystyle \frac{2A_1+iB_1}{22\mathrm{cos}\frac{(a_0+a_1)\pi }{2}}}{\displaystyle \frac{2\mathrm{cos}\frac{(a_0+a_1)\pi }{2}A_1+iB_1}{22\mathrm{cos}\frac{(a_0+a_1)\pi }{2}}}{\displaystyle \frac{\frac{(a_0+a_1)\pi }{2}}{\mathrm{sin}\frac{(a_0+a_1)\pi }{2}}}\right)`$ $`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\frac{(a_0+a_1)\pi }{2}}{\mathrm{sin}\frac{(a_0+a_1)\pi }{2}}}{\displaystyle \frac{e^{(2+4n)x_0}}{\left(2n+22\mathrm{cos}\frac{(a_0+a_1)\pi }{2}\right)}}[(2\mathrm{cos}{\displaystyle \frac{(a_0+a_1)\pi }{2}}A_1+iB_1).`$ $`.+{\displaystyle \frac{n}{1!}}(2\mathrm{cos}{\displaystyle \frac{(a_0+a_1)\pi }{2}}A_2+iB_2)+{\displaystyle \frac{n(n1)}{2!}}(2\mathrm{cos}{\displaystyle \frac{(a_0+a_1)\pi }{2}}A_3+iB_3).`$ $`.+{\displaystyle \frac{n(n1)(n2)}{3!}}(2\mathrm{cos}{\displaystyle \frac{(a_0+a_1)\pi }{2}}A_4+iB_4)]`$ $`+{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{e^{(2+4n)x_0}}{\left(2n+22\mathrm{cos}\frac{(a_0+a_1)\pi }{2}\right)}}{\displaystyle \frac{1}{2\sqrt{n^2+2n}}}\mathrm{ln}\left\{{\displaystyle \frac{n+1\sqrt{n^2+2n}}{n+1+\sqrt{n^2+2n}}}\right\}`$ $`([(2n+2)A_1+iB_1]+{\displaystyle \frac{n}{1!}}[(2n+2)A_2+iB_2].`$ $`.+{\displaystyle \frac{n(n1)}{2!}}[(2n+2)A_3+iB_3]+{\displaystyle \frac{n(n1)(n2)}{3!}}[(2n+2)A_4+iB_4])\}`$ $`+\text{other pole pieces}.`$ Note, in the above expression all the series are convergent. As we mentioned before, (6.26) must be considered (after adding to (6.14)) as a coefficient factor of (6.4) in order to constitute the type III (bulk-boundary) contribution i.e.: $$𝒞=𝒞_1\left(𝒞_2+𝒞_3\right).$$ (6.27) ## 7 Type I diagram In this section we calculate the contribution of the type I Feynman diagram to the classical reflection factor when the vertex is placed inside the bulk region. In fact, when the vertex is located at the boundary then, the corresponding contribution has been found and is given by $`{\displaystyle \frac{i\beta ^2}{8}}(\sigma _1\mathrm{coth}x_0+\sigma _0\mathrm{tanh}x_0)(a_0\mathrm{cot}a_0\pi +a_1\mathrm{cot}a_1\pi )`$ $`\times {\displaystyle }{\displaystyle \frac{d\omega }{2\pi }}e^{i\omega (t_1t_2)}e^{i\widehat{k}(x_1+x_2)}{\displaystyle \frac{i\widehat{k}+2\mathrm{coth}2(x_1x_0)}{P(\widehat{k})}}{\displaystyle \frac{i\widehat{k}+2\mathrm{coth}2(x_2x_0)}{P(\widehat{k})}}`$ (7.1) Clearly, in our case the bulk four point coupling should be considered in the interaction vertex. Moreover, as before, the combinatorial factor associated with this diagram will appear as a coefficient. Hence, the contribution has the form $$4i\beta ^2_{\mathrm{}}^{\mathrm{}}𝑑t_{\mathrm{}}^0𝑑xG(x_1,t_1;x,t)G(x,t;x,t)G(x,t;x_2,t_2)\mathrm{cosh}(\sqrt{2}\beta \varphi _0),$$ (7.2) where $$\mathrm{cosh}(\sqrt{2}\beta \varphi _0)=\left(2\mathrm{coth}^22(xx_0)1\right).$$ (7.3) In the previous section, we simplified the loop propagator to $`G(x,t;x,t)`$ $`=`$ $`{\displaystyle \frac{\left(1\mathrm{coth}^22(xx_0)\right)}{2\pi }}`$ (7.4) $`+{\displaystyle \frac{1}{2}}{\displaystyle \frac{dk^{\prime \prime }}{2\pi }\frac{1}{\sqrt{k^{\prime \prime 2}+4}}K(k^{\prime \prime })f(k^{\prime \prime },x)f(k^{\prime \prime },x)e^{2ik^{\prime \prime }x}}.`$ Also the integral part of the loop Green function is hard enough to evaluate and we found out that it is better to do this integration during the final stage. Now, let us rewrite the contribution (7.2) in the expanded form $`4i\beta ^2{\displaystyle }dt{\displaystyle _{\mathrm{}}^0}dx{\displaystyle }{\displaystyle }{\displaystyle \frac{d\omega }{2\pi }}{\displaystyle \frac{dk}{2\pi }}e^{i\omega (t_1t)}{\displaystyle \frac{i}{\omega ^2k^24+i\rho }}(f(k,x_1)f(k,x)e^{ik(x_1x)}.`$ $`.+K(k)f(k,x_1)f(k,x)e^{ik(x_1+x)})\mathrm{cosh}(\sqrt{2}\beta \varphi _0)`$ $`\times \left\{{\displaystyle \frac{\left(1\mathrm{coth}^22(xx_0)\right)}{2\pi }}+{\displaystyle \frac{1}{2}}{\displaystyle \frac{dk^{\prime \prime }}{2\pi }\frac{1}{\sqrt{k^{\prime \prime 2}+4}}K(k^{\prime \prime })f(k^{\prime \prime },x)f(k^{\prime \prime },x)e^{2ik^{\prime \prime }x}}\right\}`$ $`\times {\displaystyle }{\displaystyle }{\displaystyle \frac{d\omega ^{}}{2\pi }}{\displaystyle \frac{dk^{}}{2\pi }}e^{i\omega ^{}(tt_2)}{\displaystyle \frac{i}{\omega ^2k^24+i\rho }}(f(k^{},x)f(k^{},x_2)e^{ik^{}(xx_2)}`$ $`+K^{}(k^{})f(k^{},x)f(k^{},x_2)e^{ik^{}(x+x_2)}).`$ (7.5) Looking at (7), one can predict that the calculations will be lengthy and intricate. The starting point is to do the $`t`$ integration which allows the substitution $`\omega =\omega ^{}`$. Secondly, it is necessary to perform a transformation $`kk`$ in the first term of the first propagator. Moreover, if we multiply the first and the third propagator with each other, then obviously we will have four pole pieces and fortunately if we do the calculation for one of them (for example the first one), then the calculations corresponding to the other three pole pieces may be treated similarly with $`k+k^{}`$ replaced by one of $`kk^{}`$, $`k+k^{}`$ and $`kk^{}`$. Because of this in what follows we follow the problem only for one pole piece. Hence our problem is, in fact, the following integral $`𝒟=4i\beta ^2{\displaystyle _{\mathrm{}}^0}𝑑x{\displaystyle \frac{d\omega }{2\pi }\frac{dk}{2\pi }e^{i\omega (t_1t_2)}\frac{i}{\omega ^2k^24+i\rho }}`$ $`\times \mathrm{cosh}(\sqrt{2}\beta \varphi _0)f(k,x_1)f(k,x)e^{ik(x_1x)}`$ $`\times \left({\displaystyle \frac{\left(1\mathrm{coth}^22(xx_0)\right)}{2\pi }}+{\displaystyle \frac{1}{2}}{\displaystyle \frac{dk^{\prime \prime }}{2\pi }\frac{1}{\sqrt{k^{\prime \prime 2}+4}}K(k^{\prime \prime })f(k^{\prime \prime },x)f(k^{\prime \prime },x)e^{2ik^{\prime \prime }x}}\right)`$ $`\times {\displaystyle }{\displaystyle \frac{dk^{}}{2\pi }}{\displaystyle \frac{i}{\omega ^2k^24+i\rho }}f(k^{},x)f(k^{},x_2)e^{ik^{}(xx_2)}.`$ (7.6) In fact, the above contribution has two parts. The first part, in which the integral of the middle momentum $`(k^{\prime \prime })`$ is not involved, can be calculated by means of the formulae in Appendix B and we call this part $`𝒟_1`$. Let us write down the solution of this part. This contribution is expressed in terms of the hypergeometric function as : $`𝒟_1`$ $`=`$ $`{\displaystyle \frac{i\beta ^2}{\pi }}{\displaystyle \frac{d\omega }{2\pi }e^{i\omega (t_1t_2)}e^{i\widehat{k}(x_1+x_2)}f(\widehat{k},x_1)f(\widehat{k},x_2)\frac{1}{\widehat{k}^2}}`$ (7.7) $`\times \{{\displaystyle \frac{4i4\widehat{k}i\widehat{k}^2}{\widehat{k}2i}}e^{4x_0}F(2,{\displaystyle \frac{i}{2}}\widehat{k}+1,{\displaystyle \frac{i}{2}}\widehat{k}+2,e^{4x_0})`$ $`+{\displaystyle \frac{48i40\widehat{k}8i\widehat{k}^2}{\widehat{k}4i}}e^{8x_0}F(3,{\displaystyle \frac{i}{2}}\widehat{k}+2,{\displaystyle \frac{i}{2}}\widehat{k}+3,e^{4x_0})`$ $`+{\displaystyle \frac{176i96\widehat{k}8i\widehat{k}^2}{\widehat{k}6i}}e^{12x_0}F(4,{\displaystyle \frac{i}{2}}\widehat{k}+3,{\displaystyle \frac{i}{2}}\widehat{k}+4,e^{4x_0})`$ $`+{\displaystyle \frac{256i64\widehat{k}}{\widehat{k}8i}}e^{16x_0}F(5,{\displaystyle \frac{i}{2}}\widehat{k}+4,{\displaystyle \frac{i}{2}}\widehat{k}+5,e^{4x_0})`$ $`+{\displaystyle \frac{128i}{\widehat{k}10i}}e^{20x_0}F(6,{\displaystyle \frac{i}{2}}\widehat{k}+5,{\displaystyle \frac{i}{2}}\widehat{k}+6,e^{4x_0})\}.`$ Now it is better for the second part, which we call $`𝒟_2`$, to integrate first over $`x`$ then over $`k^{\prime \prime }`$. Meanwhile, before starting the integration, it is useful to note that if we do the partial fraction decomposition for $`K^{\prime \prime }(k^{\prime \prime })f(k^{\prime \prime },x)f(k^{\prime \prime },x)`$, then we will have four elementary partial fractions as $$\frac{1}{\left(ik^{\prime \prime }2\mathrm{cos}\frac{(a_0+a_1)\pi }{2}\right)},\frac{1}{\left(ik^{\prime \prime }2\mathrm{cos}\frac{(a_0a_1)\pi }{2}\right)},\frac{1}{ik^{\prime \prime }+2},\frac{1}{ik^{\prime \prime }2}.$$ As before, in the remaining section we continue the computations in detail for one of them (for example, $`\frac{1}{\left(ik^{\prime \prime }2\mathrm{cos}\frac{(a_0+a_1)\pi }{2}\right)}`$) because the calculations corresponding to the other three elementary partial fractions can be done in the same manner just by the substitution of $`\mathrm{cos}\frac{(a_0+a_1)\pi }{2}`$ by one of $`\mathrm{cos}\frac{(a_0a_1)\pi }{2}`$, -1, 1, respectively. So, our problem reduces to this integral $`2i\beta ^2\mathrm{cot}{\displaystyle \frac{a_0\pi }{2}}\mathrm{cot}{\displaystyle \frac{a_1\pi }{2}}\left(\mathrm{tan}^2{\displaystyle \frac{(a_0+a_1)\pi }{4}}\mathrm{cot}^2{\displaystyle \frac{(a_0+a_1)\pi }{4}}\right)`$ $`\times {\displaystyle _{\mathrm{}}^0}dx{\displaystyle }{\displaystyle }{\displaystyle \frac{d\omega }{2\pi }}{\displaystyle \frac{dk}{2\pi }}e^{i\omega (t_1t_2)}{\displaystyle \frac{i}{\omega ^2k^24+i\rho }}`$ $`\times \mathrm{cosh}(\sqrt{2}\beta \varphi _0)f(k,x_1)f(k,x)e^{ik(x_1x)}`$ $`\times {\displaystyle }{\displaystyle \frac{dk^{\prime \prime }}{2\pi }}{\displaystyle \frac{1}{\sqrt{k^{\prime \prime 2}+4}}}{\displaystyle \frac{e^{2ik^{\prime \prime }x}}{ik^{\prime \prime }2\mathrm{cos}\frac{(a_0+a_1)\pi }{2}}}(\mathrm{coth}2(xx_0)+\mathrm{cos}{\displaystyle \frac{(a_0+a_1)\pi }{2}})^2`$ $`\times {\displaystyle }{\displaystyle \frac{dk^{}}{2\pi }}{\displaystyle \frac{i}{\omega ^2k^24+i\rho }}f(k^{},x)f(k^{},x_2)e^{ik^{}(xx_2)}.`$ (7.8) The integration over $`x`$ can be evaluated by means of the formulae given in Appendix B as $`{\displaystyle _{\mathrm{}}^0}𝑑xe^{i(k+k^{}2k^{\prime \prime })x}\mathrm{cosh}(\sqrt{2}\beta \varphi _0)f(k,x)f(k^{},x)\left(\mathrm{coth}2(xx_0)+\mathrm{cos}{\displaystyle \frac{(a_0+a_1)\pi }{2}}\right)^2`$ $`={\displaystyle \frac{1}{(ik+2)(ik^{}+2)}}{\displaystyle \frac{\left(A_0^{}kk^{}+B_0^{}(k+k^{})+C_0^{}\right)}{(k+k^{}2k^{\prime \prime })}}`$ $`+{\displaystyle \frac{1}{(ik+2)(ik^{}+2)}}{\displaystyle \underset{n=1}{\overset{6}{}}}\{{\displaystyle \frac{\left(A_n^{}kk^{}+B_n^{}(k+k^{})+C_n^{}\right)}{\left(k+k^{}2k^{\prime \prime }4(n1)i\right)}}e^{4(n1)x_0}`$ $`\times F(n,{\displaystyle \frac{i}{4}}(k+k^{}2k^{\prime \prime })+n1,{\displaystyle \frac{i}{4}}(k+k^{}2k^{\prime \prime })+n,e^{4x_0})\},`$ (7.9) where the coefficients $`A_n^{},B_n^{},C_n^{};n=0,1,..,6`$ are constants and depend only on $`\mathrm{cos}\frac{(a_0+a_1)\pi }{2}`$. In fact $`A_5^{}`$, $`A_6^{}`$ and $`B_6^{}`$ are zero. Now the subsequent calculation is to integrate over $`k^{\prime \prime }`$ and it is clear that to do this, it is necessary to convert the hypergeometric function to an infinite series. Looking at (B.9), we may write down $$\frac{F(1,\frac{i}{4}(k+k^{}2k^{\prime \prime }),\frac{i}{4}(k+k^{}2k^{\prime \prime })+1,e^{4x_0})}{k+k^{}2k^{\prime \prime }}=\underset{n=0}{\overset{\mathrm{}}{}}\frac{e^{4nx_0}}{k+k^{}2k^{\prime \prime }4ni}$$ (7.10) and by differentiating both sides of the above relation with respect to $`x_0`$, then we obtain $$\frac{F(2,\frac{i}{4}(k+k^{}2k^{\prime \prime })+1,\frac{i}{4}(k+k^{}2k^{\prime \prime })+2,e^{4x_0})}{k+k^{}2k^{\prime \prime }4i}=\underset{n=1}{\overset{\mathrm{}}{}}\frac{ne^{4(n1)x_0}}{k+k^{}2k^{\prime \prime }4ni}.$$ (7.11) Similarly one can derive the infinite series forms of the other hypergeometric functions as $$\frac{F(3,\frac{i}{4}(k+k^{}2k^{\prime \prime })+2,\frac{i}{4}(k+k^{}2k^{\prime \prime })+3,e^{4x_0})}{k+k^{}2k^{\prime \prime }8i}=\frac{1}{2!}\underset{n=2}{\overset{\mathrm{}}{}}\frac{n(n1)e^{4(n2)x_0}}{k+k^{}2k^{\prime \prime }4ni},$$ (7.12) $`{\displaystyle \frac{F(4,\frac{i}{4}(k+k^{}2k^{\prime \prime })+3,\frac{i}{4}(k+k^{}2k^{\prime \prime })+4,e^{4x_0})}{k+k^{}2k^{\prime \prime }12i}}`$ $`={\displaystyle \frac{1}{3!}}{\displaystyle \underset{n=3}{\overset{\mathrm{}}{}}}{\displaystyle \frac{n(n1)(n2)e^{4(n3)x_0}}{k+k^{}2k^{\prime \prime }4ni}},`$ (7.13) $`{\displaystyle \frac{F(5,\frac{i}{4}(k+k^{}2k^{\prime \prime })+4,\frac{i}{4}(k+k^{}2k^{\prime \prime })+5,e^{4x_0})}{k+k^{}2k^{\prime \prime }16i}}`$ $`={\displaystyle \frac{1}{4!}}{\displaystyle \underset{n=4}{\overset{\mathrm{}}{}}}{\displaystyle \frac{n(n1)(n2)(n3)e^{4(n4)x_0}}{k+k^{}2k^{\prime \prime }4ni}}`$ (7.14) and $`{\displaystyle \frac{F(6,\frac{i}{4}(k+k^{}2k^{\prime \prime })+5,\frac{i}{4}(k+k^{}2k^{\prime \prime })+6,e^{4x_0})}{k+k^{}2k^{\prime \prime }20i}}`$ $`={\displaystyle \frac{1}{5!}}{\displaystyle \underset{n=5}{\overset{\mathrm{}}{}}}{\displaystyle \frac{n(n1)(n2)(n3)(n4)e^{4(n5)x_0}}{k+k^{}2k^{\prime \prime }4ni}}.`$ (7.15) Let us substitute (7.10), (7.11), (7.12), (7.13), (7.14) and (7.15) in (7.9) and obviously what remains in connection with the contribution (7.8), are the integrations over $`k^{\prime \prime }`$, $`k^{}`$ and $`k`$ . As before in previous sections, in order to integrate over the momenta $`k`$ and $`k^{}`$, it is sufficient to close the contours in the upper half-plane and pick up poles at $`\widehat{k}=k=k^{}=\sqrt{\omega ^24}`$ as all the other poles’ contributions will be exponentially damped when $`x,x^{}\mathrm{}`$ . Meanwhile, the integration over $`k^{\prime \prime }`$ is of the form $$\frac{dk^{\prime \prime }}{\sqrt{k^{\prime \prime 2}+4}}\frac{1}{\left(ik^{\prime \prime }2\mathrm{cos}\frac{(a_0+a_1)\pi }{2}\right)}\frac{1}{\left(k+k^{}2k^{\prime \prime }4ni\right)}.$$ (7.16) Now to manipulate the integral, let us choose the contour in the upper half-plane, taking care of the branch cut which runs from $`k^{\prime \prime }=2i`$ to infinity along the imaginary axis. Clearly this integral reduces to the integrals along the branch cut and we obtain $`{\displaystyle \frac{dk^{\prime \prime }}{\sqrt{k^{\prime \prime 2}+4}}\frac{1}{\left(ik^{\prime \prime }2\mathrm{cos}\frac{(a_0+a_1)\pi }{2}\right)}\frac{1}{\left(k+k^{}2k^{\prime \prime }4ni\right)}}`$ $`={\displaystyle \frac{1}{\left(k+k^{}+4i\mathrm{cos}\frac{(a_0+a_1)\pi }{2}4ni\right)}}({\displaystyle \frac{\frac{(a_0+a_1)\pi }{2}}{\mathrm{sin}\frac{(a_0+a_1)\pi }{2}}}`$ $`{\displaystyle \frac{2i}{\sqrt{\frac{(k+k^{})^2}{4}+44n^22ni(k+k^{})}}}`$ $`\times \mathrm{ln}\left\{{\displaystyle \frac{1+\frac{i(k+k^{})}{4}+n+\frac{i}{2}\sqrt{\frac{(k+k^{})^2}{4}+44n^22ni(k+k^{})}}{1+\frac{i(k+k^{})}{4}+n\frac{i}{2}\sqrt{\frac{(k+k^{})^2}{4}+44n^22ni(k+k^{})}}}\right\}).`$ (7.17) When $`n=0`$, the above formula is simplified much more, especially after doing the integration over $`k`$ and $`k^{}`$ and using the fact that $`\widehat{k}=k=k^{}=2\mathrm{sinh}\theta `$ . So the following formula can be obtained $`{\displaystyle \frac{dk^{\prime \prime }}{\sqrt{k^{\prime \prime 2}+4}}\frac{1}{\left(ik^{\prime \prime }2\mathrm{cos}\frac{(a_0+a_1)\pi }{2}\right)}\frac{1}{\left(2\widehat{k}2k^{\prime \prime }\right)}}`$ $`={\displaystyle \frac{1}{\left(2\widehat{k}+4i\mathrm{cos}\frac{(a_0+a_1)\pi }{2}\right)}}\left({\displaystyle \frac{\frac{(a_0+a_1)\pi }{2}}{\mathrm{sin}\frac{(a_0+a_1)\pi }{2}}}+{\displaystyle \frac{2}{\sqrt{\widehat{k}^2+4}}}({\displaystyle \frac{\pi }{2}}i\theta )\right).`$ (7.18) Now, (7.8) or more generally the contribution $`𝒟_2`$ can be written as: $`𝒟_2`$ $`=`$ $`{\displaystyle \frac{i\beta ^2}{4\pi }}\mathrm{cot}{\displaystyle \frac{a_0\pi }{2}}\mathrm{cot}{\displaystyle \frac{a_1\pi }{2}}\left(\mathrm{tan}^2{\displaystyle \frac{(a_0+a_1)\pi }{4}}\mathrm{cot}^2{\displaystyle \frac{(a_0+a_1)\pi }{4}}\right)`$ (7.19) $`\times {\displaystyle }{\displaystyle \frac{d\omega }{2\pi }}e^{i\omega (t_1t_2)}e^{i\widehat{k}(x_1+x_2)}\left({\displaystyle \frac{1}{\widehat{k}}}\right)^2f(\widehat{k},x_1)f(\widehat{k},x_2)`$ $`\times \{{\displaystyle \frac{2i\left(1\mathrm{cos}\frac{(a_0+a_1)\pi }{2}\right)^2}{\left(\widehat{k}+2i\mathrm{cos}\frac{(a_0+a_1)\pi }{2}\right)}}({\displaystyle \frac{\frac{(a_0+a_1)\pi }{2}}{\mathrm{sin}\frac{(a_0+a_1)\pi }{2}}}{\displaystyle \frac{2}{\sqrt{\widehat{k}^2+4}}}({\displaystyle \frac{\pi }{2}}i\theta ))`$ $`{\displaystyle \frac{1}{(i\widehat{k}+2)^2}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{e^{4nx_0}}{\left(2\widehat{k}+4i\mathrm{cos}\frac{(a_0+a_1)\pi }{2}4ni\right)}}({\displaystyle \frac{\frac{(a_0+a_1)\pi }{2}}{\mathrm{sin}\frac{(a_0+a_1)\pi }{2}}}`$ $`+{\displaystyle \frac{2i}{\sqrt{\widehat{k}^2+44n^24ni\widehat{k}}}}\mathrm{ln}\left\{{\displaystyle \frac{1+\frac{i}{2}\widehat{k}+n+\frac{i}{2}\sqrt{\widehat{k}^2+44n^24ni\widehat{k}}}{1+\frac{i}{2}\widehat{k}+n\frac{i}{2}\sqrt{\widehat{k}^2+44n^24ni\widehat{k}}}}\right\})`$ $`((A_1^{}\widehat{k}^2+2B_1^{}\widehat{k}+C_1^{})+n(A_2^{}\widehat{k}^2+2B_2^{}\widehat{k}+C_2^{})+{\displaystyle \frac{n(n1)}{2!}}(A_3^{}\widehat{k}^2+2B_3^{}\widehat{k}+C_3^{})`$ $`+{\displaystyle \frac{n(n1)(n2)}{3!}}(A_4^{}\widehat{k}^2+2B_4^{}\widehat{k}+C_4^{})+{\displaystyle \frac{n(n1)(n2)(n3)}{4!}}(2B_5^{}\widehat{k}+C_5^{})`$ $`+{\displaystyle \frac{n(n1)(n2)(n3)(n4)}{5!}}C_6^{})\}`$ $`+\text{other pole pieces}.`$ Firstly, in order to check the above solution, if we set $`a_0=a_1`$ and consider the other pole pieces then, we can derive the formula (3.10) in reference . As we mentioned before, the calculation of this reference is based on the case when the boundary parameters are equal. Secondly, in this solution, we verified that the term which depends explicitly on the rapidity of the particle ($`\theta `$) is cancelled by counterpart terms in the other pole pieces. It is evident that if we add the expressions (7.7) and (7.19) then, the contribution (7.6) will be obtained i.e. $`𝒟=𝒟_1+𝒟_2`$. ## 8 Discussion Affine Toda field theory on the whole line is an exactly solvable theory for which the S-matrices have been formulated. However, when a boundary is present then the boundary S-matrices of the theory i.e. the reflection factors, have not been completely found. The bootstrap technique does not uniquely determine the reflection factors. Fortunately perturbation theory provides the link between the expressions for the reflection factors which come from the bootstrap equations and the boundary parameters. Nevertheless, this method normally involves complicated calculations. In this paper the quantum reflection factor for the $`a_1^{(1)}`$ affine Toda field theory or sinh-Gordon model with integrable boundary conditions has been studied in low order perturbation theory when $`\sigma _0\sigma _1`$. It is found that at one loop order the quantum corrections to the classical reflection factor of the model can be expressed in terms of hypergeometric functions for most of the related Feynman diagrams. Although there is still some work to do to calculate the contributions of the remaining diagrams, it is understood that the provided procedure and some formalisms may be followed for them. The calculations corresponding to the type II Feynman diagram which are not carried out in this paper, are more difficult than the others. In this case the two middle propagators are exactly the same and this fact influences the difficulty of the computations. However some formulae that have been presented here, could be helpful for the remaining diagram. For example, consider the contribution of the type II (boundary-bulk) diagram: $`2\beta ^2(\sigma _1\mathrm{coth}x_0\sigma _0\mathrm{tanh}x_0){\displaystyle 𝑑t𝑑t^{}𝑑xG(x_1,t_1;x,t)G(x,t;0,t^{})}`$ $`\times G(x,t;0,t^{})G(0,t^{};x_2,t_2)\mathrm{sinh}(\sqrt{2}\beta \varphi _0).`$ (8.1) Now as far as the integration over $`x`$ is concerned we should obtain the following integrals $$_{\mathrm{}}^0𝑑xe^{i(k+k^{}k_1)x}\mathrm{sinh}(\sqrt{2}\beta \varphi _0)\mathrm{coth}^n2(xx_0),$$ (8.2) where n=0,1,2,3. It is better to solve: $$_{\mathrm{}}^0𝑑x\mathrm{exp}\left\{\tau +i(k+k^{}k_1)x\right\}\mathrm{sinh}(\sqrt{2}\beta \varphi _0)\mathrm{coth}^n2(xx_0)$$ (8.3) in which $`\tau `$ is a small positive quantity and will be taken to zero later. In fact, the relation (8.3) is very similar to the formula (A.1) in Appendix A. So, following the same procedure that have been followed in Appendix A, one can find the solution of (8.2) when $`n=0`$ as: $`{\displaystyle _{\mathrm{}}^0}`$ $`dx`$ $`e^{i(k+k^{}k_1)x}\mathrm{sinh}(\sqrt{2}\beta \varphi _0)`$ (8.4) $`=`$ $`{\displaystyle \frac{1}{\mathrm{sinh}2x_0}}`$ $`{\displaystyle \frac{2(k+k^{}k_1)}{k+k^{}k_12i}}`$ $`\times e^{2x_0}F(1,{\displaystyle \frac{i}{4}}(k+k^{}k_1)+{\displaystyle \frac{1}{2}},{\displaystyle \frac{i}{4}}(k+k^{}k_1)+{\displaystyle \frac{3}{2}},e^{4x_0}).`$ Then, the solutions of (8.2) for $`n=1,2,3`$ can be derived exactly in according to the Appendix A terms. But, this is not all of the problem. As we mentioned before, in type II diagram double Green functions cause the middle momenta to be linked to each other in a complicated way and the calculations become more intricate. Actually this diagram must be studied in three cases depending on the interaction vertices being located in the bulk region or at the boundary. Moreover because of the symmetry, the contribution of the type II (boundary-bulk) diagram is the same as the type II (bulk-boundary) one. When the boundary parameters are equal then, only the type I diagram is involved in the theory. As we mentioned before, in this special case the quantum corrections to the classical reflection of the model have been found and Ghoshal’s formula for the lightest breather is checked perturbatively to $`O(\beta ^2)`$. In our case, we realised that the contribution of the type I (bulk) reduces to the special case. Taking (7.7) and (7.19) expressions, if we put $`\sigma _0=\sigma _1`$ then, we obtain the same result as reference and this is a check on our calculations. Moreover, when $`\sigma _0\sigma _1`$ the following expressions for $`E`$ and $`F`$ in Ghoshal’s formula (2.8) have been conjectured to be: $$E=(a_0+a_1)(1B/2)F=(a_0a_1)(1B/2).$$ (8.5) So, it will be interesting to check the above conjecture after finding the contributions of the remaining diagrams and adding the results all together. This will lead to a deeper understanding of the quantum integrability of the theory. However, it is necessary to find simplifications of the contributions when they add among themselves in order to get Ghoshal’s formula. ## 9 Acknowledgement We would like to thank E. Corrigan and P. Bowcock for encouragement, discussions and suggestions, and the Ministry of Culture and Higher Education of Iran for financial support. ## Appendix A In this Appendix we obtain such integrals $$S_n=_{\mathrm{}}^0𝑑xe^{(2+ik)x}\mathrm{sinh}(\sqrt{2}\beta \varphi _0)\mathrm{coth}^n2(xx_0)$$ (A.1) in which $`n=0,1,2,3`$, $`\varphi _0`$ is the background solution to the equation of field so that $`\mathrm{sinh}(\sqrt{2}\beta \varphi _0)`$ is proportional to the bulk three point coupling which is given by $$\mathrm{sinh}(\sqrt{2}\beta \varphi _0)=2\mathrm{cosh}2(xx_0)\left(\mathrm{coth}^22(xx_0)1\right).$$ (A.2) Let us start with the simplest case when $`n=0`$. Using (A.2), we have $$S_0=_{\mathrm{}}^0e^{(2+ik)x}d\left(\frac{1}{\mathrm{sinh}2(xx_0)}\right)$$ (A.3) or, after integration by parts $$S_0=\frac{1}{\mathrm{sinh}2x_0}+2(2+ik)_{\mathrm{}}^0𝑑xe^{(2+ik)x}\frac{1}{e^{2(xx_0)}e^{2(xx_0)}}.$$ (A.4) Now, according to (3.6), if $`\sigma _0>\sigma _1`$ then $`x_00`$. Otherwise, it is necessary to adjust the background solution (3.5) by shifting $`x_0`$ through $`i\pi /2`$ in order to be guaranteed that $`x_00`$. The singularity in the equation (3.5) is unimportant provided $`x_0`$ is positive. So, from now on it is assumed $`\sigma _0\sigma _1`$. Therefore, $`x_0`$ is greater or equal to zero. But, $`x`$ is less than zero, so $`0<e^{4(xx_0)}<1`$ and hence $$\frac{1}{e^{2(xx_0)}e^{2(xx_0)}}=e^{2(xx_0)}\underset{n=0}{\overset{\mathrm{}}{}}e^{4n(xx_0)}.$$ (A.5) Substituting (A.5) in (A.4), we obtain $$S_0=\frac{1}{\mathrm{sinh}2x_0}2(2+ik)e^{2x_0}_{\mathrm{}}^0𝑑xe^{(4+ik)x}\underset{n=0}{\overset{\mathrm{}}{}}e^{4n(xx_0)}.$$ (A.6) Clearly, the series is uniformly convergent so the above relation becomes $$S_0=\frac{1}{\mathrm{sinh}2x_0}2(2+ik)e^{2x_0}\underset{n=0}{\overset{\mathrm{}}{}}e^{4nx_0}_{\mathrm{}}^0𝑑xe^{(4+4n+ik)x}.$$ (A.7) After integration over $`x`$, we obtain $$S_0=\frac{1}{\mathrm{sinh}2x_0}+2i(2+ik)e^{2x_0}\underset{n=0}{\overset{\mathrm{}}{}}\frac{e^{4nx_0}}{k(4+4n)i}.$$ (A.8) On the other hand, the above infinite series is a hypergeometric function. That is $$\underset{n=0}{\overset{\mathrm{}}{}}\frac{e^{4nx_0}}{ki(4+4n)}=\frac{F(1,\frac{i}{4}k+1,\frac{i}{4}k+2,e^{4x_0})}{k4i}.$$ (A.9) Therefore, we get the following relation $`S_0`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{sinh}2x_0}}`$ (A.10) $`2{\displaystyle \frac{k2i}{k4i}}e^{2x_0}F(1,{\displaystyle \frac{i}{4}}k+1,{\displaystyle \frac{i}{4}}k+2,e^{4x_0}).`$ The hypergeometric function is defined by $$F(a,b,c,z)=\underset{n=0}{\overset{\mathrm{}}{}}\frac{(a)_n(b)_n}{(c)_n}z^nc0,1,2,\mathrm{}$$ (A.11) where $$(a)_n=\frac{\mathrm{\Gamma }(a+n)}{\mathrm{\Gamma }(a)}=a(a+1)\mathrm{}(a+n1)n=1,2,3,\mathrm{}.$$ (A.12) The above series defines a function which is analytic when $`\|z\|<1`$. Also, the derivatives of the hypergeometric function are given by $$\frac{d^n}{dz^n}F(a,b,c,z)=\frac{(a)_n(b)_n}{(c)_n}F(a+n,b+n,c+n,z).$$ (A.13) Next, for $`n=1`$ using (A.2) for the bulk three point coupling we have $$S_1=2_{\mathrm{}}^0𝑑xe^{(2+ik)x}\mathrm{cosh}2(xx_0)\mathrm{coth}2(xx_0)\left(\mathrm{coth}^22(xx_0)1\right).$$ (A.14) On the other hand, if we differentiate the left hand side of (A.10) with respect to $`x_0`$, which is given by $`{\displaystyle \frac{S_0}{x_0}}={\displaystyle _{\mathrm{}}^0}dxe^{(2+ik)x}(4\mathrm{sinh}2(xx_0)(\mathrm{coth}^22(xx_0)1).`$ $`.+8\mathrm{cosh}2(xx_0)\mathrm{coth}2(xx_0)(1\mathrm{coth}^22(xx_0)))`$ (A.15) and by comparing the above formula with (A.14), then the following equation may be derived $$S_1=\frac{1}{4}\frac{S_0}{x_0}+_{\mathrm{}}^0𝑑xe^{(2+ik)x}\frac{1}{\mathrm{sinh}2(xx_0)}.$$ (A.16) The second term in the above relation can be manipulated as before. Moreover, it is evident that we need to differentiate the right hand side of (A.10) which is equal to $`{\displaystyle \frac{S_0}{x_0}}`$ $`=`$ $`{\displaystyle \frac{2\mathrm{cosh}2x_0}{\mathrm{sinh}^22x_0}}`$ (A.17) $`+4{\displaystyle \frac{k2i}{k4i}}e^{2x_0}F(1,{\displaystyle \frac{i}{4}}k+1,{\displaystyle \frac{i}{4}}k+2,e^{4x_0})`$ $`+8{\displaystyle \frac{k2i}{k8i}}e^{6x_0}F(2,{\displaystyle \frac{i}{4}}k+2,{\displaystyle \frac{i}{4}}k+3,e^{4x_0}).`$ Finally, by substituting the relation (A.17) in (A.16), doing the computation of second term in the right-hand side of (A.16) and after simplifying we obtain $`S_1`$ $`=`$ $`{\displaystyle \frac{\mathrm{cosh}2x_0}{2\mathrm{sinh}^22x_0}}`$ (A.18) $`+{\displaystyle \frac{k}{k4i}}e^{2x_0}F(1,{\displaystyle \frac{i}{4}}k+1,{\displaystyle \frac{i}{4}}k+2,e^{4x_0})`$ $`+2{\displaystyle \frac{k2i}{k8i}}e^{6x_0}F(2,{\displaystyle \frac{i}{4}}k+2,{\displaystyle \frac{i}{4}}k+3,e^{4x_0}).`$ In the same way, we may derive (A.1) when n is equal to 2 or 3, however gradually the calculations become lengthy and we only write down the results, that is $`S_2`$ $`=`$ $`{\displaystyle \frac{2}{3}}{\displaystyle \frac{1}{\mathrm{sinh}2x_0}}+{\displaystyle \frac{1}{6}}{\displaystyle \frac{\mathrm{cosh}^22x_0+1}{\mathrm{sinh}^32x_0}}`$ (A.19) $`{\displaystyle \frac{1}{3}}{\displaystyle \frac{5k8i}{k4i}}e^{2x_0}F(1,{\displaystyle \frac{i}{4}}k+1,{\displaystyle \frac{i}{4}}k+2,e^{4x_0})`$ $`{\displaystyle \frac{4}{3}}{\displaystyle \frac{2k3i}{k8i}}e^{6x_0}F(2,{\displaystyle \frac{i}{4}}k+2,{\displaystyle \frac{i}{4}}(k+3,e^{4x_0})`$ $`{\displaystyle \frac{8}{3}}{\displaystyle \frac{k2i}{k12i}}e^{10x_0}F(3,{\displaystyle \frac{i}{4}}k+3,{\displaystyle \frac{i}{4}}k+4,e^{4x_0})`$ and $`S_3`$ $`=`$ $`{\displaystyle \frac{13}{24}}{\displaystyle \frac{\mathrm{cosh}2x_0}{\mathrm{sinh}^22x_0}}{\displaystyle \frac{1}{24}}{\displaystyle \frac{\mathrm{cosh}^32x_0+5\mathrm{cosh}2x_0}{\mathrm{sinh}^42x_0}}`$ (A.20) $`+{\displaystyle \frac{1}{6}}{\displaystyle \frac{7k4i}{k4i}}e^{2x_0}F(1,{\displaystyle \frac{i}{4}}k+1,{\displaystyle \frac{i}{4}}k+2,e^{4x_0})`$ $`+{\displaystyle \frac{1}{3}}{\displaystyle \frac{13k22i}{k8i}}e^{6x_0}F(2,{\displaystyle \frac{i}{4}}k+2,{\displaystyle \frac{i}{4}}k+3,e^{4x_0})`$ $`+{\displaystyle \frac{1}{3}}{\displaystyle \frac{18k32i}{k12i}}e^{10x_0}F(3,{\displaystyle \frac{i}{4}}k+3,{\displaystyle \frac{i}{4}}k+4,e^{4x_0})`$ $`+4{\displaystyle \frac{k2i}{k16i}}e^{14x_0}F(4,{\displaystyle \frac{i}{4}}k+4,{\displaystyle \frac{i}{4}}k+5,e^{4x_0}).`$ ## Appendix B In this Appendix we find the following integrals $$C_n=_{\mathrm{}}^0𝑑xe^{ikx}\mathrm{cosh}(\sqrt{2}\beta \varphi _0)\mathrm{coth}^n2(xx_0).$$ (B.1) Here, $`\mathrm{cosh}(\sqrt{2}\beta \varphi _0)`$ is proportional to the bulk four point coupling and is given by $$\mathrm{cosh}(\sqrt{2}\beta \varphi _0)=\left(2\mathrm{coth}^22(xx_0)1\right).$$ (B.2) So, let us calculate such integrals $$I_n=_{\mathrm{}}^0𝑑xe^{ikx}\mathrm{coth}^n2(xx_0),$$ (B.3) where $`n=1,2,\mathrm{},6`$ . It is better to find the solution of the above integrals when $`n=1`$. Considering the following inequality (see Appendix A) $$0<e^{4(xx_0)}<1$$ and therefore, in what follows we will use the expanded form of $`\mathrm{coth}2(xx_0)`$ as $$\mathrm{coth}2(xx_0)=12\underset{n=0}{\overset{\mathrm{}}{}}e^{4n(xx_0)}.$$ (B.4) It turns out to be simple if we consider this integral $$_{\mathrm{}}^0𝑑xe^{(\tau +ik)x}\mathrm{coth}2(xx_0),$$ (B.5) where $`\tau `$ is a positive constant quantity which will be taken to zero at the end of the calculation. Moreover, by using (B.4), then (B.5) becomes $$_{\mathrm{}}^0𝑑xe^{(\tau +ik)x}2_{\mathrm{}}^0𝑑x\underset{n=0}{\overset{\mathrm{}}{}}e^{4n(xx_0)}e^{(\tau +ik)x}$$ (B.6) and regarding the arguments in Appendix A, we may evaluate the above relation to obtain $$\frac{i}{ki\tau }+2\underset{n=0}{\overset{\mathrm{}}{}}\frac{i}{k(\tau +4n)i}e^{4nx_0}.$$ (B.7) Now, we can write down the desired result, that is, $$I_1=\frac{i}{k}+2i\underset{n=0}{\overset{\mathrm{}}{}}\frac{e^{4nx_0}}{k4ni}.$$ (B.8) On the other hand, the above series is equal to a hypergeometric function $$\underset{n=0}{\overset{\mathrm{}}{}}\frac{e^{4nx_0}}{k4ni}=\frac{1}{k}F(1,\frac{i}{4}k,\frac{i}{4}k+1,e^{4x_0})$$ (B.9) and finally we find this formula $$I_1=\frac{i}{k}+\frac{2i}{k}F(1,\frac{i}{4}k,\frac{i}{4}k+1,e^{4x_0}).$$ (B.10) Now, let us compute (B.3) when $`n=2`$ and in order to solve it, it is sufficient to differentiate both sides of (B.10) with respect to $`x_0`$ to obtain $$I_2=\frac{i}{k}\frac{4i}{k4i}e^{4x_0}F(2,\frac{i}{4}k+1,\frac{i}{4}k+2,e^{4x_0}).$$ (B.11) We can follow a similar method to get higher order forms of (B.3) which we need in this paper so, it is appropriate to write down all of them i.e. $`I_3`$ $`=`$ $`{\displaystyle \frac{i}{k}}+{\displaystyle \frac{2i}{k}}F(1,{\displaystyle \frac{i}{4}}k,{\displaystyle \frac{i}{4}}k+1,e^{4x_0})`$ (B.12) $`{\displaystyle \frac{4i}{k4i}}e^{4x_0}F(2,{\displaystyle \frac{i}{4}}k+1,{\displaystyle \frac{i}{4}}k+2,e^{4x_0})`$ $`{\displaystyle \frac{8i}{k8i}}e^{8x_0}F(3,{\displaystyle \frac{i}{4}}k+2,{\displaystyle \frac{i}{4}}k+3,e^{4x_0}),`$ $`I_4`$ $`=`$ $`{\displaystyle \frac{i}{k}}{\displaystyle \frac{8i}{k4i}}e^{4x_0}F(2,{\displaystyle \frac{i}{4}}k+1,{\displaystyle \frac{i}{4}}k+2,e^{4x_0})`$ (B.13) $`{\displaystyle \frac{16i}{k8i}}e^{8x_0}F(3,{\displaystyle \frac{i}{4}}k+2,{\displaystyle \frac{i}{4}}k+3,e^{4x_0})`$ $`{\displaystyle \frac{16i}{k12i}}e^{12x_0}F(4,{\displaystyle \frac{i}{4}}k+3,{\displaystyle \frac{i}{4}}k+4,e^{4x_0}),`$ $`I_5`$ $`=`$ $`{\displaystyle \frac{i}{k}}+{\displaystyle \frac{2i}{k}}F(1,{\displaystyle \frac{i}{4}}k,{\displaystyle \frac{i}{4}}k+1,e^{4x_0})`$ (B.14) $`{\displaystyle \frac{8i}{k4i}}e^{4x_0}F(2,{\displaystyle \frac{i}{4}}k+1,{\displaystyle \frac{i}{4}}k+2,e^{4x_0})`$ $`{\displaystyle \frac{32i}{k8i}}e^{8x_0}F(3,{\displaystyle \frac{i}{4}}k+2,{\displaystyle \frac{i}{4}}k+3,e^{4x_0})`$ $`{\displaystyle \frac{48i}{k12i}}e^{12x_0}F(4,{\displaystyle \frac{i}{4}}k+3,{\displaystyle \frac{i}{4}}k+4,e^{4x_0})`$ $`{\displaystyle \frac{32i}{k16i}}e^{16x_0}F(5,{\displaystyle \frac{i}{4}}k+4,{\displaystyle \frac{i}{4}}k+5,e^{4x_0})`$ and $`I_6`$ $`=`$ $`{\displaystyle \frac{i}{k}}{\displaystyle \frac{12i}{k4i}}e^{4x_0}F(2,{\displaystyle \frac{i}{4}}k+1,{\displaystyle \frac{i}{4}}k+2,e^{4x_0})`$ (B.15) $`{\displaystyle \frac{48i}{k8i}}e^{8x_0}F(3,{\displaystyle \frac{i}{4}}k+2,{\displaystyle \frac{i}{4}}k+3,e^{4x_0})`$ $`{\displaystyle \frac{112i}{k12i}}e^{12x_0}F(4,{\displaystyle \frac{i}{4}}k+3,{\displaystyle \frac{i}{4}}k+4,e^{4x_0})`$ $`{\displaystyle \frac{128i}{k16i}}e^{16x_0}F(5,{\displaystyle \frac{i}{4}}k+4,{\displaystyle \frac{i}{4}}k+5,e^{4x_0})`$ $`{\displaystyle \frac{64i}{k20i}}e^{20x_0}F(6,{\displaystyle \frac{i}{4}}k+5,{\displaystyle \frac{i}{4}}k+6,e^{4x_0}).`$
warning/0004/hep-ph0004252.html	ar5iv	text	# 𝑁^∗ Spectrum in Lattice QCD 22footnote 2This work was done in collaboration with T. Blum, P. Chen, N. Christ, M. Creutz, C. Dawson, G. Fleming, R. Mawhinney, S. Ohta, S. Sasaki, G. Siegert, A. Soni, P. Vranas, M. Wingate, L. Wu and Y. Zhestkov (RIKEN-BNL-Columbia-KEK collaboration). ## 1 Introduction It is important to understand the hadron mass spectrum from the first principles of quantum chromodynamics (QCD), the fundamental theory of the strong interactions. As is well known, the only method for such a first-principle calculation is numerical lattice QCD. Recent lattice calculations of the light-hadron mass spectrum in the quenched approximation agree well with experimental values within about 5% $`^\mathrm{?}`$. However, this success is restricted to ground states and does not apply to excited states. The excited-state mass spectrum is big challenge in lattice QCD. In this study, we focus on a negative-parity nucleon excited-state $`N^{}(1535)`$, which is theoretically identified as the parity partner of the nucleon $`N(939)`$. We are especially interested in the mass splitting between $`N`$ and $`N^{}`$ from the viewpoint of parity partners. As is well known, the mass splitting between parity partners would be absent if chiral symmetry were preserved $`^\mathrm{?}`$. In other words spontaneous chiral-symmetry breaking is responsible for the absence of such parity doubling in the actual hadron spectrum. In this sense, regardless of a model or a theory, it is important to handle the chiral symmetry and its spontaneous breaking for reproducing precisely the mass splitting between parity-partner hadrons. Indeed, both non-relativistic quark models $`^\mathrm{?}`$ and bag models $`^\mathrm{?}`$, which explicitly break chiral symmetry, fail to reproduce the large mass splitting between $`N(939)`$ and $`N^{}(1535)`$ $`^\mathrm{?}`$. The non-relativistic quark models are based on a harmonic oscillator description of the orbital motion of constituent quarks. The plausible value of its oscillator quantum should be a couple of hundred MeV to reproduce the observed charge radius and magnetic moment of the nucleon $`^\mathrm{?}`$. Since this model regards $`N^{}`$ as a state with one quantum excitation in orbital motion, it indicates that $`N^{}`$ lies just a few hundred MeV above the ground state. Even worse, we face a serious problem of the wrong ordering of $`N^{}(1535)`$ and the positive-parity excited state $`N^{}(1440)`$ $`^\mathrm{?}`$ because $`N^{}`$ should be assigned two oscillator quanta in this model $`^\mathrm{?}`$. Furthermore this wrong ordering cannot be improved by the conventional one-gluon-exchange potential model of the residual interaction between constituent quarks $`^\mathrm{?}`$. In addition, this difficulty led Glozman and Riska to propose another candidate for the residual interaction $`^\mathrm{?}`$. It is easy to see that the bag models face essentially the same problem $`^\mathrm{?}`$. It is clearly an interesting question whether lattice QCD can reproduce this large $`N`$-$`N^{}`$ mass splitting. However, conventional lattice fermion schemes had some difficulty in this challenge. The Nielsen-Ninomiya no-go theorem $`^\mathrm{?}`$ dictates that either chiral symmetry or flavor symmetry or both are supposed to be violated at finite lattice spacing, while both of them are essential in this subject. The Wilson fermions explicitly break chiral symmetry at finite lattice spacing and hence are quite inadequate for the current problem. Although Kogut-Susskind (KS) fermions have a remnant $`U(1)`$ chiral symmetry at finite lattice spacing, they are still not capable of the $`N^{}`$ mass calculation. The main reason is that KS fermion has only discrete flavor symmetries which belong to a subgroup of the $`SU(4)`$ flavor symmetry $`^\mathrm{?}`$. There are only three irreducible representations, 8, 8’ and 16 for KS baryon operators due to this incomplete flavor symmetry. Two appropriate representations 8 and 16, to which $`N^{}`$(1535) belongs, involve also $`\mathrm{\Lambda }`$(1405), $`\mathrm{\Lambda }`$(1520) and $`N`$(1520) $`^\mathrm{?}`$. The study of $`N^{}`$ spectrum using KS fermions always faces this inevitable contamination from lower mass states. Several years ago, Kaplan $`^\mathrm{?}`$ advocated a new type of lattice fermion scheme, and Shamir $`^{\mathrm{?},\mathrm{?}}`$ reformulated it for lattice QCD simulations. They are called domain wall fermions (DWF), which utilize a fictitious fifth dimension. An important feature of DWF is that chiral symmetry is almost exactly preserved even at the finite lattice spacing. This is achieved because the symmetry violation is suppressed exponentially in terms of the number of lattice sites $`N_s`$ in the extra dimension $`^\mathrm{?}`$. In other words, $`N_s`$ gives us a way to control the violation. Recent lattice calculations with DWF have shown that good chiral properties are obtained for moderate $`N_s`$ like 16 if the lattice spacing is small enough, like 0.2 fm $`^\mathrm{?}`$. In addition, the flavor symmetry is also well preserved in this fermion discretization. Thus we are led to an attempt to use DWF for lattice QCD calculations of the $`N^{}`$ mass spectrum $`^\mathrm{?}`$. ## 2 Baryon operators The mass $`m_B`$ of the low-lying baryon $`B`$ is extracted from the two-point correlation function composed of the baryon interpolating operator $`𝒪_B`$: $$G_{𝒪_B}(t)=\underset{\stackrel{}{x}}{}0\|T\{𝒪_B(\stackrel{}{x},t)\overline{𝒪}_B(0,0)\}\|0,$$ (1) which behaves like $`\mathrm{exp}(m_Bt)`$ for large $`t`$ since $`G_{𝒪_B}(t)`$ is projected out at zero spatial momentum through the sum over $`\stackrel{}{x}`$ $`^\mathrm{?}`$. We focus on the nucleon channel, the spin-half iso-doublet baryons. There are two possible interpolating operators for the corresponding quantum number $`J^P=1/2^+`$ even if we restrict them to contain no derivatives and to belong to the $`(\frac{1}{2},0)(0,\frac{1}{2})`$ chiral multiplet under $`SU(2)_LSU(2)_R`$ $`^\mathrm{?}`$: $`B_1^+`$ $`=`$ $`\epsilon _{abc}(u_a^TC\gamma _5d_b)u_c,`$ (2) $`B_2^+`$ $`=`$ $`\epsilon _{abc}(u_a^TCd_b)\gamma _5u_c,`$ (3) where $`abc`$, $`ud`$, $`C`$ and $`\gamma _5`$ have usual meanings as color, flavor, charge conjugation and Dirac indices. The superscript “$`+`$” denotes the positive parity. The operator $`B_1^+`$ is preferred for use in lattice QCD to extract the signal of the nucleon ground-state. On the other hand, the operator $`B_2^+`$ is conventionally discarded since $`B_2^+`$ is expected to couple weakly to the nucleon ground-state due to the vanishing non-relativistic limit $`^\mathrm{?}`$. In our calculation, the nucleon interpolating operator is assigned to $`B_1^+`$ in the conventional way. We try to calculate the excited-state mass spectrum from $`B_2^+`$. Multiplying the left hand side of the previous positive parity operators by $`\gamma _5`$, we obtain the interpolating operators with negative parity $`^\mathrm{?}`$: $`B_1^{}`$ $`=`$ $`\gamma _5B_1^+=\epsilon _{abc}(u_a^TC\gamma _5d_b)\gamma _5u_c,`$ (4) $`B_2^{}`$ $`=`$ $`\gamma _5B_2^+=\epsilon _{abc}(u_a^TCd_b)u_c.`$ (5) The point to notice is the relation between the nucleon two-point functions with opposite parities $$G_{B^+}(t)=\gamma _5G_B^{}(t)\gamma _5.$$ (6) This means that the two-point correlation function can couple to both positive and negative parity states $`^\mathrm{?}`$: $`G_{B^+}(t)`$ $``$ $`(1+\gamma _t)G_0(t)+(1\gamma _t)G_0(t),`$ (7) $`G_B^{}(t)`$ $``$ $`(1\gamma _t)G_0(t)(1+\gamma _t)G_0(t),`$ (8) where $`G_0(t)=\theta (t)A_+e^{M_+t}+\theta (t)A_{}e^{+M_{}t}`$ <sup>2</sup><sup>2</sup>2In the quenched approximation, $`N^{}`$ can be regarded as the stable baryon like the nucleon.. $`M_\pm `$ denote masses of the opposite parity state. Since $`G_0(t)`$ possesses the contribution from both positive and negative parity lowest-lying states $`^\mathrm{?}`$, which propagate in the opposite time direction respectively, we are threatened by contamination from the backward propagating opposite parity state in the case of the baryon spectrum with some boundary condition. An appropriate boundary condition is required in the time direction to reduce the contamination of the opposite parity state. At the end of this section, let us briefly review the knowledge that non-broken chiral symmetry imposes parity doubling in the hadron spectrum $`^{\mathrm{?},\mathrm{?}}`$. For the sake of simplicity, we consider a particular transformation of the $`SU(2)`$ chiral symmetry as $`[𝒬_5,u]=+i\gamma _5u`$ and $`[𝒬_5,d]=i\gamma _5d`$. Then, one can easily find that the two-point correlator composed of each $`B_1`$ and $`B_2`$ should transform as $$[𝒬_5,B_{1,2}^+(x)\overline{B}_{1,2}^+(0)]=i\{\gamma _5,B_{1,2}^+(x)\overline{B}_{1,2}^+(0)\}$$ (9) in the chiral limit; $`[𝒬_5,H]=0`$. Suppose that the vacuum possesses chiral symmetry: $`𝒬_5\|0=0`$. Eq.(9) shows that the two-point correlation function of the nucleon and $`\gamma _5`$ have to anti-commute, $$\{\gamma _5,G_{B^+}(t)\}=0.$$ (10) Immediately, Eq.(6) gives $$G_{B^+}(t)=G_B^{}(t),$$ (11) which means that parity doubling arises in the nucleon channel due to chiral symmetry $`^\mathrm{?}`$. Of course, the chiral symmetry is spontaneously broken, $`𝒬_5\|00`$ in the QCD vacuum so that such a parity doubling never occurs in the actual spectrum $`^\mathrm{?}`$. In this sense, the spontaneous breaking of chiral symmetry is responsible for the absence of parity doubling. <sup>3</sup><sup>3</sup>3This argument ignores the consequence of the ’t Hooft anomaly condition $`^\mathrm{?}`$ and also spoils the possibility of the massless baryon because $`G_{B^+}(t)`$ and $`G_B^{}(t)`$ are defined at zero spatial momentum. ## 3 Numerical results We generate quenched QCD configurations on a $`16^3\times 32`$ lattice with the standard single-plaquette Wilson action at $`\beta =6/g^2=6.0`$. The quark propagator is calculated by using domain wall fermions with a fifth dimension of $`N_s`$=16 sites and the domain-wall height $`M`$=1.8. According to our simulations $`^\mathrm{?}`$, this corresponds to lattice units $`a^11.9`$ GeV from $`am_\rho `$=0.400(8) in the chiral limit and spatial lattice size $`La1.7`$ fm. In this work, we provide forward-type and backward-type quark propagators with the wall source at two different locations ($`t`$ = 5 and 27) and take the average of two measurements for the hadron spectrum in each configuration. We use 205 independent gauge configurations for the lightest two quark masses, $`m=0.02`$ and $`0.03`$ and 105 configurations for the heavier ones, $`m=0.040.125`$. Those bare quark masses correspond to mass ratios $`m_\pi /m_\rho 0.590.90`$. All calculations were done on the 600 Gflops QCDSP machine at the RIKEN BNL Research Center. Now let us touch upon technical details. We take a linear combination of two quark propagators with periodic and anti-periodic boundary condition in the time direction, respectively. This procedure enables us to extract the state with desired parity even in the baryon spectrum. ### 3.1 Parity partner of nucleon: $`N^{}`$ We present mass estimates of the nucleon ($`N`$) and its parity partner ($`N^{}`$) obtained by the single mass-fit method applied to the two-point functions. We first calculate the effective masses to find appropriate time ranges for fitting. The effective mass plot shows plateaus on the far side of the source ($`tt_{\mathrm{source}}`$=13-20) for $`N`$ and on the near side of the source ($`tt_{\mathrm{source}}`$=5-12) for $`N^{}`$. We choose to fit our baryon propagators from some minimum time slice $`T_{\mathrm{min}}`$ to an appropriate maximum time slice $`T_{\mathrm{max}}=20`$ for $`N`$ and $`T_{\mathrm{max}}16`$ for $`N^{}`$. To keep fitting ranges as wide as possible, $`T_{\mathrm{min}}`$ is varied from $`T_{\mathrm{max}}2`$ and selected under the condition $`\chi ^2/N_{\mathrm{DF}}1.0`$ where $`N_{\mathrm{DF}}`$ denotes the fitted degree of freedom. All our fits are reasonable in the sense that the confidence-level is larger than 0.3 and estimates from the weighted average of the effective mass agree with them within errors. In Figure 1 we show the low-lying nucleon spectrum as a function of the quark mass, $`m`$. The nucleon mass is extracted from the $`B_1^+`$ operator. We omit the point at $`m=0.02`$ for the operator $`B_2^{}`$ since a good plateau in the effective mass plot is absent. $`N^{}`$ mass estimates from $`B_1^{}`$ and $`B_2^{}`$ operators agree with each other within errors in the whole quark mass range, as expected from their common quantum numbers $`^\mathrm{?}`$. Note that this result disagrees with the $`N^{}`$ spectrum obtained in Ref 16: we cannot find any discrepancy between mass spectra from $`B_1^{}`$ and $`B_2^{}`$. Note also the same signal for $`N^{}`$ is obtained from a mixed correlation function $`0\|B_1^{}\overline{B}_2^{}+B_2^{}\overline{B}_1^{}\|0`$ in our study. The most remarkable feature in Figure 1 is that the $`N`$-$`N^{}`$ mass splitting is observed in the whole range and even for light valence quark mass values $`^\mathrm{?}`$. This mass splitting grows as the valence quark mass decreases $`^\mathrm{?}`$. To make this point clear, we compare two mass ratios, one from the baryon parity partners $`m_N^{}/m_N`$ and the other from pseudo-scalar and vector mesons $`m_\pi /m_\rho `$ in Figure 2. Experimental points are marked with stars, corresponding to non-strange (left) and strange (right) sectors. In the strange sector we use $`\mathrm{\Sigma }`$ and $`\mathrm{\Sigma }(1750)`$ as baryon parity partners and $`K`$ and $`K^{}`$ for mesons $`^\mathrm{?}`$. We find the baryon mass ratio clearly grows with decreasing meson mass ratio, toward the experimental values $`^\mathrm{?}`$. Finally, we evaluate the $`N`$ and $`N^{}`$ mass in the chiral limit. We simply take a linear extrapolation in 5 lightest quark masses for $`B_1^+`$ and $`B_1^{}`$. We find $`m_N`$=0.55(1) and $`m_N^{}`$=0.80(2) in lattice units for values in the chiral limit. If the scale is set as $`a^11.9`$ GeV from the $`\rho `$ mass $`^\mathrm{?}`$, we obtain $`m_N1.0`$ GeV and $`m_N^{}1.5`$ GeV in a good agreement with the experimental value. Above errors do not include systematic errors due to finite volume, lattice spacing and quenching effects. Such a systematic study will be addressed in future calculations. ### 3.2 Positive parity excited-state: $`N^{}`$ In principle, the mass spectrum of the “first” excited nucleon can be derived from the two-mass fits for the nucleon interpolating operator $`B_1^+`$. However, large statistics are required. Also it is difficult to control the systematic errors. Indeed, several attempts to extract $`N^{}`$ mass spectrum failed to reproduce experiment. Here, we take another approach to $`N^{}`$ spectrum. As we mentioned, the $`B_2^+`$ operator vanishes in the non-relativistic limit. Thus we expect that $`B_2^+`$ weakly couples to the nucleon ground-state since the non-relativistic description of the nucleon was quite successful in the naive quark model. On the other hand, nobody has succeeded in extracting the signal of the nucleon by using the so-called unconventional operator $`B_2^+`$ in lattice QCD $`^\mathrm{?}`$. This suggests that $`B_2^+`$ has negligible overlap with the nucleon ground-state and might give us a signal for an excited-state nucleon. Of course, this naive expectation is against the common sense that interpolating operators with the same quantum numbers should give the same mass spectrum. However, we find different plateaus in effective mass plots from $`B_1^+`$ and $`B_2^+`$ operators $`^\mathrm{?}`$. We apply the single-mass fit method described above to the two-point function composed of $`B_2^+`$. We find that masses extracted from $`B_1^+`$ and $`B_2^+`$ are quite different as shown in Figure 3. We conclude that we can identify $`B_2^+`$ with the “first” positive-parity excited-state ($`N^{}`$) of the nucleon for heavy quarks ($`m\stackrel{>}{}0.07`$$`^\mathrm{?}`$. The main reason is that we obtain mass estimates consistent with $`B_2^+`$ from the mixed correlation function $`0\|B_1^+\overline{B}_2^++B_2^+\overline{B}_1^+\|0`$. This suggests $`0\|B_2^+\|N0`$. Indeed, we see $`\|0\|B_2^+\|N/0\|B_2^+\|N^{}\|^210^3`$ from the two-mass fits for the $`B_2^+`$ correlation at $`m=0.10`$ and 0.125. This is plausible since the operator $`B_2^+`$ is expected to couple weakly to the ground state of the nucleon as we mentioned earlier $`^\mathrm{?}`$. This feature weakens in the lighter quark mass region (from around $`m=0.05`$). Unfortunately, we have no data of the mixed type correlation for $`m=0.02`$ and 0.03. To make a definite conclusion about the result of $`B_2^+`$, we need further calculations for the light quark mass. Finally, we want to touch upon the reason why we see a clear $`B_2^+`$ signal for the first time in this study while previous studies failed to do so. This should be related to mixing induced by explicit chiral symmetry breaking at finite lattice spacing which is absent in our calculation but severe in those calculations. Although $`B_1^+`$ and $`B_2^+`$ do not mix in the continuum because of different chiral structures<sup>3</sup><sup>3</sup>3To speak properly, $`B_1^+B_2^+`$ and $`B_1^++B_2^+`$ belong to two distinct $`(\frac{1}{2},0)(0,\frac{1}{2})`$ chiral multiplets under $`SU(2)_LSU(2)_R`$ $`^\mathrm{?}`$., it is known that unwanted mixing between them arises through the explicit breaking of chiral symmetry by conventional lattice fermions $`^\mathrm{?}`$. On the other hand, this breaking in DWF is expected to be suppressed exponentially with $`N_s`$. Thus DWF can drastically reduce the unwanted mixing $`^\mathrm{?}`$. Indeed what we see here suggests such reduction is significant. As a result, we are able to numerically confirm the expected feature $`0\|B_2^+\|N0`$ in the region of heavy valence-quark mass. ## 4 Conclusion We studied the mass spectrum of the nucleon excited-states in quenched lattice QCD using the domain-wall fermions (DWF) which preserves the chiral and flavor symmetries almost exactly. Most importantly we demonstrated that this method is capable of calculating the excited-state mass of $`N^{}`$. We made systematic investigation of the $`N^{}`$ spectrum by using two distinct interpolating operators, $`B_1^{}`$ and $`B_2^{}`$. We obtained mutually consistent results for the $`N^{}`$ mass spectrum from both of them. This is in contrast with the positive parity case as described below. In practice $`B_1^{}`$ is preferable to $`B_2^{}`$ in extracting the $`N^{}`$ mass signal. We found definite mass splitting between $`N^{}`$ and $`N`$ in the whole quark mass range studied. Furthermore, this splitting grows with decreasing quark mass. This is the first time such a remarkable feature has been observed in any theoretical calculation. The $`N^{}`$ mass and the $`N`$-$`N^{}`$ mass splitting in the chiral limit obtained by extrapolation is consistent with the experimental value within about 5-10%. This is very encouraging for further investigations of $`N^{}`$ physics. We observed no signal for the nucleon by using the unconventional nucleon operator $`B_2^+`$ which vanishes in the non-relativistic limit. Instead, we extracted the mass of the first excited nucleon $`N^{}`$ for heavy quarks. We need further study to make a definite conclusion about this subject for lighter quarks. ## Acknowledgments The author would like to thank Professor V. Burkert and the other organizers of NSTAR2000 for an invitation. He is also grateful to Professor L.Ya. Glozman and Professor D.O. Riska for recommending him to the organizers as an invited speaker in this workshop. We thank RIKEN, Brookhaven National Laboratory, and the U.S. Department of Energy for providing the facilities essential for the completion of this work. ## References
warning/0004/gr-qc0004026.html	ar5iv	text	# Discrete spacetime: classical causality, prediction, retrodiction and the mathematical arrow of timeTalk given at The First International Interdisciplinary Workshop on ”Studies on the Structure of Time: from Physics to Psycho(patho)logy”, 23-24 November 1999, Consiglio Nazionale delle Richerche, Area della Ricerca di Palermo, Palermo, Sicily, Italy. ## 1 Introduction The term causality is frequently used in a way which suggests that an intrinsic causal structure underlies the universe. In relativity this is reinforced by the assumption of a metric tensor with a Lorentzian signature. This gives the traditional light cone structure associated with spacelike and timelike intervals, and imposes conditions on the possible trajectories of particles and quantum field theory operator commutation relations. We shall discuss the idea that causality is a convenient account designed to satisfy and conform to the patterns of classical logic that the human Theorist wishes to believe underlies the dynamics of space, time and matter. In this approach causality need not be associated with any a priori concept of metric tensor. This view of causality has been suggested by various philosophers and scientists. Hume argued that causality is a fiction of the mind. He said that causal reasoning is due to the mind’s expectation that the future is like the past and because people have always associated an effect with a cause. Kant believed that causality is a category used to classify experience. Lentzen said that causality is a relation within the realm of conceptual objects. Lurchin said that causality is a personal way of thought that is not among our immediate sensual data, but is rather our basic way to organize that data. For Maxwell the principle of causality expresses the general objective of theoretical sciences to achieve deterministic explanations, and according to Heisenberg, causality only applies to mathematical representations of reality and not to individual mechanical systems. ## 2 The PPM view of time To avoid confusion we distinguish three sorts of time: $`i)`$ Process Time is the hypothesized time of physical reality. Although there is geological and astrophysical evidence for some sort of temporal ordering in reality , process time need not exist in any real sense and may just be a convenient way for humans to think about the Universe. $`ii)`$ Physiotime is the subjective time that humans sense and which they believe runs monotonically forwards for them. It is the end product of complex bio-dynamical processes occurring in process time. Its origins are not understood currently. Many physicists believe that this feeling is an illusion. What matters here is the undeniable existence of this feeling, because humans are driven by this sensation of an ever increasing time to believe that descriptions of reality must involve such a concept. $`iii)`$ Mathematical times are conceptual inventions of human theorists designed to model process time. Examples are Newtonian Absolute Time, relativistic coordinate times, proper time and cosmic time. Mathematical times usually involve subsets of the real line, which has an ordering property. This ordering is used to model the notions of earlier and later. This presupposes something about the nature of process time that may be unwarranted. In the Euclidean formulation of field theory for example there is no dynamical ordering parameter. ### 2.1 The Theorist We shall use the term Theorist to denote the human mind operating at its clearest and most rational in physiotime. The Theorist has the status of an observer or deity overseeing the mathematically consistent development of chosen mathematical models that are used to represent phenomena in process time. Free will enters into the discussion here as the freedom of the Theorist to choose boundary conditions in these models. Whether free will is an illusion or not is regarded here as irrelevant. ## 3 Classical causality The need to seek causal explanations stems from the peculiarities of human consciousness. Humans generally want to explain phenomena. When they do this they invariably try to invoke what may be called classical logic. This is the everyday logic that postulates that statements are either true or else false and that conclusions can be drawn from given premises. It is also the logic of vision, which generally informs the brain that an object either is in a place or else is not in that place. The rational conscious mind tends to believe that the external universe follows this logic, and this is the basis for the construction of CM (classical mechanics) and all the belief structures which it encodes into its view of reality. It is also the logic of jurisprudence and common sense. This logic served humanity extremely well for millennia, until technological advances in the early years of the twentieth century revealed that quantum phenomena did not obey this logic in detail. A CM Theorist is anyone who believes in a classical view of reality. In the mindset of a CM Theorist, reality is assumed to be strictly single valued at each and every time even in the absence of observation. Philosophers say that reality is determinate. The CM Theorist attempts to make unique predictions wherever possible, such as where a planet will be at a future time. The assumption is made that the planet will be somewhere at that time and not nowhere, and that it will not be in two or more places at that time. In general, quantum theory requires a pre-existing classical conceptual framework for a sensible interpretation. For example, relativistic quantum field theory assumes a classical Lorentzian metric over spacetime, and only the fields are quantised. For this reason, we shall focus our attention on a classical formulation of causality. ## 4 Functions and links To set up our framework for causality, it will be useful to review the definition of a function: #### Definition 1: A function $`f`$ is an ordered triple $`f(F,D,R)`$ where $`F,D`$ and $`R`$ are sets which satisfy the following: 1. $`F`$ is a subset of the Cartesian product $`D\times R`$; 2. for each element $`x`$ in $`D`$ there is exactly one element $`y`$ in $`R`$ such that the ordered pair $`(x,y)`$ is an element of $`F`$. Physicists tend to write $`y=f\left(x\right)`$ . $`D`$ is called the domain of (definition of) $`f`$ and $`R`$ is the range of $`f.`$ The image of $`f`$ is the subset $`f(D)`$ of $`R`$ such that for each element $`v`$ in $`f(D)`$ there is at least one element $`u`$ in $`D`$ such that $`v=f(u).`$ Without further information it cannot be assumed that $`f\left(D\right)=R,`$ but in our work this must be assumed to hold. Otherwise, there arises the possibility of having a CM where something could happen without a cause, i.e. an element $`z`$ of $`R`$ could exist for which there is no $`x`$ in $`D`$ such that $`z=f\left(x\right)`$. The range and domain of a function do not have a symmetrical relationship and the ordering of the sets in Definition $`1`$ is crucial. Usually, no pair of component sets in the definition of a function can be interchanged without changing the function. This asymmetry forms the basis of the time concept discussed in this article and defines what we call a mathematical arrow of time. Assuming that $`f`$ is single valued, then we may employ the language of dynamics here, though this may seem unusual. We could say that $`y=f(x)`$ is determined by $`x`$ via the process $`f`$, or that $`x`$ causes $`f\left(x\right)`$. Then $`x`$ is a cause, $`f\left(x\right)`$ is its effect and $`f`$ is the mechanism of causation. Although Definition $`1`$ carefully excludes the concept of a many-valued explicit function, such a possibility arises when we discuss implicit functions. #### Graphical notation: We shall use the following graphical notation: ### 4.1 Explicit functions: The process of mapping elements of $`D`$ into $`R`$ via $`f`$ will be represented by the LHS (left hand side) of $`Fig.1a`$, where the large circles denote domain and image sets, the small circle denotes the function, and arrows indicate the direction of the mapping. An alternative representation is given by the RHS (right hand side) of $`Fig.1a`$, where the labels $`D`$, $`R`$ and $`f`$ are understood. Here the integers $`0`$ and $`1`$ represent the ordering of the function from $`D`$ $`\left(0\right)`$ to $`R\left(1\right)`$, so that the arrows are not needed. ### 4.2 Many-to-one functions: Definition $`1`$ raises the question: which particular $`x`$ in $`D`$ caused a given $`y`$ in $`R\mathrm{?}`$ As given, nothing in Definition $`1`$ rules out many-to-one mappings, so without further information about the function $`f,`$ there may be more than one such $`x.`$ The assumption of a unique pre-image is equivalent to the belief that, given the present state of the universe, there was a unique past which gave rise to it. ’t Hooft has recently discussed this in the context of gravitation using equivalence classes of causes . This is bound up with the notion of irreversibility. If however it is the case that $`f`$ is one-to-one and onto, then its inverse function $`f^1`$ exists and so any element $`y=f\left(x\right)`$ in $`R`$ can be mapped back to a unique cause $`x`$ in $`D`$ via $`f^1.`$In such a case there would be no inherent difference in principle between the roles of $`D`$ and $`R`$. In the language of dynamics the mechanics would be reversible. ### 4.3 Implicit functions and links: Suppose now that the relationship between $`D`$ and $`R`$ is implicit rather than explicit. For example, let $`D`$ and $`R`$ each be a copy of the real line $`𝖱`$ with elements $`x`$,$`y`$ belonging to $`D`$ and $`R`$ respectively and suppose that the only dynamical information given was an implicit equation of the form $$g(x,y)=0.$$ (1) Then our graphical notation for this is given by $`Fig.1b`$, where the small circle now denotes an implicit equation or link $`g`$ relating elements of $`D`$ and $`R`$. Without further information no arrows are permitted at this stage. ### 4.4 Resolution of links: Given an implicit equation in two variables such as $`\left(\text{1}\right)`$, suppose now that it could proved that there was always a unique solution for $`y`$ in $`R`$ given any $`x`$ in $`D`$ (the solution $`y`$ of course depending on the value of $`x)`$. Then our convention is that now arrows pointing from $`D`$ into $`R`$ via the link $`g`$ may be added to indicate this possibility, giving the LHS of $`Fig.1a`$ with $`f`$ replaced by $`g`$. But now the relationship between $`D`$ and $`R`$ is formally equivalent in principle to having an explicit function $`y`$ of $`x`$ (even if $`y`$ could not be obtained analytically). In such a case we will also say that any given element of $`D`$ causes or determines a corresponding element of $`R`$ and that the link between $`D`$ and $`R`$ may be resolved from $`D`$ to (or in favour of) $`R`$. Our concept of resolution depends only on the existence of a unique solution and does not imply that a solution could actually be computed by the Theorist in practice. Computability is an attribute associated with physiotime, and is not here regarded as an essential ingredient of our version of causality. A more severe definition of causality however might be to impose the restriction of computability. We do not do this here because we wish to avoid anthropomorphism. What is important in classical mechanics is the existence of a unique resolution; the universe does not actually “compute” anything when this resolution occurs. ### 4.5 Inequalities: It is possible to consider links which are not equations but more general relations such as inequalities. For example, suppose $`D`$ and $`R`$ are copies of the real line with elements denoted by $`x`$, $`y`$ respectively and consider a link defined by $$g(x,y)x+y<0.$$ (2) Given $`x`$ there is an infinity of solutions for $`y`$, so in this case one possible interpretation would be that the link is equivalent to a many valued function of $`x`$, although this way of putting things might seem unusual. On the other hand, given a $`y`$ there is also an infinity of solutions for $`x`$. This leads to an alternative interpretation of the link as a many valued function of $`y`$. Such examples do not generate a classical TRA (temporal resolution of alternatives) and so would not occur normally in our classical spacetime dynamics. ### 4.6 Reversibility: Now suppose we were given an implicit equation for $`x`$ and $`y`$ such that we could prove the existence of a unique solution for either $`x`$ or $`y`$ given the other variable. Then the arrows could point either way and this would correspond to a choice of causation<sup>1</sup><sup>1</sup>1This choice is taken by the Theorist in physiotime.. Because cause and effect can be interchanged in such a case, it then becomes meaningful to talk about this dynamics being reversible. Clearly this is formally equivalent to having an invertible explicit function. When discussed in this way it becomes clear that in general, spacetime dynamics will be irreversible. Reversibility will occur only under very special conditions, which of course is the experience of experimentalists. ### 4.7 Generalization: In general the sets $`D`$ and $`R`$ need not be restricted to the reals. They could be any sort of sets, such as vector spaces, tensor product spaces, quaternions, operator rings, group manifolds, etc. Whatever their nature, they will always be referred to as events for convenience. In the sorts of dynamics we have in mind events will often be sets such as group manifolds. Links are defined as specific relationships between events. They may be more complicated than those in the above examples, and may have several or many components relating different components of different events. The specification of a set of events and corresponding links will be said to specify a (classical) spacetime dynamics. An important generalization is that a link might involve more than two events. Given for example five events $`R`$, $`S`$, $`T`$$`U`$ and $`V`$ with elements $`r,s,t,u`$ and $`v`$ respectively then a classical dynamics involving these events would give some set of equations or link $`g`$ of the generic form $$g(r,s,t,u,v)=0.$$ (3) This will be represented by $`Fig.2a`$. Suppose now, given such an $`g`$, that it could be proven that there is always a unique solution $`tT`$ given the other elements $`r,`$ $`s,`$ $`u`$ and $`v`$. In such a case this will be indicated by arrows pointing from $`R`$, $`S`$, $`U`$ and $`V`$ into $`g`$ and an arrow pointing into $`T`$ as in the LHS of $`Fig.2b`$. Then it will be said that $`g`$ can be (causally) resolved in (favour of) $`T`$, and $`T`$ will be called the resolved event. By definition, classical resolution is always in favour of a single resolved event, given initial data about the other events associated with the link. This does not imply anything about the possibility of resolving $`g`$ in favour of any of these other events. It may be or not be possible to do this. Suppose the Theorist could in principle resolve $`T`$ if they were given $`R,S,U`$ and $`V`$, and also resolve $`S`$ if they were given $`R,T,U`$ and $`V`$. Then the Theorist has to make a choice of resolution and choose one possible resolution and exclude the other. It would not be meaningful classically to resolve $`g`$ in favour of $`T`$ and $`S`$ at the same moment of physiotime, the reason being that these alternative resolutions employ different and inconsistent initial data sets (boundary conditions). Initial data sets are equivalent to information, and it is a self-evident premise that a Theorist can have at most one initial data set from a collection of possible and mutually inconsistent initial data sets at a given moment of physiotime. There remains one additional exotic possibility. If might be the case that given say $`R,S`$ and $`T`$, both $`U`$ and $`V`$ were implied by a knowledge of the link. For example, suppose the link was equivalent to the equation $$r+s+it+u+iv=0,$$ (4) where $`i=\sqrt{1}`$. Assuming that $`r,s,t,u`$ and $`v`$ were always required to be real, then we could always find $`u`$ and $`v`$ from a given $`r,s`$ and $`t,`$ simply by equating real and imaginary parts. A situation where a given link can be resolved in two or more events given just one initial data set at that link will be called a fluctuation process. Fluctuation processes will be excluded from our notion of classical causality. They may have a role in the QM (quantum mechanics) version of causality, which is beyond the scope of this article. One reason for excluding fluctuation processes is that this guarantees that information flows (in physiotime) from a link into a single event. This is related to the concept of cosmic time discussed below and to the idea that classical mathematical time has one dimension. When theorists discuss models with more than one parameter called a time, all but one of these has to be hidden or eliminated at the end of the day if a classical picture is to emerge. It may be argued therefore that the mechanism of classical resolution is the origin of the concept of time, and that time, like causality, is a no more than a convenient theoretical construct designed by the human mind to provide a coherent description of physical reality. This carries no implication that what we called process time is really a linear time. As we said before, process time is just a convenient label for something which may be quite different to what we believe it to be. When a choice of resolution exists and is made, then as an additional simplification and provided there is no confusion with other relations to which $`T`$ may be a party (not shown), the diagram on the LHS of $`Fig.2b`$ may be replaced by the RHS of $`Fig.2b`$. Here the numbers zero and one indicate the ordering of the resolution. Because $`T`$ may be regarded as caused by the other events, it can be regarded as later and so has a greater associated discrete time. Such times will be called dates. In general a link may be a whole collection of relations and equations. If there is just one small part of these equations which does not determine a unique solution fully, i.e., prevents a resolution, then arrows or dates are in principle not permitted. However, under some circumstances it may be reasonable to ignore some part of a dynamical relation in such a way that arrows could be justified as far as the remaining parts were concerned. For example, the microscopic laws of mechanics appear to be resolvable forwards and backwards in time (i.e., are reversible) provided the “small” matter of neutral kaon decays and the thermodynamic arrow of time (which could involve the gravitational field ) are ignored. In $`Fig.1a`$, $`D`$ is the complete cause of $`R`$; in $`Fig.2b`$, $`R`$ is a partial cause of $`T`$. The complete cause of $`T`$ is the collection of events $`R`$, $`S`$, $`U`$ and $`V`$, but only for this choice of resolution. The Theorist could decide to alter boundary conditions so that $`T`$ was no longer regarded as the resolved event. Having outlined our ideas on functions and links, we shall apply them now to discrete spacetime. ## 5 Discrete spacetime Fourier’s principle of similitude states that a system $`S^{}`$ similar to but smaller than another system $`S`$ should behave like $`S.`$ It is the physicist’s analogue of continuity in mathematics, and is of course an erroneous principle when applied to matter, as evidenced by the observation of atoms and molecules. It is generally supposed that this principle will also break down in the microscopic description of space and time. Classical GR (general relativity) may therefore be an approximation, albeit a remarkably good approximation, to some model of space and time which is not intrinsically a four-dimensional pseudo-Riemannian manifold. There have been numerous suggestions concerning the fundamental nature and meaning of space and time, such as twistor theory, point set theory, etc., and each of these suggestions makes a specific set of mathematical assumptions about spacetime. Likewise, in this article a specific view of space time and dynamics is proposed and its consequences explored. Of course, there is an important question concerning the use of classical or quantum physics here. In this paper the proposals are based on classical ideas and the ramifications of quantum physics are explored elsewhere. From before the time of Newton, physicists took the view that material objects have definite spatial positions at definite times. In the $`20^{th}`$Century theorists went further and developed to the extreme the Wellsian or block Universe view that space and time exist in some physically meaningful sense, even in the absence of matter. Whatever that sense is, be it a physical one or simply an approximate relationship of sorts between more complex attributes of reality, most physicists agree on the prime status of spacetime as the arena in which or over which physical objects exist. This is certainly the case for classical physicists and to various degrees for quantum physicists. In this paper the focus is on a classical description of a discrete spacetime structure. Discreteness is considered here for several reasons. First, as mentioned above, it would be too much to hope that Fourier’s principle of similitude should apply to spacetime and not to matter. Second, there is a strong feeling in the subject of quantum gravity that the Planck scales are significant. Third, discreteness has the advantage over continuity in being less mathematically restrictive. Theories based on discrete principles can usually encompass the properties of those based on continuous principles via appropriate limit processes, yet retain features which cannot occur in the continuum. Discrete spacetime structure and its relationship to causality has been discussed by a number of authors, notably by Sorkin et al and Finkelstein et al . A basic difference between those approaches and that taken here is that no a priori underlying spacetime manifold is assumed here. #### Proposition $`1`$: classical spacetime may be modelled by some discrete set $`𝒮`$. Elements of $`𝒮`$ will be denoted by capital letters such as $`P,Q,R,`$etc. $`𝒮`$ will be called a spacetime and its elements referred to as events even if subsequently $`𝒮`$ turns out to have only an indirect relationship with the usual four-dimensional spacetime of physics. What these events mean physically depends on the model. It is simplest to think of events as labels for mathematical structures representing the deeper physical reality associated with process time. Events are meaningful only in relationship with each other and it is meaningless to talk about a single event in spacetime without a discussion of how the Theorist relates it to other events in spacetime. This is done by specifying the links or relationships between the events. Links are as important as the events themselves and it is the totality of links and events which makes up our spacetime dynamics. This should include all attributes relating to matter and gravitation, and it is in principle not possible to discuss one without the other. The structure of our spacetime dynamics is really all there is; links and events. No preordained notion of metrical causality involving spacelike and timelike intervals is assumed from the outset. All of that should emerge as part of the implications of the theory. In classical continuous spacetime theories, on the other hand, the metric is usually assumed to exist independently of any matter, even before it is found via the equations of GR. As we said before, this metric carries with it lightcone structure and other pre-ordained attributes of causality. Discrete spacetime carries with it the astonishing possibility of providing a natural explanation for length, area and volume. According to this idea, attributed to Riemann , these are simply numerical counts of how many events lie in certain subsets of spacetime. Discreteness may also provide a natural scale for the elimination of the divergences of field theory, and permits all sorts of novelties to occur which are difficult if not impossible to build into a manifold. Recently, the study of spin networks in quantum gravity has revealed that quantization of length, area and volume can occur . ## 6 Event state space In CM the aim is usually to describe the temporal evolution of chosen dynamical degrees of freedom. These take on many possible forms, such as position coordinates or various fields variables such as scalar, vector and spinor fields. In our approach, we associate with each event $`P`$ in a given spacetime $`𝒮`$ an internal space $`\mathrm{\Omega }_P`$ of dynamical degrees of freedom called event state space. This space could be whatever the Theorist requires to model the situation. Moreover, it could be different in nature at each event $`P`$ in this spacetime. Elements of $`\mathrm{\Omega }_P`$ will be denoted by lower case Greek letters, such as $`\xi _p,\lambda _P`$ etc. and a chosen element $`\xi _P`$ of $`\mathrm{\Omega }_P`$ will be called a state of the event $`P`$. Elementary examples of event state spaces are: #### $`i)`$ Scalar fields: a real valued scalar field $`f`$ on a spacetime $`𝒮`$ is simply a rule which assigns at each event $`P𝒮`$ some real number $`f_p.`$ This may be readily generalized to complex valued functions. If no other structure is involved then obviously $`\mathrm{\Omega }_P`$ $`=`$ $`𝖱_PP𝒮,`$ $`\xi _P`$ $``$ $`f_P𝖱_P,`$ (5) where $`𝖱_P`$ is a copy at $`P`$ of the real line $`𝖱.`$ A classical configuration of spacetime in this model would then be some set $`\{f_P,f_Q,\mathrm{}\}`$. This configuration “exists” at some moment of the Theorist’s physiotime but the model itself would not necessarily have any causal ordering. That would have to be determined by the Theorist in the manner discussed below. #### $`ii)`$ Vector fields: suppose at each event $`P`$ there is a copy $`V_P`$ of some finite dimensional vector space $`V`$with elements $`𝐯V,`$ etc. Then a vector field is simply a rule which assigns at each event $`P`$ some element $`𝐯_PV_P.`$ #### $`iii)`$ Group manifolds: at each event P we choose a copy $`G_P`$ of some chosen abstract group, such as $`Z_2`$ or $`SU(3)`$ to be our event space. #### $`iv)`$ Spin networks: A spin network is a graph with edges labelled by representations of a Lie group and vertices labelled by intertwining operators. Spin networks were originally invented by Penrose in an attempt to formulate spacetime physics in terms of combinatorial techniques but they may also defined as graphs embedded in a pre-existing manifold . The state event space at each event $`P`$ is a copy of $`SU(2).`$ In this particular model the sort of events we are thinking of in spacetime would be associated with the geometrical links of a triangulation and the links (the dynamical relationship between our events) would be associated with the geometrical vertices of the triangulation, that is, with the intertwining operators. ### 6.1 Neighbourhoods and local environments For physically realistic models the number of events in the corresponding spacetime $`𝒮`$ will be vast, possibly infinite. Sorkin et al give a figure of the order $`10^{139}`$ per cubic centimetre-second, assuming Planck scales for the discrete spacetime structure. We shall find it useful to discuss some examples with a finite number of events for illustrative purposes. In our spacetime diagrams we will follow the convention established for functions in the previous section; large circles denote events and small circles denote links, with lines connecting events and links. Before any temporal resolution is attempted, no arrows can be drawn. $`Fig.3a`$ shows a finite spacetime with 14 events and 9 links. #### Definition $`2`$: The ( local) environment $`_A`$ of an event $`A`$ is the subset of links which involves $`A`$, that is, all those links to which $`A`$ is party, and the degree of an event is the number of elements in its local environment. For example, from $`Fig.3a`$, the local environment of the event labelled $`P`$ is the set of links $`_P\{f,g,h\}`$, and so $`P`$ is a third degree event. #### Definition $`3`$: The neighbourhood $`𝒩_A`$ of an event $`A`$ is the set of events linked to $`A`$ via its environment. For example, from $`Fig.3a`$ the neighbourhood of event $`P`$ is the set of events $`𝒩_P\{Q,R,S,T,U,V\},`$ and $`Q,R,S,T,U`$ and $`V`$ are the neighbours of $`P.`$ #### Definition $`4`$: The domain $`𝒟_f`$ of a link $`f`$ is the set of events involved in that link, and the order of a link is the number of elements in its domain. For example, from $`Fig.3a`$, $`𝒟_f`$ $`\{Q,R,P,U\}`$ and so $`f`$ is a fourth order link. The local environment of an event will be determined by the underlying dynamics of the spacetime, i.e. the assumed fundamental laws of physics. Currently these laws are still being formulated and discussed, so only a more general (and hence vague) discussion can be given here with some simple examples. A spacetime and its associated structure of neighbourhoods and local environments will be called a spacetime dynamics. ## 7 Kinds of spacetime dynamics There are two interrelated aspects of any spacetime dynamics: $`i)`$ the nature of the event state space associated with each event and $`ii)`$ the nature of the links connecting these events. The way in which events are related to each other structurally via the links will be called the local discrete topology. If this discrete topology has a regularity holding for all links and events, such as that of some regular lattice network, then we shall call this a homogeneous discrete topology. Otherwise it will be called inhomogeneous. We envisage three classes of classical spacetime dynamics: 1. Type $`𝐀`$: spacetime dynamics with a homogeneous and fixed discrete topology with a variable<sup>2</sup><sup>2</sup>2i.e., variable in physiotime. event state configuration which does not affect the discrete topology. This corresponds to (say) field theory over Minkowski spacetime; 2. Type $`𝐁`$: spacetime dynamics with inhomogeneous but fixed discrete topology with a variable event state configuration which does not affect the discrete topology. This corresponds to field theory over a fixed curved background spacetime, such as in the vicinity of a black hole; 3. Type $`𝐂`$: spacetime dynamics with discrete topology determined by the event state configuration; this corresponds to GR with matter. Type $`A`$ and $`B`$ spacetime dynamics are relatively easy to discuss. Once a fixed discrete topology is given this provides the template or matrix for the Theorist to “slot in” the causal patterns associated with initial event state configurations (initial data sets). In this sense, types $`A`$ and $`B`$ are not genuinely background free, but they are independent in a sense of any preordained Lorentzian metric structure. Type $`C`$ presents an altogether more interesting scenario to discuss and demonstrates the basic issue in classical GR which is that the spacetime dynamics should determine its own structure, including its topology. GR without matter of any sort does not make sense in our approach, because we need to specify event state spaces in order to define the links. Spins are needed to specify spin networks, for example. In our approach, gravitation is intimately bound up with discrete spacetime topology. A spacetime diagram need not be planar. Indeed, there need not be any concept of spacetime dimension at this stage. Sorkin et al suggest that at different scales, a given discrete spacetime structure might appear to be approximated by different continuous spacetime dimensions, such as $`26`$, $`10`$, or $`4,`$ depending on the scale. ## 8 Classical resolution and causal structure The PPM model was introduced as an approach to the modelling of process time. Working in physiotime, the Theorist first decides on a spacetime dynamics and from an initial data set then determines a consistent event state at each event in the spacetime. Because of the existence of the links, however, these event states cannot all be independent, and this interdependence induces our notion of causality, as we now explain. First, restrict the discussion to Type $`A`$ and $`B`$ spacetimes and suppose that each link is fully resolvable. By this is meant that if the order of a link is $`n`$ and the event states are specified for any $`n1`$ of the events in the domain of the link, then the remaining event space in the domain is uniquely resolved. #### Example $`1`$ An example of a fully resolvable link is the following: Let $`f`$ be a link of order $`r`$, with domain $`𝒟_f\{P_1,P_2,\mathrm{},P_r\}.`$Suppose the event space at $`P_i`$ is a copy $`G_i`$ of some group $`G`$ such as $`Z_2`$ or $`SU(n)`$ and suppose the link is defined by $$f:g_1g_2\mathrm{}g_r=e,$$ (6) where $`g_iG_i`$ and $`e`$ is the group identity<sup>3</sup><sup>3</sup>3We may use a matrix representation to define products of elements from different copies of $`G.`$. Then clearly, $$g_1=g_r^1g_{r1}^1\mathrm{}g_2^1$$ (7) and similarly for any of the other events states. In other words, we can always resolve any one of the events in $`𝒟_f`$ uniquely in terms of the others. Now suppose the Theorist chooses one link, such as $`f`$ in $`Fig.3a`$, which is a fourth order link. Then its domain $`𝒟_f`$ can be identified immediately from the spacetime dynamics to be $`\{P,Q,R,U\}`$. The Theorist is free to specify the states at three of these without any constraints, and this represents an initial data set called $`𝒮_0`$. Suppose these are events $`P,Q`$ and $`U`$. If now the structure of the link $`f`$ is such that there is only one possible solution in $`\mathrm{\Omega }\left(R\right)`$, then we have a classical resolution of alternatives<sup>4</sup><sup>4</sup>4i.e., the alternatives which form the event space $`\mathrm{\Omega }\left(R\right).`$ and the emergence of a causal structure. We can use the language of dynamics and say that events $`P,Q`$ and $`U`$ cause $`R`$ and denote it by a diagram such as $`Fig.2b`$. But it should be kept in mind that the Theorist has decided on which three sets to use as an initial data set. Our interpretation of causality is that it is dictated partly by the spacetime dynamics and partly by the choices made by the Theorist. In general, classical resolution involves using information about $`n1`$ event states at an $`n^{th}`$ order link to determine the state of the remaining event in the domain of that link. An initial data set may involve more than one link, such as shown in $`Fig.3b`$. In that diagram, events are labelled by integers representing the discrete times at which their states may be fixed. Events in the chosen initial data set are labelled with a time $`0`$ and shaded grey. These are the events $`\{O,P,Q,U\}.`$ Given an initial data set $`𝒮_0`$, the Theorist can then use the links to deduce the first (or primary) implication $`𝒮_1`$. This is the set labelled by integer time $`1`$ in $`Fig.3b`$, and consists of events $`\{R,V\}.`$ The event state at each of the events in the primary implication is determined from a knowledge of the event states on the initial data set, assuming the links do indeed permit a classical TRA. It could be the case that the primary implication is the empty set. Given a knowledge of the event states on $`𝒮_0`$ and its primary implication $`𝒮_1,`$ then the second (or secondary) implication $`𝒮_2\{W,T\}`$ can now be found and its events labelled by the discrete time $`2`$. This process is then repeated until the full implication $`𝒮_{\mathrm{}}`$ of $`𝒮_0`$ is determined. In $`Fig.3b`$ this is the set $`\{R,S,T,V,W,X,Y,Z\}.`$ In this example it is assumed that each of the links is fully resolvable. If any link is not fully resolvable, it may still be possible to construct a non-empty full implication for certain initial data sets. Several important concepts can be discussed with this example. ### 8.1 The future of an initial data set The full implication of an initial data set consists of those events whose event states follow from given initial conditions, and so it is not unreasonable to call the full implication of any initial data set the future of that set. Whilst the term “from” in the preceding sentence refers to a process of inference or implication carried out in physiotime by the Theorist, the result is that diagram $`Fig.3b`$ for example now carries dates (or equivalently arrows) as a consequence. This ordering may now be regarded as an attribute of the mathematical model rather than of physiotime, and this is the mathematical arrow of time referred to previously. The resulting structure can then be used to represent phenomena in process time. However, some caution should be taken here. We may encounter spacetime dynamics which are equivalent to reversible dynamics in continuous time mechanics. The full implication of some initial data sets for such spacetime mechanics may propagate both into what we would normally think of as the conventional future and into the conventional past. This is in accordance with the general problem encountered with any reversible dynamics: we can never be sure which direction is the real future and which is the real past unless external information is supplied to tell us. An example of a reversible discrete topology is given in $`Fig.4a`$. This is the topology of the discrete time harmonic oscillator . Assuming the links are fully resolvable, we see that the initial data set shown (shaded and labelled $`0`$) has a full implication which extends to the right and to the left of the diagram $`Fig.4b`$, i.e. into the conventional past and future. ### 8.2 Inaccessible events In $`Fig.3b`$ the set $`\{M,N\}`$ is inaccessible from the given initial data set $`\{O,P,Q,U\}`$, that is, its intersection with both the initial data set and its full implication is the empty set. Such an inaccessible set may be interpreted in a number of ways. It could be thought of as the absolute past of the initial data set $`\{O,P,Q,U\}`$ because in this particular example, the initial data set implies nothing about $`\{M,N\},`$ but specifying the states at $`M`$ and $`N`$ would fix $`Q`$. In other examples such inaccessible events could be interpreted as beyond an event horizon of some sort. The general feature of inaccessible events is that they cannot be affected by any changes to the event states in an initial data set. Archaeology is our term for the process of reconstructing portions of an absolute past from new and limited information added to some initial data set. It means the process whereby specifying one or more event states in an absolute past has the immediate consequence that the Theorist can deduce even more information about that past. #### Example $`2`$: Consider a spacetime dynamics based on the infinite regular triangular topology shown in $`Fig.5`$, with an initial data set shaded and labelled by $`0.`$ Its first, second and third implications are labelled by $`1,2`$ and $`3`$ respectively. The full implication of the initial data set in fact extends to infinity on the right. Suppose the Theorist now determines in some way or decides on the state at the event shaded and dated $`1`$ in $`Fig.5`$. Then the events in the past of the initial data set dated by $`1`$ and unshaded can now be resolved, thereby extending the original full implication one layer to the left of the initial data set. Just one extra piece of information can trigger an implication with an infinite number of elements. This process of retrodiction is called archaeology for obvious reasons<sup>5</sup><sup>5</sup>5The act of finding just a few Roman artefacts in a field may lead an Archaeologist to the conclusion that there had been a Roman settlement there., and for this spacetime could be continued indefinitely into the past. The Theorist can, by providing new initial data in this way, eventually cover the entire spacetime as the union of an extended initial data set and its full implication. ## 9 Causal Propagators and the speed of causality Suppose we have been given an initial data set $`𝒮_0`$ and have worked out its full implication $`𝒮_{\mathrm{}}`$. Now pick any event $`P`$ in $`𝒮_0`$ and change its state $`\psi _P`$. The consequence of this is to change the states in some events in $`𝒮_{\mathrm{}}`$, but not necessarily in all events in $`𝒮_{\mathrm{}}.`$The subset $`𝒫_P\left(𝒮_0\right)`$ of $`𝒮`$ consisting of all those events changed by the change in $`P`$ will be called the causal propagator associated with $`P`$ and $`𝒮_0`$. $`P`$ will be called the vertex of the propagator. We note the following: 1. for any event $`Q`$ in an initial data set $`𝒮_0,`$ the causal propagator $`𝒫_Q\left(𝒮_0\right)`$lies entirely within the full implication $`𝒮_{\mathrm{}}`$ of $`𝒮_0;`$ 2. A causal propagator depends on a vertex and on an associated initial data set; 3. A causal propagator divides spacetime into three sets: the vertex, those events which cannot be affected by any change at the vertex, and those events which could be changed. This structure is rather like the lightcone structure in special relativity which separates events into those which are timelike, lightlike, or spacelike relative to the vertex of the lightcone. In our context, we could in some sense talk about the speed of causality, analogous to the speed of light in relativity, as the limiting speed with which disturbances could propagate over our spacetime. #### Example $`3:`$ As an example we give a spacetime lattice of Type $`A`$ labelled by two integers $`m,n`$ running from $`\mathrm{}`$ to $`+\mathrm{}`$. The state space $`\mathrm{\Omega }_n^m`$ at each event $`P_n^m`$ is the set $`\{+1,1\}`$ and a state at $`P_n^m`$ will be denoted by $`\psi _n^m`$. The links are given by the equations $$\psi _n^m\psi _n^{m+1}\psi _n^{m1}\psi _{n1}^m\psi _{n+1}^m=1,\mathrm{}<m,n<\mathrm{}.$$ Now choose an initial data set $$\psi _0^m=\psi _1^m=+1,\mathrm{}<m<\mathrm{}.$$ This corresponds to selecting the index $`m`$ as a spatial coordinate and the index $`n`$ as a timelike coordinate. The initial data set is then equivalent to specifying the initial values and initial time derivatives of a scalar field on a hyperplane of simultaneity in a two dimensional spacetime. Because this spacetime dynamics is reversible, the full implication of this initial data set is the entire spacetime minus the initial data set. Now consider changing the state at $`m=n=0`$ from $`\psi _0^0=+1`$ to $`\psi _0^0=1.`$ In $`Fig.6`$ we show a bitmap plot of all those events whose state is changed by the change in the event $`(0,0)`$. The structure looks just like a lightcone, with complex fractal-like patterns developing inside the retarded and advanced parts of the lightcone. The speed of causality in this example is evidently unity if we interpret the indices in the manner discussed above. ## 10 Cosmic time For some spacetime dynamics a global temporal ordering can constructed by assigning an integer to each event as follows. If $`P`$ is earlier then $`Q`$ (i.e. P is a partial or complete cause of $`Q)`$ then some integer $`p`$ is assigned to $`P`$ and some integer $`q`$ to $`Q`$ such that $`p<q.`$These integers are called dates above. If it is possible to find a consistent ordering over the whole of $`𝒮`$ based on the above rule then we may say that a cosmic time exists for that spacetime. A cosmic time cannot be constructed for a finite spacetime dynamics if there are no events of degree $`1.`$ There are two situations where a cosmic time may be possible: either the spacetime is finite with one or more events of degree one, or the spacetime is infinite. However, these are not sufficient properties to guarantee a cosmic time can be constructed. It is possible to find spacetime dynamics for which more than one cosmic time pattern can be established. ### 10.1 Causal (timelike) loops: $`Fig.7`$ shows part of a spacetime dynamics containing a closed causal (timelike) loop. No cosmic time can be found for such a spacetime. This corresponds to the situation in GR where the existence of a closed timelike loop in a spacetime precludes the possibility of finding a global cosmic time coordinate for that spacetime. ### 10.2 Spacelike hypersurfaces In the spacetime depicted in $`Fig.3`$ the concept of metric was not introduced. Nevertheless, it is possible to give a definition of a spacelike (hyper) surface in this and other spacetime. Whether this corresponds to anything useful depends on the details of the spacetime dynamics. #### Definition 5: A spacelike hypersurface $``$ of a spacetime $`𝒮`$ is any subset of $`𝒮`$ which would have a consistent full implication if it were used as an initial data set. We have already used the term initial data set for such subsets. However, not all initial data sets need be consistent. A spacelike hypersurface is just a consistent initial data set. There may be many spacelike hypersurfaces associated with a given spacetime and they need not all be disjoint. The definition of spacelike hypersurface involves a choice of causation by the Theorist. In general there may be more than one spacelike hypersurface passing through a given event, and it is the Theorist’s choice which one to use. The possible non-uniqueness of spacelike hypersurface associated with a given event is desirable, because this is precisely what occurs in Minkowski spacetime, where there is an infinite number of spacelike hypersurfaces passing through any given event, corresponding to hyperplanes of simultaneity in different inertial frames. ## 11 Causal sets Causal sets are sets with some concept of ordering relationship, which makes them suitable for discussions concerning causality . This presupposes some pre-existing temporal structure independent of any dynamical input. This is not a feature of our dynamics, where causal structure emerges only after the Theorist has chosen the initial data set. ## 12 Concluding remarks and summary The PPM model leaves certain fundamental questions unanswered, such as the origin of physiotime and whether process time is a meaningful concept. However, once these are accepted as given and the model used in the right way, then causality and time itself emerge as observer (Theorist) oriented concepts. Many if not all of the phenomena associated with metric theories of spacetime can be recovered, which suggests that further investigation into this approach to spacetime dynamics may prove fruitful. A good question to ask is: where would Lorentzian causality come from in our approach? The answer is that it is embedded or encoded in the definition of the links. Only when full implications are worked out from given initial data sets would it be noticed that the dynamics itself naturally forces certain patterns to emerge and not others. Only at that stage would the Theorist would recognize an underlying bias in the dynamics in favour of certain more familiar interpretations. For example, given the continuous time equation $$\left(\frac{^2}{t^2}\frac{^2}{x^2}m^2\right)\phi (t,x)=0,$$ (8) the notation suggests that $`t`$is a time, so we could attempt to solve it with initial data on the hyperplane $`t=0.`$ However, it would soon emerge that evolution in $`t`$ gave runaway solutions. In other words, the equation itself would carry the information that it would be wiser to define initial data on the hyperplane $`x=0`$. We would not need to invoke the spurious concept of Lorentzian signature metric embedded in spacetime to discover this. At this stage we might suddenly realise that we had by some chance interchanged the symbols for time and space in an otherwise ordinary Klein-Gordon equation with a real mass. So rather than deal with $`\left(\text{8}\right),`$ which behaves like a Klein-Gordon equation with an imaginary mass, we would simply interchange the symbols $`t`$ and $`x`$ and then define initial data on the hyperplane now defined by $`t=0`$. ## 13 Acknowledgment The above is an expanded version of the author’s talk at The First International Interdisciplinary Workshop on “Studies on the Structure of Time: from Physics to Psycho(patho)logy”, 23-24 November 1999, Consiglio Nazionale delle Richerche, Area della Ricerca di Palermo, Palermo, Sicily, Italy. The proceedings of this conference, including this article, will be published presently by Kluwer Academic (New York). The author is grateful to the University of Nottingham for support, to the Organisers of the Conference for their assistance, to the other participants for their thoughts, and to Kluwer Academic (New York) for permission to place this article in these electronic archives.
warning/0004/cond-mat0004372.html	ar5iv	text	# First-principles study of ferroelectric and antiferrodistortive instabilities in tetragonal SrTiO3 ## I Introduction First-principles calculations are proving to be one of the most powerful tools for carrying out theoretical studies of the electronic and structural properties of materials. A particularly successful application of this technique has been its use in understanding the perovskite ferroelectric compounds. These materials have important technological applications because of their switchable macroscopic polarization and their piezoelectric properties. They are also attractive objects of fundamental study because of the rich variety of phase diagrams that they display as a function of temperature. At high temperature, the ABO<sub>3</sub> perovskites retain full cubic symmetry. However various structural phase transitions take place as the temperature is reduced. For example, BaTiO<sub>3</sub> and KNbO<sub>3</sub> undergo phase transitions from the cubic paraelectric (PE) phase to a succession of tetragonal, orthorhombic, and finally rhombohedral ferroelectric (FE) phases. In contrast, PbTiO<sub>3</sub> displays only a single transition, from the cubic PE phase to a tetragonal FE phase. In NaNbO<sub>3</sub> and PbZrO<sub>3</sub>, non-polar antiferrodistortive (AFD) or antiferroelectric (AFE) transitions take place, associated with different types of tilts of the oxygen octahedra, in addition to the FE transitions. It is understood that the two types of transitions result from the condensation of soft phonon modes at the Brillouin zone boundary with $`q0`$ and at the zone center with $`q=0`$ (Ref. ). The low-temperature behavior of SrTiO<sub>3</sub> has been an attractive subject for experimental and theoretical study. SrTiO<sub>3</sub> behaves as an incipient ferroelectric (similar to KTaO<sub>3</sub>) in the sense that it has a very large static dielectric response and is only barely stabilized against the condensation of the FE soft mode at low temperature. As the temperature is reduced, SrTiO<sub>3</sub> first undergoes a transition from the cubic to a tetragonal AFD phase at 105K, but this transition is of non-polar character and has little influence on the dielectric properties. The static dielectric response closely obeys a Curie-Weiss law of the form of $`ϵ(TT_c)^1`$ at temperatures above about 50K, but the divergence at a critical temperature $`T_c`$ 36K that would be expected from this formula is not observed. Instead, the susceptibility saturates at an enormous value of $`2\times 10^4`$ as $`T`$ approaches zero. Because the system is so close to a ferroelectric state, it is not surprising to find that it can be induced to become ferroelectric, either by the application of electric field, uniaxial stress or by the substitution of Ca ions on the Sr sublattice. Finally, the SrTiO<sub>3</sub> system also displays puzzling phonon anomalies and electrostrictive response in the low-temperature regime. This peculiar behavior, especially the failure of the system to condense into a FE phase at low $`T`$, has been the subject of considerable theoretical study and speculation. Recently, efforts have focused on the so-called “quantum paraelectric state” postulated by Müller and Burkard, who suggested that quantum fluctuations of the atomic positions could suppress the FE transition and lead to a stabilized paraelectric state. This hypothesis has received dramatic support from a recent experiment showing that isotopically exchanged SrTi<sup>18</sup>O<sub>3</sub> appears to become ferroelectric at 23K, suggesting that normal SrTi<sup>16</sup>O<sub>3</sub> must be very close indeed to the ferroelectric threshold. First-principles calculations have already contributed significantly to the understanding of the structural properties of SrTiO<sub>3</sub>. Calculations of two groups confirmed that SrTiO<sub>3</sub>, in its high-symmetry cubic structure at $`T=0`$, is unstable to both FE and AFD distortions when the atomic coordinates are treated classically. Using classical Monte Carlo simulations on an effective-Hamiltonian fitted to the first-principles calculations, Zhong and Vanderbilt predicted that SrTiO<sub>3</sub> would first undergo the AFD transition at about 130 K, and then a further transition into a state with simultaneous AFD and FE character at 70 K. Anharmonic interactions between the AFD and FE modes were found to be competitive, in the sense that the presence of the AFD distortion was found to reduce the FE transition temperature by about 20%. The same authors later showed that, when a quantum-mechanical treatment of the atomic positions was included via a quantum path-integral Monte Carlo simulation, the AFD transition temperature was shifted very close to the experimental one at 105 K, and the FE transition was suppressed down to the lowest temperatures that could be studied ($``$5 K), consistent with the experimental absence of a transition. On the other hand, LaSota and coworkers have recently performed a first-principles calculation of the ground state structural properties and interactions between the FE and AFD instabilities in SrTiO<sub>3</sub> using an LAPW approach. These authors found that the AFD tetragonal structure is stable against the FE distortions, indicating no interaction between the AFD and FE modes. The conclusion is in contrast to the previous theory. However, LaSota et al. made certain approximations to the FE eigenmodes, as will be discussed below. An additional motivation for a detailed theoretical study of SrTiO<sub>3</sub> is the opportunity to make contact with the remarkably systematic experimental study of Uwe and Sakudo. These authors made careful measurements of the anisotropic dielectric susceptibilities and Raman mode frequencies as a function of uniaxial stress applied along different crystal orientations. They also fitted their results, plus those of previous experimental studies, to obtain a phenomenological description of the couplings between the AFD, FE, and strain degrees of freedom of the crystal. In particular, their fit contains an anisotropic coupling which, in the tetragonal AFD phase, tends to favor FE distortions that are perpendicular to the tetragonal axis over those that are parallel. However, previous theoretical work has given an unclear picture of this anisotropy. On the one hand, Vanderbilt and Zhong found that the interaction between the FE and AFD modes, which was mainly through the on-site anharmonic coupling, would tend to favor FE modes polarized perpendicular to the AFD tetragonal axis, in accord with experiment. On the other hand, the Monte Carlo calculations previously referred to indicated a sequence of transitions with decreasing temperature in which the FE order parameter of a $`z`$-polarized mode was found to develop before the $`x`$ or $`y`$-polarized ones, indicating that the $`z`$-polarized mode goes soft first. Finally, the recent work of LaSota et al. suggests that there is very little anisotropic coupling at all. In view of these apparently conflicting theoretical results, we felt it worthwhile to clarify this situation by carefully studying the anisotropy of the FE distortion energy in the AFD ground state. With these motivations, we have carried out a thorough analysis of the ground-state structural and dynamical properties of tetragonal SrTiO<sub>3</sub> using first-principles density-functional calculations. This approach is based on a classical treatment of the nuclear motions, and so is obviously unable to take into account the quantum zero-point motion of the ionic positions which becomes critically important at low temperature. Nevertheless, from such a calculation one is still able to compute interaction parameters for comparison with experiment, to identify the effects which tend to suppress the FE instability in the presence of the AFD state, and to obtain qualitatively a picture of the dielectric anisotropy connected with the splitting of the differently polarized FE modes in the AFD state. The calculations are carried out using a plane-wave basis and ultrasoft pseudopotentials. The theoretical equilibrium AFD structure is obtained by minimizing the energy with respect to cell volume, $`c/a`$ ratio, and internal parameters. Frozen-phonon calculations are then used to obtain the frequencies of $`\mathrm{\Gamma }`$-point and $`R`$-point phonon modes, including FE soft modes, in the AFD ground state. For this purpose, we make use of a point-group symmetry analysis to reduce the complexity of the distortions that need to be studied. In order to interpret the results in terms of a phenomenological description involving FE and AFD mode distortions and strains, we use an approach similar to that underlying the effective-Hamiltonian scheme first developed for the BaTiO<sub>3</sub> system and later applied to SrTiO<sub>3</sub>. That is, we use the LDA calculations to compute the values of the Taylor expansion coefficients of the total energy with respect to these distortions, and compare with the experimental determinations of Uwe and Sakudo. The rest of the manuscript is organized as follows. In Sec. II we briefly describe the technique employed for the first-principles calculations. In Sec. III we present and discuss the results of the calculations. We begin with the determination of the theoretical tetragonal AFD structure, and then proceed to study the energies of AFD and FE distortions about this reference structure, with special attention to the anisotropies of the AFD and FE mode frequencies. Finally, we conclude in Sec. IV. ## II Theoretical details Our ab-initio plane-wave pseudopotential calculations are based on the Hohenberg-Kohn-Sham density-functional theory (DFT) within the local-density approximation (LDA). Ultrasoft Vanderbilt pseudopotentials are used, with the O($`2s`$), O($`2p`$), Ti($`3s`$), Ti($`3p`$), Ti($`3d`$), Ti($`4s`$), Sr($`4s`$), Sr($`4p`$), and Sr($`5s`$) states included in the valence. The exchange-correlation energy is of the Ceperley-Alder form with Perdew-Zunger parameterization. A conjugate-gradient minimization scheme is used to minimize the Kohn-Sham energy, using a plane-wave cutoff of 30 Ry for all calculations. Unless otherwise stated, our calculations are carried out at the theoretical equilibrium lattice constant of 7.303 a.u. which is $``$1% less than the experimental value of 7.365 a.u., the discrepancy representing the inherent LDA error. Cubic SrTiO<sub>3</sub> has a simple cubic 5-atom unit cell with a common lattice parameter $`a`$ along the , , and directions. We will briefly discuss some calculations carried out for a doubled unit cell corresponding to the condensation of a soft AFD mode at the (110)$`\pi /a`$ or $`M`$ point of the Brillouin zone (BZ) boundary. However, most of our attention will be focused on the ground-state tetragonal phase obtained by freezing in an AFD phonon mode at the (111)$`\pi /a`$ or $`R`$ point of the BZ boundary. This triply-degenerate phonon mode corresponds to the rotation of the TiO<sub>6</sub> octahedra in opposite directions from one cubic unit cell to the next, followed by a small tetragonal strain. (Note that it is conventional to label the phonon modes with respect to the simple-cubic BZ, even when they condense to lower the symmetry.) Taking the rotation to be about the $`z`$ axis, we adopt a 10-atom tetragonal unit cell with lattice vectors of length $`\sqrt{2}a`$, $`\sqrt{2}a`$, and $`c`$ along the , \[$`\overline{1}`$10\], and directions, respectively. (That is, in our convention, $`c/a`$ is close to 1, not $`1/\sqrt{2}`$.) The rotation of the oxygen atoms in the Ti$``$O plane is shown in Fig. 1. Throughout this paper, we will use $`x^{}`$ and $`y^{}`$ to denote the original cubic directions ( and , respectively), while $`x`$ and $`y`$ are taken as parallel to the tetragonal lattice vectors along and \[$`\overline{1}`$10\], respectively. That is, the $`x`$-$`y`$ frame is rotated by 45 relative to the $`x^{}`$-$`y^{}`$ frame ($`z`$ axes are congruent). In all cases, we use a k-point set that is equivalent to the $`6\times 6\times 6`$ Monkhorst-pack mesh in the BZ of the simple cubic cell, corresponding to 108 k-points in the full BZ of the tetragonal cell. The irreducible BZ then contains 6 k-points for the undistorted cubic structure; 10 k-points for the tetragonal ground-state structure, with or without additional $`A_{1g}`$ or $`A_{2u}`$ mode displacements; and 20 k-points for the tetragonal structure with additional $`E_u`$ mode displacement. ## III Results and Discussions ### A AFD instability in cubic unit cell To establish notation, we let the energy of an AFD phonon mode per 10-atom cell of the cubic perovskite structure be expanded up to fourth order in $`\varphi _z`$, $$E=E_0+\frac{1}{2}\kappa \varphi _z^2+A_x\varphi _z^4,$$ (1) where $`\varphi _z=(a/2)\mathrm{sin}\theta _z`$ is the magnitude of the oxygen-atom displacement associated with the rotation of the oxygen octahedra around the axis, as shown in Fig. 1. In Fig. 2 we show the computed values of the total energy versus rotation angle for AFD modes at both the $`M`$ and $`R`$-points of the BZ (corresponding to in-phase or out-of-phase rotations in neighboring planes of octahedra along $`z`$ respectively). These were computed at zero strain, i.e., with the lattice vectors fixed to be those of the theoretical equilibrium cubic structure ($`a`$=7.303 a.u.). As can be seen, the computed total energy versus rotation angle can be fitted very well by the quartic Eq. (1). Defining the mode stiffness $`\kappa `$=$`^2E/\varphi _z^2`$, we find that $`\kappa <0`$ for both $`M`$\- and $`R`$-point modes as shown in Fig. 2. Nevertheless, the magnitude of $`\kappa `$ for the $`M`$-point mode is only $``$10% that of the $`R`$-point mode, indicating that the instability at the $`R`$ point is much stronger than that at the $`M`$ point. Consequently, for the remainder of this paper we will limit our discussion to $`R`$-point distortions only. As can be seen from Figs. 2 and 3, the equilibrium octahedral rotation angle is found to be $`\theta _z`$=5.5, significantly larger than the zero-temperature experimental value of 2.1. Since the theoretical equilibrium lattice constant (7.303 a.u.) is somewhat smaller than the experimental one, we also carried out similar total-energy calculations at the extrapolated zero-temperature (7.365 a.u.) and room-temperature (7.38 a.u.) experimental lattice constants. The results shown in Figure 3 confirm that increasing the lattice constant or crystal volume tends to suppress the AFD instability, as expected from previous work. However, the resulting variation of the equilibrium rotational angle is too small to explain the experimental observation, changing only marginally to 4.89 and 4.69 at the zero- and room-temperature experimental lattice constants, respectively. These results demonstrate that the underestimate by $`1\%`$ of the lattice constant by the LDA is not the primary factor responsible for the theoretical overestimate of the rotation angle. Moreover, we shall see in the next subsection that the inclusion of strain relaxation effects only acts to increase (slightly) the theoretical equilibrium rotation angle. Thus, we think that the smaller observed value of the AFD rotation angle can most likely be attributed to the quantum fluctuations associated with the motion of the oxygen atoms. This effect is not included in the theory, and should act to reduce the amplitude of symmetry-breaking distortions. While previous work has indicated that the quantum fluctuations should have a weaker effect on the AFD modes than upon the FE ones, the effect on the AFD modes could still be quite significant. An alternate possibility is simply that the underestimate is a result of LDA error not associated with the lattice constant. In any case, we have chosen to complete our theoretical investigations by considering distortions about our theoretical ground-state AFD structure, keeping in mind that the results should be interpreted with the overestimate of the rotation angle in mind. ### B AFD modes in the tetragonal structure To study the ground-state tetragonal structure, the lattice strains also need to be taken into account. We adopt the usual Voigt notation $`x_i`$ for the strain tensor, but set $`x_4=x_5=x_6=0`$ because such off-diagonal shear strains will not enter into our considerations. We take $`i`$=1, 2, 3 corresponding to the $`x^{}`$, $`y^{}`$, and $`z`$ pseudocubic axes for both strains $`x_i`$ and rotations $`\varphi _i`$. Expanding the energy up to quartic order in AFD amplitudes, quadratic order in strain, and leading order in the strain-AFD coupling, the symmetry-allowed contributions are $`E=E_0`$ $`+`$ $`{\displaystyle \frac{1}{2}}\kappa {\displaystyle \underset{i}{}}\varphi _i^2+A_x{\displaystyle \underset{i}{}}\varphi _i^4+A_x^n{\displaystyle \underset{i<j}{}}\varphi _i^2\varphi _j^2`$ (2) $`+`$ $`{\displaystyle \frac{1}{2}}c_{11}{\displaystyle \underset{i}{}}x_i^2+c_{12}{\displaystyle \underset{i<j}{}}x_ix_j`$ (3) $``$ $`b_{11}{\displaystyle \underset{i}{}}x_i\varphi _i^2b_{12}{\displaystyle \underset{i<j}{}}x_i\varphi _j^2).`$ (4) We choose the tetragonal ground state to be oriented along the $`z`$-axis, with $`x_1=x_20`$, $`\varphi _1=\varphi _2=0`$, and $`\varphi _30`$. In this ground state, phonon modes corresponding to additional oscillations of $`\varphi _i`$ belong either to the $`E_g`$ ($`i`$=1 or 2) or the $`A_{1g}`$ ($`i`$=3) representation of the tetragonal $`D_{4h}`$ point group. For symmetry-preserving ($`A_{1g}`$) distortions, it is convenient to re-express the three strain components in terms of a volume strain $`\overline{x}`$ and a shear strain $`v`$ according to $`\begin{array}{cc}x_1=x_2=\overline{x}v,\hfill & \\ & \\ x_3=\overline{x}+2v.\hfill & \end{array}`$ (8) In terms of these variables, the second and third lines of Eq. (4) can then be rewritten as $$E^{\mathrm{elastic}}(\overline{x},v)=\frac{3}{2}\alpha \overline{x}^2+3\beta v^2$$ (9) and $$E^{\mathrm{coupling}}(\overline{x},v,\varphi _i)=\eta \overline{x}\underset{i}{}\varphi _i^2\gamma v(2\varphi _3^2\varphi _1^2\varphi _2^2),$$ (10) where $$\begin{array}{cccc}\alpha =c_{11}+2c_{12},\hfill & & & \\ & & & \\ \beta =c_{11}c_{12},\hfill & & & \\ & & & \\ \eta =b_{11}+2b_{12},\hfill & & & \\ & & & \\ \gamma =b_{11}b_{12}.\hfill & & & \end{array}$$ (11) To find the equilibrium rotation angle in the $`A_{1g}`$ ground state, we hold $`\varphi _1`$=$`\varphi _2`$=0 and minimize $`E`$ in Eq. (4) with respect to $`\overline{x}`$ and $`v`$ at a fixed $`\varphi _3=\varphi _z`$. The minimizing values are $$\begin{array}{cc}\overline{x}^{\mathrm{eq}}=\frac{\eta }{3\alpha }\varphi _z^2,\hfill & \\ & \\ v^{\mathrm{eq}}=\frac{\gamma }{3\beta }\varphi _z^2.\hfill & \end{array}$$ (12) Substituting into Eq. (4), $$E(\varphi _z)=E_0+\frac{1}{2}\kappa \varphi _z^2+A_X\varphi _z^4,$$ (13) where $$A_X=A_x\frac{\eta ^2}{6\alpha }\frac{\gamma ^2}{3\beta }.$$ (14) Thus, when the strain relaxation is taken into account, the equilibrium rotation angle is given by $$\varphi _z^{\mathrm{eq}}=\sqrt{\frac{\kappa }{4A_X}}.$$ (15) We determine all the interaction parameters $`\kappa `$, $`A_x`$, $`A_x^n`$, $`c_{11}`$, $`c_{12}`$, $`b_{11}`$, and $`b_{12}`$ via a series of finite-difference calculations of total energies and forces within the LDA. Table I lists our results and compares them with the corresponding values determined by Uwe and Sakudo by fitting to experiment (all units have been converted to atomic units). We find very good agreement overall. The fact that $`A_X<A_x`$ implies that the inclusion of strain relaxations strengthens the AFD instabilities at anharmonic order. The equilibrium rotation angle increases to 6.0 (compared with 5.5 for the cubic strain state), while the equilibrium values of $`\overline{x}`$ and $`v`$ are found to be $`0.10\%`$ and $`0.23\%`$, respectively. In the tetragonal AFD ground state, the frequencies of the soft phonon modes associated with additional rotations of the oxygen octahedra are given by evaluating $$m_\varphi \omega _i^2=\frac{^2E}{\varphi _i^2}\|_{\mathrm{eq}},$$ (16) where $`m_\varphi =4m_\mathrm{O}`$ is the mass factor associated with the oxygen rotational mode and the derivative is to be evaluated under conditions of fixed strain at the equilibrium structure (i.e., at the equilibrium values of $`\overline{x}`$, $`v`$, and $`\varphi _z`$). Eq. 16 gives the frequencies of the $`E_g`$ and $`A_{1g}`$ modes to be, respectively, $$\begin{array}{cc}\omega _1^2=\omega _2^2=\kappa (A_X^n/2m_\varphi A_X),\hfill & \\ & \\ \omega _3^2=2\kappa A_x/m_\varphi A_X,\hfill & \end{array}$$ (17) where $`A_X^n=A_x^n+\gamma ^2/\beta `$. The values of the two frequencies are 45 cm<sup>-1</sup> and 130.7 cm<sup>-1</sup> respectively, so that the frequencies of the $`A_{1g}`$ and $`E_g`$ modes are roughly in the ratio 3:1. This is consistent with observed ratios of $``$2.5:1 in pressure-dependent experiments and $``$3:1 in temperature-dependent experiments . Up to this point, the analysis has been done at the theoretical equilibrium lattice constant. However, it is well known that the LDA tends to underestimate the lattice constants of perovskites by $``$1%. Moreover, past experience has shown that the displacement patterns associated with soft modes may depend critically on the lattice constant and strains. To take these effects into account, we adopted a strategy of applying a negative hydrostatic pressure to the lattice to restore the experimental lattice constant. Using Eq. (4) and minimizing the Gibbs free energy $$G=E+3\overline{x}P$$ (18) with respect to $`\overline{x}`$ at fixed pressure $`P`$, we find $$\overline{x}=\frac{\eta }{3\alpha }\varphi _z^2\frac{P}{\alpha },$$ (19) while $`v`$ is independent of pressure. Then $$G(\varphi _z,P)=\frac{1}{2}\kappa _{\mathrm{eff}}\varphi _z^2+A_X\varphi _z^4\frac{3P^2}{2\alpha }$$ (20) where the effective harmonic coefficient is $`\kappa _{\mathrm{eff}}=\kappa +2\eta {\displaystyle \frac{P}{\alpha }}.`$ Thus, at fixed pressure, the equilibrium rotation angle is $$\varphi _z^{\mathrm{𝑒𝑞}}=\sqrt{\frac{\kappa _{\mathrm{eff}}}{4A_X}}.$$ (21) As one can see, the harmonic coefficient $`\kappa _{\mathrm{eff}}`$ depends upon the external pressure variable. It would thus be possible, in principle, to adjust $`P`$ so as to fit the resulting rotation angle to the experimental angle of 2.1. However, the pressure needed to achieve this, $``$14.4GPa, would expand the lattice constant to 7.48 a.u., which is much larger than the experimental value. Instead, we adjust $`P`$ so as to fit the experimental lattice constant. That is, we adjust $`P`$ so that the volume strain is $`\overline{x}=(a_{\mathrm{exp}}a_{\mathrm{theo}})/a_{\mathrm{theo}}=0.849\%`$, where $`a_{\mathrm{exp}}`$=7.365 a.u. and $`a_{\mathrm{theo}}`$ are the zero-temperature experimental and theoretical lattice constants, respectively. Substituting (21) into (19), we obtain $$\overline{x}(P)=x_0\frac{P}{\alpha _{\mathrm{eff}}}$$ (22) where $$\begin{array}{cc}x_0=\frac{\kappa \eta }{12\alpha A_X},\hfill & \\ & \\ \alpha _{\mathrm{eff}}=\alpha \left[1+\frac{\eta ^2}{6\alpha A_X}\right]^1.\hfill & \end{array}$$ (23) Inserting $`\overline{x}=0.849\%`$ leads to $`P=`$5.26GPa. The strains along the tetragonal and planar axes are found to be 1.135% and 0.705%, respectively. Relative to a cubic cell with the experimental lattice parameter 7.365 a.u., the tetragonal cell we obtained is thus expanded along while compressed along the and directions. In this circumstance, the equilibrium rotation angle is found to be 4.93. Under these conditions, the $`A_{1g}`$ and $`E_g`$ soft-mode frequencies of Eq. (17) now become $$\begin{array}{cc}\omega _1^2=\omega _2^2=\kappa _{\mathrm{eff}}(A_X^n/2m_\varphi A_X),\hfill & \\ & \\ \omega _3^2=2\kappa _{\mathrm{eff}}A_x/m_\varphi A_X.\hfill & \end{array}$$ (24) From the coefficients in Table I, $`A_x/A_X`$=$`1.187`$ and $`2A_X^n/A_X`$=$`1.11`$, so that clearly $`\|\omega (E_g)\|<\|\omega (A_{1g})\|`$. In fact, the two frequencies are calculated to be $`\omega (E_g)`$=37 cm<sup>-1</sup> and $`\omega (A_1g)`$=109 cm<sup>-1</sup>. (For comparison, $`\omega (A_1g)`$=124 cm<sup>-1</sup> in Ref. , while the measured $`\omega (E_g)`$ and $`\omega (A_1g)`$ in Ref. are 15 cm<sup>-1</sup> and 48 cm<sup>-1</sup> respectively). To summarize the results so far, we have found that the AFD mode condenses at the $`R`$ point of the cubic BZ, associated with a triply-degenerate phonon of $`\mathrm{\Gamma }_{25}`$ symmetry. As a consequence of the transition from the cubic to the tetragonal state, the degenerate $`R`$-point modes split strongly into an $`A_{1g}`$ singlet and an $`E_g`$ doublet, the latter having a softer frequency than the former. This is in good qualitative agreement with experiment. The expansion approach used above for the AFD modes makes an implicit assumption that the phonon eigenvectors from the cubic structure are a good approximation to those in the tetragonal AFD structure. For the AFD modes, where the anisotropy is large, we do not expect this approximation to be at all serious. However, our next task will be to analyze the FE mode anisotropy in the AFD state. As will be seen below, this turns out to be much more delicate than for the AFD modes. Thus, we have chosen to take a different approach for the FE modes, in which the normal modes in the AFD ground state are directly computed. The symmetry analysis needed to do this is given in the next subsection, and the FE mode analysis is then given in the concluding subsections. ### C Symmetry analysis of normal modes in AFD tetragonal phase In this section, we present some details of the point-group symmetry analysis of the normal modes in the AFD tetragonal structure, needed for the calculation of the frequencies of transverse optical (FE) modes at the Brillouin zone center. To harmonic order, the displacement energy can be expressed as $$E=\frac{1}{2}\underset{i,j,\alpha \beta }{}\mathrm{\Phi }_{i,j}^{\alpha ,\beta }u_i^\alpha u_j^\beta ,$$ (25) where $`u_i^\alpha `$ is the displacement of sublattice $`i`$ in Cartesian direction $`\alpha `$, and the force constant matrix $`\mathrm{\Phi }`$ obeys the symmetry conditions $`\mathrm{\Phi }_{i,j}^{\alpha ,\beta }=\mathrm{\Phi }_{j,i}^{\beta ,\alpha }`$ and $`_j\mathrm{\Phi }_{ij}=0`$. The dynamical matrix is just related to the force constant matrix by a diagonal mass tensor. It is well known that the vibrational modes at a given $`𝐤`$-point in the BZ of a crystal transform according to the corresponding irreducible representations of the symmetry group for that $`𝐤`$-point. Such an analysis, which has previously been used to construct the force-constant matrices for the FE modes in the cubic perovskite structure and for all modes in the tetragonal FE structure of PbTiO<sub>3</sub>, is applied here to the zone-center modes in the AFD tetragonal structure of SrTiO<sub>3</sub>. When SrTiO<sub>3</sub> condenses from the cubic into the tetragonal AFD phase, the point group lowers from $`O_h`$ to $`D_{4h}`$. Because we are interested in zone-center modes, the symmetry group of $`𝐤`$ is just the $`D_{4h}`$ point group itself, which contains 16 symmetry operations that can all be generated from a fourfold rotation $`C_4`$ and two mirror reflections $`\sigma _h`$ and $`\sigma _d`$. The character table is shown in Table II. There are 10 irreducible representations (irreps), of which two are two-dimensional. The AFD soft modes originating from the cubic $`\mathrm{\Gamma }_{25}(R)`$ phonons now belong either to the $`A_{1g}`$ or $`E_g`$ irreps, depending on whether the octahedron rotation axis is parallel or perpendicular to the tetragonal axis, respectively. Similarly, the modes originating from cubic $`\mathrm{\Gamma }_{15}`$ FE modes are now either $`A_{2u}`$ or $`E_u`$, depending on whether the polarization is parallel or perpendicular to the AFD axis. To label all of the displacement patterns associated with the 10-atom cell (30 degrees of freedom), we use the notation $`T\alpha K`$ to denote a displacement associated with atom type $`T`$ in pseudocubic Cartesian direction $`\alpha `$ and having “phase relation” $`K`$. The five atoms types are abbreviated as ‘$`S`$’ for Sr, ‘$`T`$’ for Ti, ‘3’ for oxygen atoms making Ti-O chains in the $`z`$ direction, and ‘1’ and ‘2’ for oxygen atoms in TiO<sub>2</sub> $`x`$-$`y`$ planes. The “phase relation” $`K`$ is either ‘$`\mathrm{\Gamma }`$’ or ‘$`R`$’ depending on whether the two atoms of the same type in the 10-atom cell move in-phase or out-of-phase (that is, whether they originate from $`\mathrm{\Gamma }`$-point or $`R`$-point modes of the parent cubic structure). Note, however, that some individual displacements contribute to more than one irrep (e.g., 1$`z\mathrm{\Gamma }`$ and 2$`z\mathrm{\Gamma }`$ contribute to both the $`A_{2u}`$ and $`B_{2u}`$ irreps). Thus, for the two in-plane oxygens we introduce alternative “type” designations $`E`$ (‘even’) and $`O`$ (‘odd’) in place of ‘1’ and ‘2’, where $$\begin{array}{cc}u_{E\alpha K}=\frac{1}{\sqrt{2}}(u_{1\alpha K}+u_{2\alpha K}),\hfill & \\ & \\ u_{O\alpha K}=\frac{1}{\sqrt{2}}(u_{2\alpha K}u_{1\alpha K}).\hfill & \end{array}$$ (26) (six degrees of freedom for each atom type), classifying them according to the irreducible representations to which they belong. Note that the four ($`E`$,$`O`$)($`x`$,$`y`$)$`R`$ modes do not have a simple one-to-one correspondence with the $`A_{1g}`$, $`A_{2g}`$, $`B_{1g}`$, and $`B_{2g}`$ modes to which they give rise; these are indicated in the table with just the notation $`R`$. Since we have access to the Hellmann-Feynman forces in our first-principles ultrasoft-pseudopotential approach, it is convenient to compute the force-constant matrix elements from finite differences as $$\mathrm{\Phi }_{i,j}^{\alpha ,\beta }=\frac{F_i^\alpha }{u_j^\beta },$$ (27) where $`F`$ is the force that results from a sufficiently small displacement $`u`$. We can use the symmetry analysis to identify the set of sublattice displacements that may participate in a given normal mode, and to calculate the forces that arise at first order with each such displacement. For example, for the $`A_{2u}`$ FE mode, we find the four displacements of types $`Sz\mathrm{\Gamma }`$, $`Tz\mathrm{\Gamma }`$, $`Ez\mathrm{\Gamma }`$, and $`3z\mathrm{\Gamma }`$ may participate. For each, a displacement amplitude of 0.2% of the lattice constant is chosen so that the harmonic approximation is still well satisfied. From each such calculation, the resulting force vector is projected onto the same set of four displacements, thus building up the $`4\times 4`$ force-constant matrix. This matrix is then symmetrized and diagonalized. In a similar way, the FE $`E_u`$ mode is represented in a 6$`\times `$6 subspace with basis $`Sx\mathrm{\Gamma }`$, $`Tx\mathrm{\Gamma }`$, $`TyR`$, 3$`x\mathrm{\Gamma }`$, $`Ex\mathrm{\Gamma }`$, and $`Oy\mathrm{\Gamma }`$. ### D FE instability in the AFD tetragonal phase Previous work of Zhong and Vanderbilt and LaSota et al. has indicated that cubic SrTiO<sub>3</sub>, in the absence of any AFD distortion, shows a FE instability. We have confirmed this result here by performing frozen-phonon calculations and building up the $`4\times 4`$ force-constant matrix for $`\mathrm{\Gamma }_{15}`$ modes polarized along $`z`$ (essentially the same procedure outlined for the $`A_{2u}`$ modes at the end of the last subsection, except performed for the cubic 5-atom cell). The non-zero eigenvalues of the corresponding dynamical matrix are given in Table IV for both the theoretical equilibrium (7.303 a.u.) and expanded experimental (7.365 a.u.) lattice constants, together with other theoretical and experimental results for comparison. Our results at the theoretical volume indicate that the FE soft phonon mode frequency remains real, although it was imaginary according to the earlier linear-response and plane-wave calculations which used slightly different lattice constants. However, at the zero-temperature experimental volume, the frequency is found to be imaginary, indicating that the FE instability is indeed very sensitive to the crystal volume. Increasing the crystal volume enhances the FE instabilities, in agreement with experiment. This is in contrast to the AFD instability, which weakens with increasing volume as shown in Sec. III.A. These results suggests that there is an inherent competition between AFD and FE modes. To understand the behavior of the FE modes in the zero-temperature AFD tetragonal structure, we next perform frozen-phonon calculations for this structure. We use the AFD state described at the end of Sec. III.B, i.e., $`P=`$5.26GPa, $`V`$ =798.989 a.u.<sup>3</sup>, $`c/a=`$ 1.004, and $`\varphi _z`$=4.93. Both $`E_u`$ and $`A_{2u}`$ modes are calculated using the method presented in Sec. III.C. Tables V and VI present the resulting eigenfrequencies and eigenvectors for the three lowest modes of each symmetry. It can be seen that both the $`E_u`$ and $`A_{2u}`$ soft-mode frequencies are imaginary (96$`i`$ and 90$`i`$ cm<sup>-1</sup>, respectively). Total-energy calculations are then performed with the corresponding eigenmodes at different mode amplitudes. The double-well energy curves are as shown in Fig. 4. One can see that the well depth of the $`E_u`$ mode is substantially greater than that of the $`A_{2u}`$ mode, indicative of a much stronger instability for a distortion of $`E_u`$ symmetry. With respect to the FE mode frequencies, the comparison with experiment is problematic because the theoretical values are imaginary (indicating an instability) while the experimental values are not (consistent with a non-FE $`T`$=0 ground state). As indicated in the Introduction, the experimental stabilization of the non-FE AFD structure is understood to result from quantum fluctuations of the atomic coordinates. It should be emphasized that the purpose of the present calculations is to study the low-temperature structural instabilities in a classical framework, i.e., in the absence of quantum fluctuations. Thus, the existence of the FE instabilities in our DFT-LDA ground state calculations, taken together with earlier quantum Monte Carlo simulations, tend to corroborate the hypothesis that the quantum fluctuations are responsible for the experimental absence of a real FE phase. Nevertheless, we can make the following comparisons. In the stress-induced Raman scattering measurement of Uwe and Sakudo, the $`E_u`$ and $`A_{2u}`$ soft mode frequencies were found to be 9.1$`\pm `$0.6 cm<sup>-1</sup> and 19$`\pm `$1 cm<sup>-1</sup> respectively (this appreciable splitting between the two modes has frequently been overlooked). The two squared frequencies produce a difference of 3$`\times `$10<sup>2</sup> cm<sup>-2</sup>. Our calculation shows the difference between the two frequencies to be 6 cm<sup>-1</sup>, roughly the same order of magnitude as the difference found experimentally. However, the theoretical difference of squared frequencies is 10$`\times `$10<sup>2</sup> cm<sup>-2</sup>, or roughly three times larger than the experimental value. This is simply because of the large magnitude of the imaginary frequencies. In summary, we found imaginary frequencies for both the FE $`E_u`$ and $`A_{2u}`$ soft modes in the AFD tetragonal phase. There is an apparent splitting between the two modes, $`\omega ^2(E_u)<\omega ^2(A_{2u})`$, suggesting that the FE structure of $`E_u`$ symmetry is more energetically favorable than the $`A_{2u}`$ one. This result is consistent with the fact that the $`A_{2u}`$ FE mode is less easily observed in neutron scattering experiments because its energy is higher than that of the $`E_u`$ phonon. ### E Influences of structural distortions on the stabilities In the previous subsection, we have calculated the FE phonon frequencies for either a cubic structure in the absence of the AFD distortion, or for a tetragonal AFD phase as observed experimentally. To understand the interaction of different distortions and their roles in affecting the FE instabilities, we performed frozen-phonon calculations for the cubic reference structure at the experimental lattice constant both with and without the tetragonal strain, and with and without the AFD rotation. In Table VII we present the results for the FE phonon frequencies for each of these scenarios. Observe that the two FE modes instabilities with different symmetries depend on the strain distortion and AFD distortion in an opposite sense. For the $`A_{2u}`$ mode, a non-zero shear strain increases the FE instability while an AFD rotation distortion reduces it. The $`E_u`$ mode, however, depends on the two distortions in the opposite way. Thus, the final sign of the frequency splitting (i.e., the sign of the anisotropy) ultimately depends sensitively on a delicate partial cancellation of the two contributions. The final result is that the $`E_u`$ mode is energetically slightly more unstable than the $`A_{2u}`$ mode. Finally, we have tested whether the use of the cubic eigen-mode distortion in place of the true tetragonal one is a good approximation for understanding the coupling of the FE and AFD instabilities. We calculated the expectation value $`\xi _c\|\mathrm{\Phi }_T\|\xi _c`$ , where $`\mathrm{\Phi }_T`$ is the dynamical matrix in the tetragonal ground state structure (fourth row of Table VII) and $`\xi _c`$ is the cubic eigen-mode distortion vector. The resulting imaginary frequencies are 85.4$`i`$ and 86.9$`i`$ cm<sup>-1</sup> for the A<sub>2u</sub> and E<sub>u</sub> modes respectively, compared with values of 90$`i`$ and 96$`i`$ cm<sup>-1</sup> in Table VII. Thus, the splitting is much smaller when the expectation value is used. This makes it clear that the nature of the FE anisotropy is quite sensitive to the tetragonal lattice strain associated with the AFD rotation distortion. This suggests a probable explanation for the fact that the classical MC simulations of Ref. , which are based on use of a cubic mode eigenvector, predicted the anisotropy incorrectly – i.e., the $`A_{2u}`$ mode was found to go soft before the $`E_u`$ one. ## IV Conclusion In this work, we have investigated both the AFD and FE structural instabilities in the tetragonal phase of SrTiO<sub>3</sub>. A unique aspect of this work is that we have studied the FE frequencies carefully using the exact eigenmode distortion obtained from the ground-state tetragonal structure, whereas previous studies have made the approximation of using the cubic eigenmode distortion instead. We show that the instabilities have a sensitive dependence on the crystal volume. We also found that the existence of the FE instabilities are affected by the coupling to the shear strain and the rotational distortion. Both types of distortion contribute, but with opposite sign, so that it is a subtle cancellation between them that determines the splitting of frequencies between the $`E_u`$ and $`A_{2u}`$ modes. The $`E_u`$ mode is found to be the more unstable of the two in the ground-state AFD phase. The degree of tilt of the oxygen octahedra is still overestimated in our calculations, compared to the experimental results. We attribute this to the quantum fluctuations of the oxygen coordinates, which tend to suppress the average AFD distortion amplitude in the experiments. The AFD and FE instabilities show opposite trends with an increase of crystal volume (weakening and strengthening, respectively). However, despite this competitive behavior, the FE instability is found to coexist in the zero-temperature AFD phase. Once again, quantum zero-point fluctuations must be invoked to explain the experimental observation that the AFD phase is stable against FE distortions. ###### Acknowledgements. Support for this work was provided by ONR Grant N00014-97-1-0048.
warning/0004/hep-ph0004190.html	ar5iv	text	# Probing top flavour-changing neutral scalar couplings at the CERN LHC ## 1 Introduction The search for flavour-changing neutral (FCN) current processes both at high and low energies is one of the major tools to test the Standard Model (SM), with the potential for either discovering or putting stringent limits on new physics. In the SM, there are no gauge FCN couplings at tree level due to the GIM mechanism, and scalar couplings are also automatically flavour-diagonal provided only one Higgs doublet is introduced. However, there is no fundamental reason for having only one Higgs doublet, and for two or more scalar doublets FCN couplings are generated at tree level unless an ad hoc discrete symmetry is imposed . There are stringent limits on the FCN couplings between light quarks arising from the smallness of the $`K^0\overline{K}^0`$, $`B^0\overline{B}^0`$ and $`D^0\overline{D}^0`$ mass differences , as well as from the CP violation parameter $`ϵ`$ of the neutral kaon system. If there is no suppression of scalar FCN couplings, the masses of the neutral Higgs bosons have to be of the order of 1 TeV. However, if FCN couplings between light quarks are suppressed, being for instance proportional to the quark masses or to some combination of Cabibbo-Kobayashi-Maskawa (CKM) mixing angles , the lightest neutral Higgs may have a mass $`M_H`$ of around 100 GeV. In this case, the effects of scalar FCN currents involving the top quark can be quite sizeable. The interactions among the top, a light quark $`q=u,c`$ and a Higgs boson $`H`$ will be described by the effective Lagrangian $$=\frac{g_W}{2\sqrt{2}}g_{tq}\overline{q}(c_v+c_a\gamma _5)tH+\mathrm{h}.\mathrm{c}.$$ (1) The axial and vector parts of the coupling are normalized to $`\|c_v\|^2+\|c_a\|^2=1`$. The natural size of $`g_{tc}`$ in a two Higgs doublet model (2HDM) is $`g_{tc}0.20`$, leading to large FCN current effects observable at the CERN Large Hadron Collider (LHC). The natural size of $`g_{tu}`$ ranges from 0.012 to 0.025, depending on the model considered. This is too small to be observed, but if this coupling is enhanced by a factor $`2`$, it could also be detected. Models with new heavy exotic quarks can also have large tree level $`Htq`$ couplings, with the particularity that they are proportional to the FCN $`Ztq`$ couplings . Present limits on the latter allow for $`g_{tu}0.31`$ or $`g_{tc}0.35`$, both observable at LHC. The large top quark mass opens the possibility of having relatively large effective vertices $`Htq`$ induced at one loop by new physics. In the Minimal Supersymmetric Standard Model (MSSM) there are large regions of the parameter space where one can have $`g_{tc}0.04`$ <sup>1</sup><sup>1</sup>1The value of $`g_{tu}`$ is much smaller in principle.. In contrast, in the SM the effective $`Htc`$ vertex is very small, $`g_{tc}10^6`$ due to a strong GIM cancellation . The large top quark mass also suggests that the top quark could play a fundamental rôle in the electroweak symmetry breaking mechanism and in this case top FCN scalar couplings would be large, $`g_{tc}0.30.4`$ or $`g_{tc}0.52.9`$ , depending on the model considered. Hence the importance of measuring its couplings, in particular to the Higgs boson . Present experimental bounds on top FCN scalar (and gauge) couplings are very weak. Some processes have been proposed to measure the $`Htq`$ vertices, for instance $`t\overline{c}\nu _e\overline{\nu }_e`$ and $`t\overline{c}e^+e^{}`$ production at a linear collider with a center of mass (CM) energy of $`12`$ TeV and $`t\overline{c}`$ production at a muon collider or at hadron colliders via gluon fusion $`ggHt\overline{c}`$ . These processes can probe the FCN couplings of the heavier mass eigenstates, but do not provide useful limits for the light scalars. In this Letter we will show that top decays $`tHq`$ at LHC provide the best limits on the couplings of a light scalar of a mass slightly above 100 GeV. The importance of this mass range stems from the fact that global fits to electroweak observables in the SM seem to prefer $`M_H100`$ GeV, as well as the MSSM prediction that the mass of the lightest scalar has to be less than 130 GeV . On the other hand, direct searches at the CERN $`e^+e^{}`$ collider LEP at CM energies up to 202 GeV imply the bound $`M_H>107.9`$ GeV for the mass of a SM Higgs, with a 95% confidence level (CL) . For the lightest Higgs of the MSSM the bounds depend on the choice of parameters and typically are reduced to $`M_H>88.3`$ GeV. In the following we will perform a detailed analysis of the sensitivity of the LHC to FCN scalar couplings in top decays $`tHq`$ . We will also study $`Ht`$ production, which also serves to constrain the $`Htq`$ vertex . However, the bounds obtained from this process are much less restrictive. Finally, we will comment on the potential of the LHC to discover these signals at the rates predicted by some models. ## 2 Limits on FCN couplings from top decays To obtain constraints on the coupling $`g_{tq}`$ we will first consider $`t\overline{t}`$ production, with the top quark decaying to $`W^+bl\nu b`$, $`l=e,\mu `$ and the antitop decaying to $`H\overline{q}`$. With a good $`\tau `$ identification the decays $`W^+\tau ^+\nu `$ can be included, improving the statistical significance of the signal. The decay mode of the Higgs boson depends strongly on its mass. For our evaluation, we take $`M_H=120`$ GeV, well above the LEP limits, and decaying into $`b\overline{b}`$ pairs. In principle, there is no reason to assume that the $`Hbb`$ coupling is suppressed. In a general 2HDM, both doublets couple to the $`b`$ quark and, unless some fine-tuned cancellation occurs, the coupling of the physical light Higgs to the $`b`$ quark is proportional to the $`b`$ mass, which we take $`m_b=3.1`$ GeV at the scale $`\mu =M_H`$. To be conservative, we assume that the Higgs couplings to the gauge bosons are not suppressed and are the same as in the SM. Thus for $`M_H=120`$ GeV $`H`$ decays predominantly into $`b\overline{b}`$ pairs, with a branching ratio of $`0.7`$. For $`M_H>130`$ GeV, the $`W^+W^{}`$ decay mode, with one of the $`W`$’s off-shell, begins to dominate <sup>2</sup><sup>2</sup>2Of course, if the couplings to the gauge bosons are suppressed, for instance if $`H`$ is a pseudoscalar, the signal will get an enhancement factor of roughly $`9/7`$ and the $`b\overline{b}`$ decay mode will dominate also for larger values of $`M_H`$.. The signal is then $`l\nu bb\overline{b}j`$, with three $`b`$ quarks in the final state. Clearly, $`b`$ tagging is crucial to separate the signal from the backgrounds. We assume a $`b`$ tagging efficiency of 50%, and a mistag rate of 1%, similar to the values recently achieved at SLD . We evaluate the signal using the full $`gg,q\overline{q}t\overline{t}W^+bH\overline{q}l\nu bb\overline{b}q`$ matrix elements with all intermediate particles off-shell. We calculate the matrix elements using HELAS code generated by MadGraph and modified to include the decay chain. For definiteness we assume $`c_v=1`$, $`c_a=0`$, but we check that the difference when using other combinations is below $`1\%`$. We normalize the cross section with a $`K`$ factor of 2.0 . The main background to this signal comes from $`t\overline{t}`$ production with standard top and antitop decays, $`tW^+bl\nu b`$, $`\overline{t}W^{}\overline{b}jj\overline{b}`$. This process mimics the signal if one of the jets resulting from the $`W^{}`$ decay is mistagged as a $`b`$ jet. The matrix elements for this process are calculated with HELAS and MadGraph in the same way as the signal, taking again $`K=2.0`$. Another background to be considered is $`Wb\overline{b}jj`$ production, for which we use VECBOS modified to include energy smearing and trigger and kinematical cuts. For this process we take $`K=1.1`$ . $`Wjjjj`$ production is negligible after using $`b`$ tagging. For the phase space integration we use MRST structure functions set A with $`Q^2=\widehat{s}`$. After generating signals and backgrounds, we simulate the detector resolution effects with a Gaussian smearing of the energies of jets $`j`$ and charged leptons $`l`$ , $$\frac{\mathrm{\Delta }E^l}{E^l}=\frac{10\%}{\sqrt{E^l}}0.3\%,\frac{\mathrm{\Delta }E^j}{E^j}=\frac{50\%}{\sqrt{E^j}}3\%,$$ (2) where the energies are in GeV and the two terms are added in quadrature. We then apply detector cuts on transverse momenta $`p_T`$, pseudorapidities $`\eta `$ and distances in $`(\eta ,\varphi )`$ space $`\mathrm{\Delta }R`$: $$p_T^l15\mathrm{GeV},p_T^j20\mathrm{GeV}$$ $$\|\eta ^{l,j}\|2.5,\mathrm{\Delta }R_{jj,lj}0.4.$$ (3) For the events to be triggered, we require both the signal and background to fulfill at least one of the trigger conditions . These conditions are different in the low 10 fb<sup>-1</sup> (L) and high 100 fb<sup>-1</sup> (H) luminosity Runs. For our processes, the relevant conditions are * one jet with $`p_T180`$ GeV (L) / 290 GeV (H), * three jets with $`p_T75`$ GeV (L) / 130 GeV (H), * one charged lepton with $`p_T20`$ GeV (L) / 30 GeV (H), * missing energy $`E_T\overline{)}50`$ GeV (L) / 100 GeV (H) and one jet with $`p_T50`$ GeV (L) / 100 GeV (H). We reconstruct the signal as follows. There are three tagged $`b`$ jets $`b_{13}`$ in the final state and a non-$`b`$ jet $`j`$. There are three possible pairs $`b_i,b_j`$, with invariant masses $`M_{b_ib_j}`$, one of which results from the decay of the Higgs boson and has an invariant mass close to $`M_H`$. The invariant mass of this pair and the jet $`j`$, $`M_{b_ib_jj}`$, is also close to the top mass. Hence we choose the pair $`b_i,b_j`$ which minimizes $`(M_{b_ib_j}M_H)+(M_{b_ib_jj}m_t)`$, defining the Higgs reconstructed mass as $`M_H^{\mathrm{rec}}=M_{b_ib_j}`$ and the antitop reconstructed mass as $`m_{\overline{t}}^{\mathrm{rec}}=M_{b_ib_jj}`$. In Figs 12 we plot the kinematical distributions of these variables for the signal and for the $`t\overline{t}`$ background. The $`m_{\overline{t}}^{\mathrm{rec}}`$ distribution is slightly broader for $`t\overline{t}`$ due to the signal reconstruction method. The remaining jet $`b_k`$ and the charged lepton result from the decay of the top quark. To reconstruct its mass, we make the hypothesis that all missing energy comes from a single neutrino with $`p^\nu =(E^\nu ,p_T\overline{)},p_L^\nu )`$, and $`p_T\overline{)}`$ the missing transverse momentum. Using $`(p^l+p^\nu )^2=M_W^2`$ we find two solutions for $`p^\nu `$, and we choose that one making the reconstructed top mass $`m_t^{\mathrm{rec}}\sqrt{(p^l+p^\nu +p^{b_k})^2}`$ closest to $`m_t`$. In Fig. 3 we plot the kinematical distribution of this variable for the signal and $`t\overline{t}`$ background. The reconstruction of the signal and the requirement $`m_t^{\mathrm{rec}},m_{\overline{t}}^{\mathrm{rec}}m_t`$, $`M_H^{\mathrm{rec}}M_H`$ are sufficient to eliminate the $`Wb\overline{b}jj`$ background, but hardly affect $`t\overline{t}`$, as can be seen in Figs. 13. To improve the signal to background ratio we reject events when one pair of jets, $`b_ij`$ or $`b_jj`$, seems to be the product of the hadronic decay of a $`W`$ boson, as happens for the $`t\overline{t}`$ background. We define the reconstructed $`W`$ mass as the invariant mass $`M_{b_ij}`$, $`M_{b_jj}`$ closest to $`M_W`$ (see Fig. 4) and impose a veto cut on events with $`M_W^{\mathrm{rec}}M_W`$. The complete set of kinematical cuts for both Runs is summarized in Table 1. We also require a large transverse energy $`H_T`$. Alternatively, we can use an artificial neural network (ANN) as a classifier to distinguish between the signal and the $`t\overline{t}`$ background. In addition to $`M_H^{\mathrm{rec}}`$, $`m_{\overline{t}}^{\mathrm{rec}}`$, $`m_t^{\mathrm{rec}}`$ and $`M_W^{\mathrm{rec}}`$, we consider $`p_T^H`$, the transverse momentum of the reconstructed Higgs boson, $`p_T^{\mathrm{max}}`$ and $`p_T^{\mathrm{min}}`$, the maximum and minimum transverse momentum of the jets, and $`p_T^{b,\mathrm{max}}`$, the maximum transverse momentum of the two $`b`$’s which reconstruct the Higgs boson. The kinematical distributions of these variables are shown in Figs. 58. We observe that these variables are clearly not suitable to perform kinematical cuts on them. However, their inclusion as inputs to an ANN greatly improves its ability to classify signal and background events. We have not found any improvement adding other variables, and in some cases the results are worse. To construct the ANN we use JETNET 3.5 . We find convenient using a three-layer topology of 8 input nodes, 12 hidden nodes and 1 output node. However, we do not claim that either the choice of variables or the network topology are the best ones. The performance of the ANN is not very sensitive to the number of hidden nodes. The input variables are normalized to lie approximately in the interval $`[1,1]`$ to reduce the training time. The network output $`r`$ is set to one for the signal and zero for $`t\overline{t}`$. We train the network using two sets with roughly 40000 signal and 40000 background events and the standard backpropagation algorithm. To test the ability of the network to classify never seen data we use two other sets with the same size. One frequent problem is overtraining. To avoid it, we keep training while the network error evaluated on the test sets, $`E_t^{(N)}`$, decreases for the successive training epochs $`N`$. When it begins to increase, we stop training. To avoid fluctuations, we keep track of $`E_t^{(N)}`$ for previous epochs. When we reach an epoch $`N_0`$ when $`E_t^{(N_0100)}`$ is smaller than all the following values, we stop training and start again until epoch $`N_0100`$. We repeat the same procedure using four other different training and test sets of the same size, and find that the results are fairly stable. The final training is done including the 4 training and 4 test sets. At the end of the training, most of the signal and background events in the test samples concentrate in very narrow intervals $`[0.99,1]`$ and $`[0,0.01]`$, respectively. To classify signal and background we use two other signal and background sets with never seen before data, 50 and 100 times larger, respectively, than the ones used in training, to avoid possible statistical fluctuations. To separate signal from background we simply require $`r>0.998`$. This maintains 29% of the signal while it rejects 99.9% of the $`t\overline{t}`$ background. It is not necessary to train another ANN to distinguish the $`Wb\overline{b}jj`$ background, the same network with $`r>0.998`$ completely eliminates it. In Table 2 we collect the number of events without cuts, with the standard cuts in Table 1 and with the ANN cut $`r>0.998`$, for LHC Runs L and H. We normalize the signal to $`g_{tq}=0.2`$. Although for definiteness we train the ANN with a vector FCN coupling ($`c_v=1`$, $`c_a=0`$), we check that the same ANN correctly recognizes the signal events when we consider an axial, left- or right-handed FCN coupling. The differences between the cross sections are in all cases below $`1\%`$, and comparable to the statistical Monte Carlo uncertainty. The differences between the cross sections after standard kinematical cuts are also smaller than $`1\%`$. To derive upper bounds on the FCN couplings we follow the Feldman-Cousins construction . For the numerical evaluation of the confidence intervals we use the PCI package . If there is no evidence of this process, i. e., if the number of events observed $`n_0`$ is equal to the expected background $`n_b`$, we obtain independent limits on the couplings $`g_{tu}`$, $`g_{tc}`$. Using the standard cuts we obtain $`g_{tq}0.071`$ at Run L and $`g_{tq}0.041`$ at Run H. Using the ANN cut improves these limits, $`g_{tq}0.064`$ at Run L and $`g_{tq}0.035`$ at Run H. All bounds are calculated with a 95% CL. Although the improvement may not seem very significant, in Section 4 we will see that the luminosity required to see a positive signal is reduced to one half with the help of the neural network. It is worth to give here also the results for the Fermilab Tevatron. This process at Run II with a CM energy of 2 TeV and a luminosity of 20 fb<sup>-1</sup> gives only $`g_{tq}0.38`$ (using standard kinematical cuts) due to the smaller statistics available. ## 3 Limits on FCN couplings from $`Ht`$ production The existence of top FCN Higgs couplings leads to $`Ht`$ production through the diagrams in Fig. 9. This process has a negligible cross section in the SM . As in the previous section, we consider only the leptonic decays of the top quark, and the signal is $`l\nu bb\overline{b}`$. Note that $`Ht`$ production is completely analogous to $`Zt`$ production via $`Ztq`$ anomalous couplings with $`Zb\overline{b}`$ , and has the same large backgrounds as this decay mode. The most dangerous is $`t\overline{t}`$ production, with $`tW^+bl\nu b`$, $`\overline{t}W^{}\overline{b}jj\overline{b}`$, when one of the jets resulting from the $`W^{}`$ decay is missed by the detector and the other is mistagged as a $`b`$ jet. Another small background is $`Wb\overline{b}j`$ production. The $`Ht`$ process has a smaller cross section than $`t\overline{t}`$ with subsequent decay $`\overline{t}H\overline{q}`$ for the same values of $`g_{tq}`$, but also receives an additional contribution from the latter when $`q`$ is missed by the detector. This correction is not very significant for the analogous case of $`Zt`$ production, about $`+10\%`$ for the up quark and $`+50\%`$ for the charm, but in this case with $`M_H>M_Z`$ the additional contribution becomes more important. In our calculations we will then take into account both processes. The signals and backgrounds are generated as for the $`l\nu bb\overline{b}j`$ case. To estimate the number of events in which a jet is missed by the detector, we consider that a jet is missed when it has a low transverse momentum $`p_T^j<20`$ GeV or a very large pseudorapidity $`\eta ^j>3`$. The signal is reconstructed choosing first the pair of $`b`$ jets whose invariant mass is closest to $`M_H`$. The remaining $`b`$ jet is assigned to the decay of the top quark, whose mass is reconstructed as in the previous Section. To improve the signal to background ratio we use the kinematical cuts shown in Table 3. In this case we cannot reduce the background imposing a veto cut to reject events in which a pair of jets has an invariant mass around $`M_W`$, because one of these jets is missed by the detector. Using an ANN does not improve the situation. The number of events for the $`guHt`$ and $`gcHt`$ signals, the additional contribution $`t\overline{t}Ht(q)`$ and their backgrounds is collected in Table 4. The limits obtained from the $`Ht`$ process are much less restrictive than from top decays. If no signal is observed, we obtain from Table 4 $`g_{tu}0.13`$, $`g_{tc}0.19`$ at Run L and $`g_{tu}0.076`$, $`g_{tc}0.11`$ at Run H. The same is true for Tevatron Run II, where the bounds obtained are very poor, $`g_{tu}0.40`$, $`g_{tc}1.2`$ ## 4 LHC discovery potential In the previous Sections we have considered that the signal is not seen, and have obtained 95% CL upper bounds on the size of the FCN couplings. These bounds imply the limits in the branching ratios of $`tHq`$ in Table 5. The limits obtained for different values of $`M_H`$ are very similar. For a lighter Higgs, the larger phase space (in the case of top decays) or the larger structure functions (in $`Ht`$ production) increase the signal cross section and hence allow a better determination of $`g_{tq}`$. This improvement disappears when expressing the limit in terms of a branching ratio. After repeating the analysis using $`M_H=110`$, $`\mathrm{Br}(Hb\overline{b})=0.8`$, we find limits $`3\%`$ smaller from top decays and $`5\%`$ smaller from $`Ht`$ production. A heavier Higgs has the opposite behaviour, with the disadvantage of a smaller branching ratio $`Hb\overline{b}`$. With $`M_H=130`$, $`\mathrm{Br}(Hb\overline{b})=0.45`$ the limits from top decays and $`Ht`$ production are 1.5 and 1.3 times larger, respectively. Let us turn now to the possibility of observing these anomalous decays. The most common criterion to estimate the discovery potential is to use the discovery significance defined as the ratio of the expected signal divided by the expected background fluctuation, $`sn_s/\sqrt{n_b}`$. This is adequate when $`n_b>5`$ and the Poisson distribution of the background can be approximated by a Gaussian with standard deviation $`\sqrt{n_b}`$. One has “evidence” ($`3\sigma `$) when $`s>3`$, and “discovery” ($`5\sigma `$) occurs with $`s>5`$. In a general 2HDM, the most common ansatz for the FCN couplings is to assume that they scale with the quark masses. Under this assumption, the expected size of these couplings is $`g_{tc}\sqrt{m_tm_c}/M_W=0.20`$, $`g_{tu}\sqrt{m_tm_u}/M_W=0.012`$. An $`Htc`$ coupling of this size (leading to $`\mathrm{Br}(tHc)=1.5\times 10^3`$) will easily be detected and precisely measured. In top decays, at Run H we have $`s=48`$ using standard cuts and $`s=68`$ using the network cut. The benefit of the ANN is clear: the luminosity required to discover a signal is reduced to one half. In $`Ht`$ production, the significance is much smaller, $`s=7`$. The minimum coupling which can be discovered with $`5\sigma `$ is $`g_{tc}=0.054`$ ($`\mathrm{Br}(tHc)=1.1\times 10^4`$), and the minimum coupling which can be seen with $`3\sigma `$ is $`g_{tc}=0.042`$ ($`\mathrm{Br}(tHc)=6.5\times 10^5`$). An $`Htu`$ coupling $`g_{tu}=0.012`$ is too small to be seen, because in top decays we have only $`s=0.23`$, and in $`Ht`$ production $`s=0.046`$. Note that for a 2HDM this ansatz implies that there are other interesting FCN current processes, for instance $`H\mu \tau `$ at a muon collider . Another possibility is to consider that the FCN couplings between two quarks are proportional to some CKM mixing angles involving these quarks. Rephasing invariance implies that these couplings are of the form $`g_{qq^{}}=\lambda V_{qi}V_{q^{}i}^{}`$, with $`i=d,s,b`$ . Choosing $`i=b`$, the hierarchy among the mixing angles guarantees a small value for $`g_{uc}=\lambda V_{ub}V_{cb}^{}`$. With $`\lambda =5`$ we recover the values of $`g_{uc}`$ and $`g_{tc}`$ of the previous ansatz, but $`g_{tu}`$ is a factor of two larger, $`g_{tu}=0.025`$. An Htu coupling of this size is now near the detectable level: $`s=1.1`$ in top decays and in $`s=0.22`$ in $`Ht`$ production, and if it is enhanced by a factor $`2`$ it would also be detected. In models with exotic quarks the top FCN scalar couplings can be large, $`g_{tu}0.31`$ or $`g_{tc}0.35`$ (but not both simultaneously). These are the only models where the top can naturally have a large mixing with the up quark. However, the effects of top mixing in these models are most clearly seen with the appearance of tree level $`Ztq`$ couplings, leading to top anomalous decays $`tZq`$ and $`Zt`$ production . If these processes are not observed at LHC, the $`Htq`$ couplings (proportional to the FCN $`Ztq`$ couplings) must be small, $`g_{tu}0.024`$, $`g_{tc}0.030`$, below the detectable level. To analyze the case of top decays $`tHq`$ in the MSSM it is necessary to recalculate the signal and backgrounds and train again the ANN with the new data. The reason is that in the region of parameters of interest the Higgs width is much larger due to the enhanced $`Hbb`$ coupling, and the $`M_H^{\mathrm{rec}}`$ distribution is broader in principle. For consistency we will use the parameters of Ref. : $`\mathrm{tan}\beta =35`$, a pseudoscalar mass $`M_A=120`$ GeV and $`c_v=c_a=1/\sqrt{2}`$. With these parameters, the Higgs mass is in the interval $`M_H=105120`$ GeV . The kinematics is fairly independent of the particular value of $`M_H`$ used, as we have shown above, the only difference being the Higgs decay rate to $`b\overline{b}`$ which varies slightly in this mass range. For our analysis we take $`M_H=110`$ GeV, $`\mathrm{Br}(Hb\overline{b})=0.9`$ . Our limits on $`\mathrm{Br}(tHq)`$ are proportional to $`1/\mathrm{Br}(Hb\overline{b})`$. To separate signal from background we select events with network output $`r>0.994`$. There are large regions in the MSSM parameter space where $`\mathrm{Br}(tHc)`$ is greater than $`1.1\times 10^4`$ and even reaches $`4\times 10^4`$ . These decays can be discovered at LHC Run H with $`5\sigma 18\sigma `$. Also, a rate $`\mathrm{Br}(tHc)=6.7\times 10^5`$ is observable with $`3\sigma `$. The decays $`tHu`$ have a too small rate to be seen. On the other hand, if no signal is observed we obtain a limit $`g_{tq}0.054`$ at Run L and $`g_{tq}0.031`$ at Run H. Note that although these limits seem better than those obtained for $`M_H=120`$, when translated into branching ratios they give very similar bounds, $`\mathrm{Br}(tHq)1.4\times 10^4`$ at Run L and $`\mathrm{Br}(tHq)4.7\times 10^5`$ at Run H. We are indebted to J. Guasch for discussions and for providing us with data for the MSSM Higgs analysis. We have also benefited from discussions with F. del Aguila, Ll. Ametller and J. Seixas. J. A. A. S. thanks the members of the CFIF for their hospitality during the realization of this work. This work was partially supported by CICYT under contract AEN96–1672 and by the Junta de Andalucía, FQM101.
warning/0004/cond-mat0004192.html	ar5iv	text	# Zero-temperature equation of state of quasi-one dimensional H2 ## Abstract We have studied molecular hydrogen in a pure 1D geometry and inside a narrow carbon nanotube by means of the diffusion Monte Carlo method. The one-dimensionality of H<sub>2</sub> in the nanotube is well maintained in a large density range, this system being closer to an ideal 1D fluid than liquid <sup>4</sup>He in the same setup. H<sub>2</sub> shares with <sup>4</sup>He the existence of a stable liquid phase and a quasi-continuous liquid-solid transition at very high linear densities. The experimental finding in 1991 of carbon nanotubes opened brand new possibilities both in technology and in fundamental physics. The nanoscale of these new materials has led to the discovery of novel mechanical, chemical and electrical properties that suggest exciting new technological applications. One of the more relevant features of these new materials is their large adsorption capability compared with that of a graphite planar substrate. In fact, the enhanced Van der Waals interaction when particles are adsorbed inside a single nanotube, or in the interstitial sites of single wall carbon nanotube (SWCN) bundles, may allow the density of the adsorbed substrate to be high enough for the bundles to serve as gas storage devices. Special interest exists in the physisorption of hydrogen in the quest for a fuel cell efficient enough to be used as a pollution-free energy carrier. It has been recently suggested that SWCN with diameters of the order of a nanometer can be the best candidates to approach the pursued level of packing . From a more fundamental point of view, the strong confinement of particles adsorbed in the carbon channels of a SWCN bundle, with diameters ranging form 7 to 40 Å and an width to length ratio of $``$ 1000, offers the possibility of an experimental realization of a quasi-one dimensional system. Moreover, if the temperature is low enough one is dealing with a unique opportunity of studying a nearly one-dimensional quantum fluid. In a recent experiment, Teizel et al. have unambiguously observed the quasi-one dimensional behavior of <sup>4</sup>He adsorbed in a SWCN bundle by measuring its desorption rate. On the other hand, theoretical studies in the limit of zero temperature and strictly one dimension have proved the existence of a liquid state with a binding energy in the milli-Kelvin scale . Using the diffusion Monte Carlo (DMC) method, we also compared the ideal 1D geometry system with a realistic situation in which <sup>4</sup>He is adsorbed in a narrow nanotube . The aim of the present work is to extend the theoretical study to the appealing case of H<sub>2</sub> that, besides its technological relevance, might offer the existence of a homogeneous liquid phase at zero temperature. It is worth noticing that both bulk and two-dimensional H<sub>2</sub>, are solid in this temperature limit. Liquid phases have only been observed in theoretical calculations of small clusters and in 2D geometries with localized alkaline impurities . We have studied molecular hydrogen at zero temperature in a one-dimensional (1D) array and inside a single walled carbon nanotube (T) of radius $`R`$ = 3.42 Å, which corresponds to a (5,5) armchair tube . The technique used is the DMC method, which has become in the last decades one of the most efficient theoretical tools, from the microscopic point of view, to deal with quantum fluids. The DMC method solves the $`N`$-body imaginary-time Schrödinger equation $$\frac{f(𝐑,t)}{t}=D_R^2f(𝐑,t)+D_R(𝐅f(𝐑,t))+(E_L(𝐑)E)f(𝐑,t),$$ (1) in which $`f(𝐑,t)`$ = $`\psi (𝐑)\mathrm{\Psi }(𝐑,t)`$, $`D=\mathrm{}^2/2m`$, $`E_L(𝐑)=\psi (𝐑)^1H\psi (𝐑)`$ is the local energy, and $`𝐅(𝐑)=2\psi (𝐑)^1\psi (𝐑)`$ is the so-called drift force. The time-independent wave function $`\psi (𝐑)`$ acts as an important sampling function and is chosen as a trial model usually adjusted at a variational Monte Carlo (VMC) level. In bosonic systems like the present one the DMC method provides exact results apart from statistical uncertainties. More specific details of DMC may be drawn from Refs. . A relevant issue in a microscopic study is the nature of the interspecies interaction. We have considered the H<sub>2</sub> molecules interacting via the isotropic semiempirical potential from Silvera and Goldman (SG) that has been extensively used in path integral Monte Carlo (PIMC) and DMC calculations of bulk , clusters and H<sub>2</sub> films . The SG is a pair potential that incorporates to some extent the effect of three-body interactions by means of an effective two-body term of the form $`C_9/r^9`$. On the other hand, the isotropy of the potential is well justified if one considers that at very low temperatures almost all the H<sub>2</sub> molecules are para-hydrogen species, i.e., they are in the J=0 rotational state. In the simulations of H<sub>2</sub> inside a nanotube, we consider a cylindrically symmetric potential as suggested by Stan and Cole . In that simplified model, the interactions between C atoms and H<sub>2</sub> molecules are axially averaged out resulting in a potential which only depends on the distance to the center of the tube. It has been proved that the differences between that smoothed potential and a potential which is built up as an explicit sum of individual C-H<sub>2</sub> interactions are not significant and surely smaller than the relative uncertainty in the ($`\sigma `$,$`ϵ`$) Lennard-Jones parameters. Considering $`\sigma `$ = 2.97 Å and $`ϵ`$ = 42.8 K, the symmetric potential felt by a H<sub>2</sub> molecule in a (5,5) tube has a depth of 42 $`ϵ`$, three times larger that the attraction of the same molecule in a flat graphitic surface. The use of importance sampling in DMC requires the introduction of a trial wave function that guides the diffusion process to relevant regions of the walkers phase space. In the 1D system, we consider $$\mathrm{\Psi }^{1\mathrm{D}}(𝐑)=\mathrm{\Psi }_\mathrm{J}(𝐑),$$ (2) with $`\mathrm{\Psi }_\mathrm{J}(𝐑)=_{i<j}\mathrm{exp}\left[\frac{1}{2}\left(\frac{b}{r_{ij}}\right)^5\right]`$ a Jastrow wave function with a McMillan two-body correlation factor. Inside the nanotube, H<sub>2</sub> molecules interact with the walls of the cylinder and therefore we have added an additional one-body term $$\mathrm{\Psi }^\mathrm{T}(𝐑)=\mathrm{\Psi }_\mathrm{J}(𝐑)\mathrm{\Psi }_\mathrm{c}(𝐑)$$ (3) with $`\mathrm{\Psi }_\mathrm{c}(𝐑)=_i^N\mathrm{exp}(cr_i^2)`$, $`r_i`$ being the radial distance of the particle to the center, to avoid the hard core of the H<sub>2</sub>-nanotube interaction. Theoretical calculations of 1D <sup>4</sup>He agree in predicting a liquid-solid phase transition at high linear densities . This transition which is only possible at absolute zero temperature looks like a nearly continuous one without a measurable difference between the melting and freezing densities. Following the same procedure that in our previous work in helium, we have explored the existence of such a transition in H<sub>2</sub>. In this ordered phase, we modify the trial wave function in both the 1D system and inside the nanotube by multiplying them by a $`z`$-localized factor $`\mathrm{\Psi }_\mathrm{s}(𝐑)=_i^N\mathrm{exp}(a(z_iz_{is})^2)`$. The sites $`z_{is}`$ are equally-spaced points in both the 1D line and the axial direction of the nanotube. The variational parameters $`a`$, $`b`$ and $`c`$ have been optimized by means of VMC calculations. In the liquid phase and near the equilibrium density $`b`$= 3.759 Å, with a slight increase with density (at $`\lambda `$ = 0.277 Å<sup>-1</sup>, $`b`$ = 3.789 Å), whereas $`c`$ = 4.908 Å<sup>-2</sup> is kept fixed for all $`\lambda `$ values. In the solid phase, $`b`$ = 3.404 Å, $`a`$ = 0.799 Å<sup>-2</sup> and $`c`$ = 5.136 Å<sup>-2</sup>, with a negligible $`\lambda `$ dependence in the region analyzed. The possible existence of a liquid-solid phase transition at high linear density has been studied in both 1D and inside a narrow nanotube. In Table I, results for the energy per particle of both systems are reported for the liquid ($`a`$ = 0) and solid ($`a`$ 0) phases. The comparison between the energies of both phases at the same density shows that their difference changes sign in going from $`\lambda `$ = 0.312 Å<sup>-1</sup> to $`\lambda `$ = 0.304 Å<sup>-1</sup> in 1D and from $`\lambda `$ = 0.320 Å<sup>-1</sup> to $`\lambda `$ = 0.312 Å<sup>-1</sup> in the tube. Above these densities, the system prefers to be localized in a solid-like structure with a difference $`\|E(s)E(l)\|`$ that increases with $`\lambda `$. When the density decreases the liquid phase is energetically preferred and again the size of the difference $`\|E(s)E(l)\|`$ increases when $`\lambda `$ diminishes. The density value at which this difference becomes zero is estimated to be $`\lambda =0.309`$ Å<sup>-1</sup> in 1D and $`\lambda =0.315`$ Å<sup>-1</sup> in the tube, being not possible to distinguish between freezing and melting densities. As previously studied in <sup>4</sup>He it appears to be a nearly continuous phase transition located at a density close to the inverse of the location of the minimum of the respective pair potential ($`r_m^1`$ = 0.337 Å<sup>-1</sup> vs $`\lambda _s`$ = 0.358 Å<sup>-1</sup> for helium, and $`r_m^1`$ = 0.291 Å<sup>-1</sup> vs $`\lambda _s`$ = 0.309, 0.315 Å<sup>-1</sup> for molecular hydrogen). Inside the nanotube (T), the energies are much more negative that in 1D due to the strong attraction of the carbon substrate: the binding energy of a single H<sub>2</sub> molecule in the tube is $`E_b=1539.87\pm 0.11`$ K. Looking at the T-energy results contained in Table I one realizes that also in this case a transition occurs at a density very close to the 1D one. It is remarkable that both in 1D and T, H<sub>2</sub> remains liquid below the liquid-solid transition density, and thus a homogeneous liquid phase at zero pressure is predicted. That result contrasts with the theoretically and experimentally well established solid phase in 3D and the 2D solid phase predicted by a PIMC calculation . The equations of state of liquid H<sub>2</sub> near the equilibrium density for both the 1D and T systems are shown in Fig. 1. In order to make the energy scales compatible we have subtracted the single binding energy $`E_b`$ to the T results. The lines in the figure correspond to the third-degree polynomial fits in the form $$\frac{E}{N}=e_0+A\left(\frac{\lambda \lambda _0}{\lambda _0}\right)^2+B\left(\frac{\lambda \lambda _0}{\lambda _0}\right)^3.$$ (4) The best set of parameters $`e_0`$, $`\lambda _0`$, $`A`$ and $`B`$ are reported in Table II. The equilibrium densities in both systems are the same considering their respective uncertainties but the binding energy $`e_0=e(\lambda _0)`$ is larger when H<sub>2</sub> is inside the nanotube. The difference between the 3D geometry (T) and the idealized one (1D) can be quantified by means of the adimensional parameter $$\mathrm{\Delta }^\mathrm{T}=\frac{(E^\mathrm{T}E_b^\mathrm{T})E_{1\mathrm{D}}}{(E^\mathrm{T}E_b^\mathrm{T})}$$ (5) In the present system, around $`\lambda _0`$, $`\mathrm{\Delta }^\mathrm{T}`$ = 3.5 % which emphasizes the proximity between the real system and the idealized one. It is worth noting that in <sup>4</sup>He inside the same nanotube we obtained $`\mathrm{\Delta }^\mathrm{T}=90`$ % , and therefore H<sub>2</sub> seems a better candidate to experimentally achieve a 1D condensed phase. That significant difference between helium and hydrogen may be understood taking into account that the difference between the hard-core size of the C-He interaction ($`\sigma _{CHe}`$ = 2.74 Å) and the C-H<sub>2</sub> one ($`\sigma _{CH_2}`$ = 2.97 Å) is magnified in a (5,5) tube because its small radius ($`R`$ = 3.42 Å). In Fig. 2, the density dependence of the pressure for both the 1D and T systems is reported from equilibrium up to the liquid-solid transition density. As a matter of comparison, the same results for <sup>4</sup>He are also plotted. Both in H<sub>2</sub> and <sup>4</sup>He the pressure increases faster in the 1D geometry ($`P_\lambda `$) than in the tube (P) due the transverse degree of freedom that particles have in the latter case (notice the proportionality between the scales of $`P`$ and $`P_\lambda `$ in Fig. 2, $`P_\lambda /P=\pi R^2`$). At a given density $`\lambda `$, the difference between the T and 1D pressures is smaller in H<sub>2</sub> than in <sup>4</sup>He. For example, at the respective transition densities that difference is more than one and a half times larger in helium than in hydrogen. Therefore, the one-dimensionality of H<sub>2</sub> inside the nanotube is well maintained in all the liquid regime in contrast with <sup>4</sup>He, in which the departure from such an idealized model already appears around the equilibrium density and increases significantly with $`\lambda `$. Also apparent form Fig. 2 is a much smaller compressibility in H<sub>2</sub> than in <sup>4</sup>He. In the 1D geometry at $`\lambda `$ = $`\lambda _0`$ the velocity of sound in H<sub>2</sub> is $`c`$ = 736.1 $`\pm `$ 0.2 m/sec to be compared with $`c`$ = 7.98 $`\pm `$ 0.07 m/sec in <sup>4</sup>He at $`\lambda `$ = $`\lambda _0`$(<sup>4</sup>He) = 0.062 Å<sup>-1</sup>. The velocity of sound drops to zero at the spinodal point that, according to the equation of state (Table II), is located at densities $`0.210\pm 0.001`$ Å<sup>-1</sup> and $`0.209\pm 0.001`$ Å<sup>-1</sup> for the 1D and T systems, respectively. The spatial structure of molecular hydrogen in the 1D array and inside the nanotube has also been analyzed by means of the two-body radial distribution function, $`g(z)`$, and its Fourier transform, the static structure factor $`S(k)`$. In Fig. 3, 1D results for $`g(z)`$ are reported at both the equilibrium density for the liquid phase, and in the solid-liquid transition region ($`\lambda `$ = 0.312 Å<sup>-1</sup>). The corresponding results for H<sub>2</sub> inside the nanotube are indistinguishable from the 1D results in the scale shown in Fig. 3. The solid $`g(z)`$ shows a strongly localized order around the equally spaced $`z`$-sites that decreases very slowly when $`z`$ increases. The result for $`g(z)`$ at $`\lambda `$ = $`\lambda _0`$ manifests the nature of a dense fluid with an appreciable structure that decreases slowly and shows residual ordering up to large $`z`$ distances. A comparison with previous DMC results on helium at the same density shows that the height of the <sup>4</sup>He peaks is approximately an 80 % of the corresponding ones in H<sub>2</sub>, but with a $`z`$-ordering that also survives up to large distances; the additional 20 % comes from the sizeable difference between the well depths of the respective interaction potentials. The significant differences in structure between the low and high density regimes are reflected even more clearly in the static structure factors. In Fig. 4, results for $`S(k)`$ corresponding to the same densities reported in Fig. 3 are shown. At high density, a characteristic result for a solid phase is obtained with a regular $`k`$-spacing according to the only periodicity allowed by the 1D geometry. At $`\lambda `$ = $`\lambda _0`$, $`S(k)`$ shows a first peak reflecting the localization observed in $`g(z)`$ (Fig. 3), and a subsequent very smoothed maximum as expected in a homogeneous liquid phase. In conclusion, we have studied the zero-temperature equation of state of molecular hydrogen in a 1D geometry and inside a narrow nanotube by means of the diffusion Monte Carlo method. The 1D calculation predicts the existence of a self-bound system with a binding energy of -4.8 K and a quasi-continuous liquid-solid transition at high density. The comparison with a real system, hydrogen in a nanotube, points to a close proximity between its properties and the ones of the 1D limit. The prediction of a liquid H<sub>2</sub> phase inside a (5,5) carbon nanotube is one of the main conclusions of the present work. The high one-dimensionality of this system would preclude a superfluid behavior but the use of wider nanotubes can provide a proper setup to observe the so-long desired superfluidity in molecular hydrogen. We expect the present work could encourage experimentalists to explore such a intriguing possibility. This research has been partially supported by DGES (Spain) Grants N<sup>0</sup> PB96-0170-C03-02 and PB98-0922, and DGR (Catalunya) Grant N<sup>0</sup> 1999SGR-00146. M. C. G. thanks the Spanish Ministry of Education and Culture (MEC) for a postgraduate contract. We also acknowledge the supercomputer facilities provided by the CEPBA.
warning/0004/quant-ph0004098.html	ar5iv	text	# Partitioning optical solitons for generating entangled light beams ## Abstract It is shown that bipartition of optical solitons can be used to generate entangled light beams. The achievable amount of entanglement can be substantially larger for $`N`$-bound solitons ($`N`$ $`=`$ $`2,3`$) than for the fundamental soliton ($`N`$ $`=`$ $`1`$). An analysis of the mode structure of the entangled beams shows that just $`N`$ modes are essentially entangled. In particular, partitioning of the fundamental soliton effectively produces $`2`$-mode squeezed light. Quantum coherent dynamics of atomic and photonic systems has been a subject of increasing interest stimulated by the rapidly developing field of quantum information processing and quantum computing . Intrinsic parallelism of quantum state evolution renders it possible to overcome the problem of exponential time required by a classical computer for solving complex problems such as integer factorization , discrete logarithm , and search . High sensitivity of quantum coherent dynamics to eavesdropping can be used for developing secure encryption protocols realized in fiber optical systems . One of the most remarkable features of quantum coherence is entanglement, on which, e.g., quantum teleportation is based . Entanglement can be regarded as being the nonclassical contribution to the overall correlation between two parts of a system (see, e.g., ). Typically, discrete systems that are built up by qubits have been considered. Concepts that use continuous quantum variables have recently attracted increasing publicity . An illustrative example is the teleportation of coherent states by means of entangled light of the squeezed-vacuum type . A way to generate entangled light is to use two fields each of which is prepared in some nonclassical state and combine them at a linear four-port device like a beam splitter or a coupler. In particular, the entangled squeezed light used in is produced by combining two light beams independently squeezed via $`\chi ^{(2)}`$ parametric processes. The possibility of realizing entangled pulse sources by combining two separately squeezed optical solitons are discussed in . Alternatively, a nonlinear coupling between two fields can prepare them in an entangled (continuous-variable) state. For example, the $`\chi ^{(2)}`$ process of parametric down-conversion has been extensively used as a source of entangled beams and pulses (see, e.g., ). In this letter we show that bipartition of (fundamental and $`N`$-bound) optical solitons formed in a medium with a Kerr-type $`\chi ^{(3)}`$ nonlinearity yields entangled light fields, because of the internal quantum correlations of the solitons. From classical optics it is well known that the Kerr nonlinearity can compensate for dispersion-assisted pulse spreading or diffraction-assisted beam broadening (see, e.g., ). In both cases, the evolution of the (complex) field amplitude $`a(x,t)`$ can be described by the nonlinear Schrödinger equation (NSE) $`i_ta(x,t)=\frac{1}{2}\omega ^{(2)}_{xx}a(x,t)+\chi \|a(x,t)\|^2a(x,t),`$ (1) ($`t`$, propagation variable; $`x`$, “transverse” coordinate; $`\omega ^{(2)}`$, second-order dispersion or diffraction constant; $`\chi `$, nonlinear-coupling constant). Bright temporal solitons can be formed either in focusing media with anomalous dispersion ($`\chi `$ $`<`$ $`0`$, $`\omega ^{(2)}`$ $`>`$ $`0`$) or in defocusing media with normal dispersion ($`\chi `$ $`>`$ $`0`$, $`\omega ^{(2)}`$ $`<`$ $`0`$). Spatial solitons always require a focusing nonlinearity. Instead of working in $`x`$-space, we can also turn to the $`\omega `$-space, $$\underset{¯}{a}(\omega ,t)=\frac{1}{\sqrt{2\pi }}_{\mathrm{}}^{\mathrm{}}𝑑xa(x,t)e^{i\omega x}.$$ (2) In what follows the coordinates in the $`x`$\- and $`\omega `$-domains are scaled by the initial soliton width $`x_0`$ and the spectral width $`\omega _0`$ $`=`$ $`1/x_0`$ respectively. Propagation distances are measured in dispersion (diffraction) lengths $`t_\mathrm{d}`$ $`=`$ $`x_0^2/\|\omega ^{(2)}\|`$. The NSE (1) can be solved by the inverse scattering method . In particular, the bell-shaped fundamental soliton is extremely stable and can propagate over distances far beyond the limit given by the dispersion- (diffraction-) assisted spreading. Scaling the power of the fundamental soliton by a factor of $`N^2`$, $`N`$ $`=`$ $`2,3,\mathrm{}`$, yields so-called $`N`$-bound solitons . The two classes of solutions of the NSE have been realized experimentally (see, e.g., ). In quantum theory, the complex field amplitude $`a(x,t)`$ and $`a^{}(x,t)`$ become field operator $`\widehat{a}(x,t)`$ and $`\widehat{a}^{}(x,t)`$, respectively, with $`[\widehat{a}(x,t),\widehat{a}^{}(x^{},t)]=\delta (xx^{}),`$ (3) and the NSE becomes an operator-valued equation . In order to solve the quantum mechanical problem, it is commonly assumed that the initial state is a multimode coherent state such that $`\widehat{a}`$ $``$ $`N\text{sech}(x/x_0)`$ is the initial shape of the classical $`N`$-bound-soliton solution. For not too large propagation distances the state then evolves into a multimode squeezed state of Gaussian type, which can be calculated numerically within the framework of appropriate discretization . The nonlinear soliton motion leads to internal quantum correlations, which have been studied both experimentally and theoretically ). From the results the question is suggested of whether or not it is possible to produce two entangled light beams by appropriate partitioning of quantum solitons. Let us consider bipartition as sketched in Fig. 1, which may be realized experimentally using spectral filtering or spatial filtering . Since the soliton quantum state remains pure during the propagation, the entanglement $`E`$ of the bipartite system is given by the von Neumann entropy of a subsystem i.e., $$E=S_1=\mathrm{Tr}_1(\widehat{\rho }_1\text{ln}\widehat{\rho }_1)=S_2=\mathrm{Tr}_2(\widehat{\rho }_2\text{ln}\widehat{\rho }_2),$$ (4) where $`\widehat{\rho }_i`$ is the density operator of the $`i`$th subsystem, and $`\mathrm{Tr}_i`$ means the trace with regard to the $`i`$th subsystem ($`i`$ $`=`$ $`1,2`$). In order to calculate the entropy of a subsystem prepared in a Gaussian state, let us consider the $`x`$-domain and assume that the subsystem extends over a region $`\mathrm{\Xi }`$ (the calculation in the $`\omega `$-domain is analogous). We introduce new bosonic operators $$\widehat{b}_k=_\mathrm{\Xi }𝑑x\left[\mu _k(x)\mathrm{\Delta }\widehat{a}(x)+\nu _k(x)\mathrm{\Delta }\widehat{a}^{}(x)\right]$$ (5) \[$`\mathrm{\Delta }\widehat{a}(x)`$ $`=`$ $`\widehat{a}(x)`$ $``$ $`\widehat{a}(x)`$\] and impose on the $`\mu _k(x)`$ and $`\nu _k(x)`$ the condition that $`\widehat{b}_k^{}\widehat{b}_k^{}=\sigma _k\delta _{kk^{}},\widehat{b}_k\widehat{b}_k^{}=0.`$ (6) For notational convenience, we have omitted the argument $`t`$. From Eqs. (5) and (6) it then follows that the parameters $`\sigma _k`$ and the functions $`\mu _k(x)`$ and $`\nu _k(x)`$ solve the equations $`[\begin{array}{c}\left(\sigma _k+\frac{1}{2}\right)\delta _{kk^{}}\\ 0\end{array}]={\displaystyle _\mathrm{\Xi }}𝑑x{\displaystyle _\mathrm{\Xi }}𝑑x^{}[\begin{array}{cc}\mu _k(x)& \nu _k(x)\\ \nu _k^{}(x)& \mu _k^{}(x)\end{array}]`$ (16) $`\times [\begin{array}{cc}C(x,x^{})^{}& B(x,x^{})\\ B(x,x^{})^{}& C(x,x^{})\end{array}][\begin{array}{c}\mu _k^{}^{}(x^{})\\ \nu _k^{}^{}(x^{})\end{array}],`$ where $`B(x,x^{})=\mathrm{\Delta }\widehat{a}(x)\mathrm{\Delta }\widehat{a}(x^{}),`$ (17) $`C(x,x^{})=\mathrm{\Delta }\widehat{a}^{}(x)\mathrm{\Delta }\widehat{a}(x^{})+\frac{1}{2}\delta (xx^{}).`$ (18) Note that the functions $`\mu _k(x)`$ and $`\nu _k(x)`$ must satisfy the conditions $`{\displaystyle _\mathrm{\Xi }}𝑑x\left[\mu _k(x)\mu _k^{}^{}(x)\nu _k(x)\nu _k^{}^{}(x)\right]=\delta _{kk^{}},`$ (19) $`{\displaystyle _\mathrm{\Xi }}𝑑x\left[\mu _k(x)\nu _k^{}(x)\nu _k(x)\mu _k^{}(x)\right]=0,`$ (20) because of the commutation relations $`[\widehat{b}_k,\widehat{b}_k^{}^{}]`$ $`=`$ $`\delta _{kk^{}}`$ and $`[\widehat{b}_k,\widehat{b}_k^{}]`$ $`=`$ $`0`$. With regard to the new bosonic operators, the Gaussian state (of the subsystem) is obviously a multimode thermal state, with $`\sigma _k`$ being the mean photon number of the $`k`$th mode. Hence, the entanglement $`E`$ is given by $$E=_kS_{\mathrm{th}}(\sigma _k),$$ (21) with $$S_{\mathrm{th}}(\sigma )=\text{ln}\left[\frac{(\sigma +1)^{\sigma +1}}{\sigma ^\sigma }\right]$$ (22) being the entropy of a single-mode thermal state. Note that when the $`\mathrm{\Xi }`$-region comprises the whole soliton prepared in a pure Gaussian state, then all the $`\sigma _k`$ vanish and thus $`E`$ $`=`$ $`0`$, as expected. Figure 2 presents examples of the entanglement obtained by symmetrical bipartition of the fundamental soliton and the 2- and 3-bound solitons in both the $`x`$-domain \[solid lines in Fig. 2$`(b)`$\] and the $`\omega `$-domain \[solid lines in Fig. 2$`(c)`$\] as a function of the propagation distance. The figure clearly shows that highly entangled parties can be realized by bipartition of optical solitons produced in Kerr-like media. The (numerical) calculations have shown that the maximum entanglement is typically observed for symmetrical bipartition (i.e., at $`x`$ $`=`$ $`0`$ or $`\omega `$ $`=`$ $`0`$). For the system considered in the figure, an asymmetrical bipartition of the $`3`$-bound soliton in the $`x`$-domain at $`\|x\|`$ $``$ $`0.4x_0`$ can improve the entanglement by $`5.3\%`$ for certain propagation distances. The entanglement realized by bipartition of the fundamental soliton increases monotonously with the propagation distance. The maxima and minima of the entanglement realized by bipartition in the $`x`$-domain of $`2`$\- and $`3`$-bound solitons follow quite exactly the periodic change between soliton compression and expansion \[compare Figs. 2$`(a)`$ and 2$`(b)`$\]. In the $`\omega `$-domain this correlation is less pronounced \[compare Figs. 2$`(a)`$ and 2$`(c)`$\]. It is worth noting that the amount of entanglement that can be achieved increases with the order parameter $`N`$. Formally, the sum in Eq. (21) runs over all modes. In fact, only a few modes substantially contribute to the entanglement, as it is seen from Figures 2 and 3 (the modes are numbered such that the single-mode contribution to the entanglement decreases with increasing mode index). In particular, it is seen that the number of relevant modes is just given by the soliton-order parameter $`N`$. A possible reason for this fact can be seen in the multicomponent structure of the classical $`N`$-soliton solution as discussed in . In summary, the results offer novel possibilities of using optical solitons as sources of entangled light beams. Although bipartition of solitons again yields multimode objects, only a few modes are involved in the entanglement. In particular, fundamental solitons are suited for generation of entangled light like $`2`$-mode squeezed light, which is an alternative to the standard $`\chi ^{(2)}`$ sources. ###### Acknowledgements. This work was supported by the Deutsche Forschungsgemeinschaft. We are grateful to Tomaš Opatrný for valuable discussions.
warning/0004/hep-th0004090.html	ar5iv	text	# 1 2T Physics ## 1 2T Physics Two-time (2T) physics has by now been shown to provide a reformulation of all possible one-time (1T) particle dynamics, including interactions with Yang-Mills, gravitational and other fields -. 2T physics has mainly been developed in the context of particles, but some advances have also been made with strings and p-branes , and some insights for M-theory have already emerged . In the case of particles, there exists a general worldline formulation with background fields , as well as a field theory formulation , both described in terms of fields that depend on $`d+2`$ coordinates $`X^M`$. The 2T point of view has been useful in bringing new insights into 1T physics, first by revealing previously unnoticed hidden symmetries in 1T dynamical systems, and second by providing a unification of classes of 1T dynamics that are different in 1T physics, but are regarded as being gauge equivalent to each other under the local gauge symmetries of a unique 2T system. It is important to realize that the Standard Model of particle physics and gravity (i.e. all the physics we can verify) can be rewritten as a 2T field theory, modulo a reformulation of the mass term in the Higgs potential. The 2T approach is not a naive extension of the number of timelike dimensions, as this would be disastrous for unitarity and causality, but it does imply the existence of a more subtle higher structure of dimensional unification, including an extra timelike and an extra spacelike dimension, plus gauge symmetries that tend to reduce the number of effective dimensions (worldline Sp$`(2,R)`$ or OSp$`\left(n\|2\right),`$ and local spacetime generalized bosonic/fermionic kappa symmetries). Thus, for a fixed set of 2T background fields (e.g. flat background) one can find many related interacting 1T systems by gauge fixing the underlying gauge symmetry. The non-trivial aspects of 2T physics is due to the fact that $`d`$ dimensions (1T) can be embedded in many ways inside $`d+2`$ dimensions (2T) (also such concepts extended to superspace). The 1T observers see the same 2T system as different dynamics from the point of view of the chosen $`d`$ dimensions, but all 1T observers are gauge equivalent from the point of view of the 2T observer. This gauge equivalence is translated to a kind of “duality” among the 1T dynamical systems that are derived from a given 2T system. The origins of the field theory formalism goes back to Dirac’s work in 1936 on conformal symmetry, but this approach historically was used exclusively as a reformulation of conformal symmetry in field theory -. An early approach on the worldline is also aimed to conformal symmetry . The unification of different dynamics in the form of 2T physics was not understood until the work of , which reached the Sp$`(2,R)`$ gauge theory formulation on the worldline by following a very different path and motivation (coming from 2T signals and dualities in M-theory, F-theory, S-theory, and influenced directly by the formalism in ). Also, in the field theory formalism, it was realized only recently that the 2T field equations are a reflection of the Sp$`(2,R)`$ gauge symmetries that underlie 2T physics on the worldline, and that the gauge fixing procedure which produce the various 1T dynamics in the worldline formulation can be carried out in the field theory formalism as well. Although the 2T unification of 1T systems can be examined in either the worldline or field theory forms, the worldline approach provides a better understanding of the underlying gauge symmetries at this stage, while the field theory formulation provides a familiar approach for interactions among fields. In this paper we will further extend the gauge symmetries of 2T physics on the worldline by further developing some concepts already introduced in that were necessary for the 2T formulation of supersymmetry. The gauge symmetry is associated with the spacetime SO$`(d,2)`$, some internal symmetry and an extended concept of bosonic/fermionic kappa supersymmetry. We will use the Sp$`(2,R)`$ gauge symmetry together with the extra gauge symmetry to show that particle dynamics in one of the Sp$`(2,R)`$ 1T gauges (the SO$`(d1,1)`$ relativistic particle gauge<sup>2</sup><sup>2</sup>2The same idea could be applied in other Sp$`(2,R)`$ 1T gauges, but we will refrain from discussing other Sp$`(2,R)`$ gauges in this paper.) can be cast either in terms of particle coordinates and momenta $`(x^\mu ,p^\mu )`$ or in terms of twistor variables à la Penrose or supertwistors à la Ferber . This result will be extended to the conformal superparticle cases for $`d=3,4,6`$ with $`N`$ supersymmetries using the supergroups OSp$`\left(N\|4\right),`$ SU$`(2,2\|N),`$ OSp$`\left(8^{}\|2N\right)`$ respectively. It will be shown that the (super)particle description $`(x,p,\theta )`$ is gauge equivalent to the (super)twistor description when the SO$`(d,2)`$ and Sp$`(2,R)`$ gauge symmetries are utilized. Furthermore, it will be shown that the quantization of the model can be carried out in terms of oscillators associated with the (super)twistors, and the unitarity established. The main motivation for developing this formalism came from trying to understand the quantum theory of a toy M-model introduced in . The methods presented here will be applied in to the case of the toy M-model based on the supergroup OSp$`\left(1\|64\right).`$ ## 2 1T superparticle Since our ultimate goal includes spacetime supersymmetry we will directly discuss the massless superparticle; the reader can delete the fermions to specialize to the purely bosonic case. The action for the superparticle in $`d`$ dimensions, with $`N`$ supersymmetries, is traditionally described by the following well known Lagrangian written in terms of position $`x^\mu \left(\tau \right)`$, momentum $`p^\mu \left(\tau \right)`$ and $`N`$ spacetime spinors $`\theta _\alpha ^a\left(\tau \right)`$ ($`\alpha `$ denotes the spinor and $`a`$=1,2,$`\mathrm{},N`$) $$\mathrm{\pounds }=\frac{1}{2A^{22}}\left(\dot{x}^\mu +\stackrel{~}{\theta }_a\gamma ^\mu _\tau \theta ^a\right)^2=\dot{x}p\frac{1}{2}A^{22}p^2+\stackrel{~}{\theta }_a\gamma p_\tau \theta ^a.$$ (1) We used the symbol $`A^{22}\left(\tau \right)`$ for the einbein because of its relation to the Sp$`(2,R)`$ gauge fields $`A^{ij}\left(\tau \right)`$ that we will see below. The general variation of the fields has the form (up to total derivatives) $`\delta L`$ $`=`$ $`_\tau \left(\delta x\delta \theta _a\gamma \theta ^a\right)p+2\delta \theta _a\left(p\gamma \right)\dot{\theta }^a`$ (3) $`+\delta p\left(\dot{x}^\mu +\stackrel{~}{\theta }_a\gamma ^\mu _\tau \theta ^aA^{22}p\right)\delta A^{22}{\displaystyle \frac{p^2}{2}}`$ For general $`d`$ and $`N,`$ using the appropriate spinor, this Lagrangian is symmetric ($`\delta L=0`$ up to total derivatives$`)`$ under local $`\tau `$-reparametrization, local kappa supersymmetry and global super Poincaré symmetry. In particular the global supersymmetry with parameters $`\epsilon _\alpha ^a`$ is given by $$\delta _\epsilon x^\mu =\delta _\epsilon \theta \gamma ^\mu \theta ,\delta _\epsilon \theta =\epsilon ,\delta _\epsilon p^\mu =0,\delta _\epsilon A^{22}=0,$$ (4) and the local kappa transformation with parameters $`\kappa _\alpha ^a\left(\tau \right)`$ is $$\delta _\kappa x^\mu =\delta _\kappa \theta \gamma ^\mu \theta ,\delta _\kappa \stackrel{~}{\theta }_a=\stackrel{~}{\kappa }_a\gamma p,\delta _\kappa p^\mu =0,\delta _\kappa A^{22}=4\stackrel{~}{\kappa }\dot{\theta }.$$ (5) For the special dimensions $`d=3,4,6`$ this Lagrangian is also invariant under dilations, conformal transformations and special superconformal transformations, such that the full global symmetry is given by the supergroups $$G=OSp\left(N\|4\right),SU(2,2\|N),OSp\left(8^{}\|N\right)$$ (6) respectively, as shown in for $`d=3,4`$ and in for $`d=6.`$ The conformal subgroups in these dimensions are SO$`(3,2)=Sp\left(4\right),`$ SO$`(4,2)=SU(2,2),`$ and SO$`(6,2)=Spin\left(8^{}\right)`$ respectively. For the purely bosonic case this Lagrangian has global conformal symmetry SO$`(d,2)`$ for any $`d.`$ ## 3 2T formulation The conformal symmetry SO$`(d,2)`$ is a giveaway for 2T physics. Indeed all of the above cases correspond to a special Sp$`(2,R)`$ gauge choice of a 2T formulation (the SO$`(d1,1)`$ relativistic particle gauge) in which the SO$`(d,2)`$ Lorentz symmetry in flat backgrounds in $`d+2`$ dimensions gets interpreted as the conformal symmetry in 1T (it acquires other interpretations in other Sp$`(2,R)`$ gauges). The 2T reformulation of the superparticle in $`d=3,4,6`$ dimensions requires $`d+2`$ coordinates $`X^M\left(\tau \right)`$ and momenta $`P^M\left(\tau \right),`$ and a supergroup element $`g\left(\tau \right)G`$ that contains fermions $`\mathrm{\Theta }_{\stackrel{~}{\alpha }}^a\left(\tau \right)`$ where $`\stackrel{~}{\alpha }`$ denotes the spinor in $`d+2`$ dimensions. This spinor has double the size of the spinor $`\theta _\alpha ^a\left(\tau \right)`$ in $`d`$ dimensions, which is of course necessary if the SO$`(d,2)`$ is to be realized linearly in the 2T formulation. Thus, compared to the 1T formulation there are extra degrees of freedom in $`X,P,\mathrm{\Theta }`$ and in the bosonic sectors in $`g\left(\tau \right).`$ If the covariant 2T formulation is to be equivalent to the 1T formulation there has to be various gauge symmetries and extended kappa supersymmetries to cut down the degrees of freedom to the correct set. Following this is beautifully achieved as follows. The 2T Lagrangian is $$L=\dot{X}_1X_2\frac{1}{2}A^{ij}X_iX_j\frac{1}{s}Str\left(\mathrm{\Gamma }_{MN}\dot{g}g^1\right)L^{MN},$$ (7) where $`X_i^M=(X^M,P^M)`$ is the Sp$`(2,R)`$ doublet, $`A^{ij}`$ is the Sp$`(2,R)`$ gauge potential, $`\mathrm{\Gamma }_M`$ are gamma matrices and $`\mathrm{\Gamma }_{MN}=\frac{1}{2}[\mathrm{\Gamma }_M,\mathrm{\Gamma }_N]`$ are the SO$`(d,2)`$ generators in the spinor representation of dimension $`s,`$ the Cartan connection $`\dot{g}g^1`$ projected in the direction of the SO$`(d,2)G`$ is coupled to the SO$`(d,2)`$ orbital angular momentum $`L^{MN}=\epsilon ^{ij}X_i^MX_j^N=X^MP^NX^NP^M`$ which is Sp$`(2,R)`$ gauge invariant. Although our interest in this paper is on the supergroups $`G`$ listed in (6) the discussion of local symmetries below applies also to any group or supergroup that contains SO$`(d,2)`$ as a subgroup. In fact, for the toy M-model in the case of interest is $`G=OSp\left(1\|64\right).`$ In particular, for the purely bosonic particle one can simply take $`G=SO(d,2).`$ The following is an improvement of the symmetry discussion given in . From the extensive discussions in - we already know that the Lagrangian above has local symmetry under Sp$`(2,R)`$. Beyond this, it obviously has global symmetry under $`G`$ for the transformation of $`g\left(\tau \right)`$ by right multiplication by a global group element $`g_RG`$ $$X_i^MX_i^M,A^{ij}A^{ij},ggg_R.$$ (8) Furthermore, it has local SO$`(d,2)`$ Lorentz symmetry with parameters $`\epsilon ^{MN}\left(\tau \right)`$ under left multiplication of $`g`$ in the spinor representation and transformation of $`X_i^M`$ in the vector representation $$\delta X_i^M=\epsilon ^{MN}X_{iN},\delta g=\frac{1}{4}\epsilon ^{MN}\left(\mathrm{\Gamma }_{MN}g\right),\delta A^{ij}=0.$$ (9) The time derivatives of $`\dot{\epsilon }^{MN}`$ produced by the two kinetic terms in (7) cancel each other. Moreover, for the cases in which there is another subgroup in $`G`$ (such as the SO$`\left(N\right),`$ SU$`\left(N\right),`$ Sp$`\left(N\right)`$ in (6)) with generators $`T_A`$ that satisfy Str$`\left(T_A\mathrm{\Gamma }_{MN}\right)=0,`$ there is a local symmetry with parameters $`\epsilon ^A\left(\tau \right)`$ under left multiplication of $`g`$ $$\delta g=\epsilon ^A\left(T_Ag\right),\delta X_i^M=0,\delta A^{ij}=0.$$ (10) The time derivative $`\dot{\epsilon }^A`$ as well other dependence on $`\epsilon ^A`$ drops because Str$`\left(T_A\mathrm{\Gamma }_{MN}\right)=0`$ and $`[T_A,\mathrm{\Gamma }_{MN}]=0.`$ Finally there is a local bosonic/fermionic extended kappa (super)symmetry under left multiplication of $`g`$ with infinitesimal coset elements $`K`$ of the form (take OSp$`\left(N\|2M\right)`$ as an example for the matrix notation) $$\delta g=Kg,K=\left(\begin{array}{cc}0& \xi \left(\tau \right)\\ \stackrel{~}{\xi }\left(\tau \right)& 0\end{array}\right),Str\left(\mathrm{\Gamma }_{MN}K\right)=0,Str\left(T_AK\right)=0,$$ (11) provided $`\delta A^{ij}`$ is non-zero as specified below, and $`\xi _{\stackrel{~}{\alpha }}^a\left(\tau \right)`$ has the form $$\xi _{\stackrel{~}{\alpha }}^a\left(\tau \right)=X_i^M\left(\mathrm{\Gamma }_M\kappa ^{ia}\right)_{\stackrel{~}{\alpha }},$$ (12) where the local $`\kappa _{\stackrel{~}{\alpha }}^{ia}\left(\tau \right)`$ are unrestricted local parameters. Under such a transformation we have $$\delta \mathrm{\pounds }=0+\frac{2}{s}L^{MN}Str\left([\mathrm{\Gamma }_{MN},K]_\tau gg^1\right)\frac{1}{2}\delta A^{ij}X_iX_j,$$ One must have $`L^{MN}\left(\mathrm{\Gamma }_{MN}\xi \right)`$ proportional to $`X_iX_j`$ so that $`\delta A^{ij}`$ can be chosen to cancel the contribution from the first term. Indeed, the general form in (12) has this property $$L^{MN}\left(\mathrm{\Gamma }_{MN}\xi \right)=\epsilon ^{kj}X_k^MX_j^NX_i^R\left(\mathrm{\Gamma }_{MN}\mathrm{\Gamma }_R\kappa ^i\right)=2\epsilon ^{kj}X_k^M\left(\mathrm{\Gamma }_M\kappa ^i\right)X_jX_i$$ (13) since the three gamma term in $`\mathrm{\Gamma }_{MN}\mathrm{\Gamma }_R=\mathrm{\Gamma }_{MNR}+\mathrm{\Gamma }_M\eta _{NR}\mathrm{\Gamma }_N\eta _{MR}`$ drops out due to the fact that the indices $`i,j,k`$ take only two possible values. Let us now specialize to a few cases of interest and use the gauge symmetries to cut down the degrees of freedom to those in $`d`$ dimensions given in the beginning of the previous section. * We start with the purely bosonic particle. Using the SO$`(d,2)`$ local symmetry we can choose $`g\left(\tau \right)=1`$ for all $`\tau .`$ The Lagrangian (7) is now expressed only in terms of $`X^M,P^M.`$ We work in the basis $`X^M=(X^+^{},X^{^{}},x^\mu )`$ where $`X^\pm ^{}=\left(X^0^{}\pm X^1^{}\right)/\sqrt{2}`$ are lightcone type coordinates for the extra two dimensions (similarly for $`P^M)`$. Using the Sp$`(2,R)`$ local symmetry we can choose $`X^+^{}\left(\tau \right)=1`$ and $`P^+^{}\left(\tau \right)=0,`$ and solve the two constraints $`X^2=0,`$ $`XP=0,`$ to give $`X^{^{}}=x^2/2`$ and $`P^{^{}}=xp.`$ The Lagrangian reduces to the bosonic particle Lagrangian in (1) without the fermions. When $`g=1`$ the global SO$`(d,2)`$ that acts on the right of $`g`$ must be compensated by a global transformation on the left of $`g`$ which also acts on $`X_i^M.`$ Thus the left/right transformations on $`g`$ get coupled and become the 2T spacetime global SO$`(d,2)`$ Lorentz symmetry. When the Sp$`(2,R)`$ gauges are also fixed, this global symmetry acts non-linearly on the remaining variables $`(x^\mu ,p^\mu )`$ and is then interpreted as the conformal symmetry of the massless particle. * Next we discuss the superparticle in $`d=3,4,6`$ as in . Using the local SO$`(d,2)`$ and local internal symmetries the form of $`g,`$ for $`G`$ given in (6), can be gauge fixed to $`gt=\mathrm{exp}\left(fermioniccoset\right)`$ parametrized by the spinor $`\mathrm{\Theta }_{\stackrel{~}{\alpha }}^a\left(\tau \right).`$ This was the starting point in . Using the kappa gauge $`\mathrm{\Gamma }^+^{}\mathrm{\Theta }=0,`$ and Sp$`(2,R)`$ gauges $`X^+^{}=1,P^+^{}=0`$ it was shown in that the 2T Lagrangian reduces precisely to the superparticle Lagrangian in (1) for $`d=3,4,6.`$ The global symmetry $`G`$ becomes the non-linearly realized superconformal symmetry of the massless superparticle in these special dimensions. * For other supergroups the use of the local symmetries can never reduce the degrees of freedom to only the superspace variables $`(x,p,\theta ).`$ Generally there are more bosonic degrees of freedom. It was speculated in that the extra degrees of freedom can be associated with collective coordinates that describe D-brane degrees of freedom in the particle limit. This is because the superalgebra has central extensions with relations among the charges such that BPS conditions are satisfied, as is the case for D-branes. This point will be further elaborated in for the case of OSp$`\left(1\|64\right).`$ ## 4 Twistors, supertwistors, oscillators To make the presentation as explicit as possible we will concentrate on the $`d=3`$ superparticle, rewritten in the $`d+2=5`$ two-time formalism. Hence we will take $`G=OSp\left(N\|4\right)`$ where Sp$`\left(4\right)=SO(3,2).`$ However, we will begin more generally with $`G=`$OSp$`\left(M/2N\right)`$ where Sp$`\left(2N\right)`$ has an SO$`(d,2)`$ subgroup whose spinor representation has dimension $`s=2N`$. For example for $`d=11`$ in the toy M-model, $`SO(11,2)`$ has a spinor with 64 components, and therefore we would consider OSp$`\left(1\|64\right).`$ In the basis $`X^M=(X^+^{},X^{^{}},X^\mu )`$, we define the following SO$`(d,2)`$ gamma matrices $$\mathrm{\Gamma }^\pm ^{}=\pm \tau ^\pm \times 1,\mathrm{\Gamma }^\mu =\tau _3\times \gamma ^\mu ,\{\mathrm{\Gamma }^M,\mathrm{\Gamma }^N\}=2\eta ^{MN}.$$ (14) For the case of $`d=3`$ we take the following explicit form for the SO$`(2,1)`$ gamma matrices $$\gamma ^\mu =(i\sigma _2,\sigma _1,\sigma _3)$$ (15) The group $`G=`$OSp$`\left(M/2N\right)`$ is characterized by $`gG`$ of the form $$g^1=\widehat{C}g^{st}\widehat{C}^1$$ (16) where $`g^{st}`$ is the supertranspose of $`g`$ and $`\widehat{C}`$ is the metric of OSp$`\left(M/2N\right)`$. In supermatrix notation we can write $$\widehat{C}=\left(\begin{array}{cc}1_M\hfill & 0\hfill \\ 0\hfill & C_{2N}\hfill \end{array}\right),g=\left(\begin{array}{cc}\alpha \hfill & \stackrel{~}{f}_1\hfill \\ f_2\hfill & A\hfill \end{array}\right),g^{st}=\left(\begin{array}{cc}\alpha ^t& f_2^t\\ \stackrel{~}{f}_1^t& A^t\end{array}\right),g^1=\left(\begin{array}{cc}\alpha ^t& \stackrel{~}{f}_2\\ f_1& \stackrel{~}{A}\end{array}\right)$$ (17) with the definitions $$\stackrel{~}{f}=f^tC_{2N}^1,\stackrel{~}{A}=C_{2N}A^tC_{2N}^1.$$ (18) Here $`C_{2N}`$ is the antisymmetric metric for Sp$`\left(2N\right)`$, with properties $$C_{2N}^t=C_{2N}=C_{2N}^1,\left(C_{2N}\right)^2=1_{2N}.$$ (19) Note the extra minus sign in $`\stackrel{~}{f}_1^t`$ in the definition of supertranspose $`g^{st}`$. This is necessary so that the supertranspose operation has the property $`\left(g_1g_2\right)^{st}=g_2^{st}g_1^{st}.`$ The necessity of the extra minus sign in $`\stackrel{~}{f}^t`$ can be traced to the extra minus sign that arises from anticommuting two fermions under transposition. Thus the parameters are constrained by the relation $`1`$ $`=`$ $`g^1g=\left(\begin{array}{cc}\alpha ^t& \stackrel{~}{f}_2\\ f_1& \stackrel{~}{A}\end{array}\right)\left(\begin{array}{cc}\alpha & \stackrel{~}{f}_1\\ f_2& A\end{array}\right)`$ (24) $`\alpha ^t\alpha +\stackrel{~}{f}_2f_2`$ $`=`$ $`1_N,\alpha ^t\stackrel{~}{f}_1+\stackrel{~}{f}_2A=0,f_1\alpha +\stackrel{~}{A}f_2=0,f_1\stackrel{~}{f}_1+\stackrel{~}{A}A=1.`$ (25) For infinitesimal group parameters the solution of these constraints are $$\delta \alpha ^t=\delta \alpha ,\delta f_1=\delta f_2\delta f,\delta \stackrel{~}{A}=\delta A$$ (26) Thus, one may write $$g=\mathrm{exp}\left(\begin{array}{cc}\delta \alpha & \delta \stackrel{~}{f}\\ \delta f& \delta A\end{array}\right)$$ (27) where $`\delta f`$ is any $`M\times \left(2N\right)`$ fermionic matrix, $`\delta \alpha `$ is any $`M\times M`$ antisymmetric matrix and $`\delta A`$ is any $`\left(2N\right)\times \left(2N\right)`$ symplectic matrix that satisfies $`C_{2N}\delta A^tC_{2N}^1=\delta A`$. ### 4.1 The global current Using Noether’s theorem one finds that the global OSp$`\left(N\|2M\right)`$ current of our model is $$J=g^1Lg.$$ (28) where $`L=\frac{1}{2}L_{MN}\mathrm{\Gamma }^{MN}`$. At the classical level the current satisfies $$J^2=g^1Lgg^1Lg=g^1L^2g0.$$ (29) The vanishing is because of the Sp$`(2,R)`$ constraints $`X^2=P^2=XP=0`$ at the classical level. As discussed elsewhere the treatment of constraints at the quantum level modifies this result such that the Casimir operators Str$`\left(J^n\right)`$ computed from the global currents are all non-zero but fixed at definite constants that define a specific representation. In the supermatrix notation given above the current is given by $$J=\left(\begin{array}{cc}\alpha ^t& \stackrel{~}{f}_2\\ f_1& \stackrel{~}{A}\end{array}\right)\left(\begin{array}{cc}0\hfill & 0\hfill \\ 0\hfill & L\hfill \end{array}\right)\left(\begin{array}{cc}\alpha & \stackrel{~}{f}_1\\ f_2& A\end{array}\right)=\left(\begin{array}{cc}\stackrel{~}{f}_2Lf_2\hfill & \stackrel{~}{f}_2LA\hfill \\ \stackrel{~}{A}Lf_2\hfill & \stackrel{~}{A}LA\hfill \end{array}\right)$$ (30) The parameters $`\alpha `$ have dropped out because of the internal gauge symmetry. We will now choose further gauges that completely eliminate all degrees of freedom from $`X^M,P^M,`$ shifting the dynamics to the degrees of freedom in $`g\left(\tau \right).`$ We will show that the remaining degrees of freedom are supertwistors. Using the local SO$`(d,2)`$ and Sp$`(2,R)`$ gauge symmetries we can map the two vectors $`X^M\left(\tau \right),P^M\left(\tau \right)`$ to the constants $`X^+^{}=1,`$ $`P^{}=1`$ with all other components zero $`X^M=\delta _+^{}^M,`$ $`P^M=\delta _{}^M.`$ These satisfy the constraints $`X^2=P^2=XP=0.`$ In this gauge we have $$L=\frac{1}{2}\mathrm{\Gamma }^{MN}L_{MN}=\mathrm{\Gamma }^{^{}}\mathrm{\Gamma }^+=\tau ^{}\times \gamma ^+=\left(\begin{array}{cc}0\hfill & 0\hfill \\ \gamma ^+\hfill & 0\hfill \end{array}\right),$$ (31) Plugging this form into the Lagrangian we find that the remaining degrees of freedom are described by the following dynamics $$\mathrm{\pounds }Tr\left(\stackrel{~}{A}\tau ^{}\gamma ^+_\tau A+\stackrel{~}{f}_2\tau ^{}\gamma ^+_\tau f_2\right).$$ (32) The current given above satisfies the commutation rules of OSp$`\left(N\|2M\right)`$ if we take the basic commutation rules $$\{f_\alpha ^i,f_\beta ^j\}=i\delta ^{ij}\widehat{L}_{\alpha \beta },[A_\alpha ^\gamma ,A_\beta ^\delta ]=i\widehat{L}_{\alpha \beta }\left(C^1\right)^{\gamma \delta }$$ (33) where $`\widehat{L}`$ is defined by the relation $`L\widehat{L}L=L,`$ and is given by $$\widehat{L}=\tau ^+\times \gamma ^{}.$$ (34) In the basis of gamma matrices we have chosen, the charge conjugation matrix is $`C_4=\tau _1\times C_2`$. Using this $`C_4`$ we also parametrize the 4-column $`f_2`$ and the 4$`\times `$4 matrix $`A`$ in terms of 2 dimensional blocks $$f_2^i=\left(\begin{array}{c}\xi ^i\\ \chi ^i\end{array}\right),A=\left(\begin{array}{cc}a\hfill & b\hfill \\ c\hfill & d\hfill \end{array}\right),\stackrel{~}{f}_{2i}=\left(\begin{array}{cc}\stackrel{~}{\chi }_i\hfill & \stackrel{~}{\xi }_i\hfill \end{array}\right),\stackrel{~}{A}=\left(\begin{array}{cc}\stackrel{~}{d}\hfill & \stackrel{~}{b}\hfill \\ \stackrel{~}{c}\hfill & \stackrel{~}{a}\hfill \end{array}\right)$$ (35) where $`\stackrel{~}{\xi }_i=\xi _i^TC_2^1`$ and $`\stackrel{~}{a}=C_2a^TC_2^1`$, etc. We may now compute the Lagrangian and the global OSp$`\left(8\|4\right)`$ currents more explicitly $`\mathrm{\pounds }`$ $``$ $`{\displaystyle \frac{1}{2}}\stackrel{~}{\xi }^i\gamma ^+_\tau \xi ^j+{\displaystyle \frac{1}{2}}Tr\left(\stackrel{~}{b}\gamma ^+_\tau a+\stackrel{~}{a}\gamma ^+_\tau b\right)`$ (36) $`J`$ $`=`$ $`\left(\begin{array}{ccc}\stackrel{~}{\xi }^i\gamma ^+\xi ^j\hfill & \stackrel{~}{\xi }^i\gamma ^+a\hfill & \stackrel{~}{\xi }^i\gamma ^+b\hfill \\ \stackrel{~}{b}\gamma ^+\xi ^j\hfill & \stackrel{~}{b}\gamma ^+a\hfill & \stackrel{~}{b}\gamma ^+b\hfill \\ \stackrel{~}{a}\gamma ^+\xi ^j\hfill & \stackrel{~}{a}\gamma ^+a\hfill & \stackrel{~}{a}\gamma ^+b\hfill \end{array}\right)=\left(\begin{array}{ccc}I^{ij}& \stackrel{~}{Q}^i& \stackrel{~}{S}^i\\ S^j& J& K\\ Q^i& P& \stackrel{~}{J}\end{array}\right)`$ (43) The canonical pairs are easily determined from the Lagrangian. The currents are quadratic in the canonical pairs. They are reminiscent of the oscillator formalism for supergroups and indeed we will make this connection much clearer in the following paragraphs. The interpretation of the various components of the currents are as follows. The 8$`\times 8`$ block $`I^{ij}=\stackrel{~}{\xi }^i\gamma ^+\xi ^j`$ corresponds to the generators of SO$`\left(8\right),`$ The 2$`\times 8`$ blocks $`Q^i=\stackrel{~}{a}\gamma ^+\xi ^i`$ and $`S^i=\stackrel{~}{b}\gamma ^+\xi ^i`$ represent 8 supercharges and 8 superconformal charges respectively. The 4$`\times 4`$ block including $`J,P,K`$ represents Sp$`\left(4\right)=SO(3,2)`$, with $`J,P,K`$ corresponding to Lorentz transformations and dilatations, translations and conformal transformations respectively. Note that the parameters $`\chi ,c,d`$ have dropped out from the Lagrangian and the physical global currents. This is because of the extended local kappa supersymmetry and bosonic gauge symmetries. The remaining parameters $`\xi ,a,b`$ are constrained by $`gg^1=1`$, or $`f_2\stackrel{~}{f}_2+A\stackrel{~}{A}=1`$ or $$\xi ^i\stackrel{~}{\xi }^i+a\stackrel{~}{b}+b\stackrel{~}{a}=0.$$ (44) This also guarantees $`J^2=0.`$ Furthermore, because $$\gamma ^+=\left(\begin{array}{cc}0& 1\\ 0& 0\end{array}\right)$$ (45) is a projection operator, the currents depend only on a few of the parameters in the 2$`\times 2`$ matrices $`a,b`$ and column matrix $`\xi .`$ This too is because of the extended gauge symmetry. To further simplify the expressions we define the entries in each matrix $$\xi ^i=\left(\begin{array}{c}\lambda ^i\\ \theta ^i\end{array}\right),a=\left(\begin{array}{cc}a_3& a_4\\ a_1& a_2\end{array}\right),b=\left(\begin{array}{cc}b_3& b_4\\ b_1& b_2\end{array}\right),$$ We find that the currents are given by the unrestricted parameters $`(\theta ^i,a_1,a_2,b_1,b_2)`$ corresponding to the second row of the matrices above, while the remaining parameters drop out from the Lagrangian and currents. We will write $`a_\alpha =(a_1,a_2),`$ $`a^\alpha =(a_2,a_1),`$ raising/lowering indices $`\alpha `$ by using the Levi-Civita symbol $`\epsilon ^{\alpha \beta }`$ which is the charge conjugation matrix. In this notation the Lagrangian and currents take the form $$\mathrm{\pounds }\frac{1}{2}\theta ^i\dot{\theta }^i+b^\alpha \dot{a}_\alpha ,J=\left(\begin{array}{ccc}\theta ^i\theta ^j\hfill & \theta ^ia_\beta \hfill & \theta ^ib_\beta \hfill \\ b^\alpha \theta ^j\hfill & b^\alpha a_\beta \hfill & b^\alpha b_\beta \hfill \\ a^\alpha \theta ^j\hfill & a^\alpha a_\beta \hfill & a^\alpha b_\beta \hfill \end{array}\right)=\left(\begin{array}{ccc}I^{ij}& \stackrel{~}{Q}^i& \stackrel{~}{S}^i\\ S^j& J& K\\ Q^i& P& \stackrel{~}{J}\end{array}\right)$$ (46) The basic non-zero commutation rules $$\{\theta ^i,\theta ^j\}=i\delta ^{ij},[a_\alpha ,b^\beta ]=i\delta _\alpha ^\beta $$ (47) determine the algebra of the conserved charges $`J`$. It is easily seen that the charges form the OSp$`\left(8\|4\right)`$ superalgebra. In fact the form of this construction is identical to the oscillator construction of superalgebras . In the present case the real fermions $`\theta ^i`$ are SO$`\left(8\right)`$ spinors and the canonical conjugates $`a_\alpha `$, or $`b^\alpha `$ are real Sp$`\left(2\right)`$ spinors, where Sp$`\left(2\right)=SO(2,1)`$ is the Lorentz subgroup of the conformal group Sp$`\left(4\right)=SO(3,2)`$. In the usual oscillator construction of one chooses the maximal compact subgroup of OSp$`\left(8\|4\right),`$ which is SU$`\left(4\right)\times SU\left(2\right)\times \left(U\left(1\right)\right)^2,`$ and takes a complex quartet of fermionic oscillators in the fundamental representation of SU$`\left(4\right)\times U\left(1\right)`$ and a complex doublet of bosonic oscillators in the compact subgroup SU$`\left(2\right)\times U\left(1\right),`$ plus their hermitian conjugates. The relation between these two is just a change of basis such that the Sp$`\left(4\right)`$ quartet given in the Sp$`\left(2\right)`$ basis $`(a_1,a_2,b_1,b_2)`$ is rexpressed in terms of the oscillators $`A_\alpha =\left(a_\alpha +ib^\alpha \right)/\sqrt{2}`$ as a Sp$`\left(4\right)`$ quartet in an SU$`\left(2\right)\times U\left(1\right)`$ basis $`(A_1,A_2,A_2^{},A_1^{}).`$ Similarly the real SO$`\left(8\right)`$ spinor $`\theta ^i`$ is re-expressed in the complex SU$`\left(4\right)\times U\left(1\right)`$ basis in terms of fermionic oscillators. This construction proves the unitarity of the representation, and identifies the super conformal particle with the super doubleton representation of OSp$`\left(8\|4\right).`$ It is evident that the same general arguments hold for OSp$`\left(N\|4\right).`$ Returning to the Sp$`\left(2\right)=SO(3,2)`$ spinors $`a_\alpha `$ or $`b^\alpha `$, we can interpret them as the twistor representation of the particle canonical coordinates $`x^\mu ,p^\mu `$ à la Penrose as follows. To avoid complications due to quantum ordering, we will discuss only the classical version of this interpretation. Also, we concentrate only on the purely bosonic Sp$`\left(4\right)`$ to avoid complications with the supergroup. Starting with the generators $`P_\beta ^\alpha =a^\alpha a_\beta `$ which form a traceless 2$`\times 2`$ matrix, we identify the momentum as $$a^\alpha a_\beta =\left(\gamma ^\mu \right)_\beta ^\alpha p_\mu ,$$ (48) Evidently it describes a massless particle, since $`p^2\delta _\beta ^\alpha =a^\alpha a_\gamma a^\gamma a_\beta =0.`$ The coordinate $`x^\mu `$ is defined by the following relation between the two spinors $`a_\alpha ,b_\beta `$ $$b_\alpha =x^\mu \left(\gamma _\mu \right)_\alpha ^\beta a_\beta .$$ (49) To show that this is indeed the case, consider the conformal generator $`K^\mu `$ written in the form $`K_\beta ^\alpha =K^\mu \left(\gamma ^\mu \right)_\beta ^\alpha =b^\alpha b_\beta `$ and insert the expression for $`b_\alpha `$ $`K^\mu `$ $`=`$ $`{\displaystyle \frac{1}{2}}b^\alpha \left(\gamma ^\mu \right)_\alpha ^\beta b_\beta ={\displaystyle \frac{1}{2}}x^\lambda a^\alpha \left(\gamma _\lambda \gamma ^\mu \gamma _\nu \right)_\alpha ^\beta a_\beta x^\nu `$ (50) $`=`$ $`{\displaystyle \frac{1}{2}}x^\lambda p_\sigma Tr\left(\gamma ^\sigma \gamma _\lambda \gamma ^\mu \gamma _\nu \right)_\alpha ^\beta x^\nu ={\displaystyle \frac{1}{2}}x^2p_\mu xpx^\mu .`$ (51) This result is the well known expression for the conformal generator (avoiding quantum ordering, or fermionic contributions). Similarly we compute the dimension operator $`D=\frac{1}{2}b^\alpha a_\alpha `$ and the Lorentz generator $`J^{\mu \nu }=b^\alpha \left(\gamma ^{\mu \nu }\right)_\alpha ^\beta a_\beta `$ $`D`$ $`=`$ $`{\displaystyle \frac{1}{2}}b^\alpha a_\alpha ={\displaystyle \frac{1}{2}}x^\mu a^\alpha \left(\gamma _\mu \right)_\alpha ^\beta a_\beta ={\displaystyle \frac{1}{2}}x^\mu p^\nu Tr\left(\gamma _\mu \gamma _\nu \right)=xp`$ (52) $`J^{\mu \nu }`$ $`=`$ $`b^\alpha \left(\gamma ^{\mu \nu }\right)_\alpha ^\beta a_\beta =x^\lambda a^\alpha \left(\gamma _\lambda \gamma ^{\mu \nu }\right)_\alpha ^\beta a_\beta =x^\lambda p_\sigma Tr\left(\gamma _\lambda \gamma ^{\mu \nu }\gamma _\sigma \right)=x_\mu p_\nu x_\nu p_\mu `$ (53) These are the correct expressions for the massless bosonic particle. This makes it evident that the basis we have defined in terms of the super variables $`(\theta ^i,a_\alpha ,b_\alpha )`$ is indeed the twistor basis for the conformal superparticle. In the presence of fermions the relation between the spinors $`b_\alpha ,a_\alpha `$ is considerably more complicated such that the correct generators $`K^\mu ,J^{\mu \nu }`$ emerge, as given in . The generalized relation including the fermions is obtained by applying the kappa and other gauge transformations that take the model from the fixed gauge $`g\left(\tau \right)=t\left(\tau \right)`$ used in to the fixed gauge of the twistor formalism described in this paper. ## 5 Discussion and generalizations From the previous section, it is evident that the approach is applied in a straightforward manner in $`d=4,6`$ using the supergroups $`SU(2,2\|N)`$ and OSp$`\left(8^{}\|N\right)`$ respectively. The main point is the shifting of the degrees of freedom from $`X,P`$ to $`g`$ or vice-versa by using the gauge symmetries Sp$`(2,R)`$ and SO$`(d,2).`$ When $`X,P`$ are eliminated, the remaining gauge symmetries reduce further the degrees of freedom in $`g`$ to the super twistor variables for those dimensions. In particular, for $`d=4`$ we find agreement with . Similarly, twistors can be gauge transformed to coordinate representation to describe the same system. Twistors were invented as a means of describing conformal systems covariantly in linear realizations. We have now shown that they are related to another linear realization, namely Lorentz transformations in 2T physics. In this paper we have mainly concentrated on a specific Sp$`(2,R)`$ gauge choice, namely $`X^+^{}=1,`$ $`P^+^{}=0`$ which relates to the SO$`(d1,1)`$ covariant massless relativistic particle. This is a particular embedding of $`d`$ dimensions inside $`d+2`$ dimensions, and corresponds to a particular 1T physics interpretation of the 2T theory. As we know from - there are many other 1T embeddings with different 1T physics interpretations. Such other Sp$`(2,R)`$ gauges may now be combined with the present techniques of shifting particle variables to twistor-like variables embedded in $`g,`$ and thus find new twistor-like realizations of 1T dynamical systems, as well as establish duality-type relations among them. The reader will notice that the spacetime and internal subgroups of OSp$`\left(8\|4\right),`$ SU$`(2,2\|4),`$ OSp$`\left(8^{}\|4\right)`$ were treated in an asymmetric manner in the coupling introduced in (7). Since these supergroups describe the supersymmetries in AdS$`{}_{7}{}^{}\times S_4,`$ AdS$`{}_{5}{}^{}\times S_5,`$ AdS$`{}_{4}{}^{}\times S_7,`$ respectively, one may wonder if there is a more symmetric treatment of the AdS$`\times `$S spaces that would apply to these cases. In fact, in addition to the spacetime $`X^M,P^M`$ phase space we may introduce internal $`Y^I,K^I`$ phase space, define the internal angular momentum $`L^{IJ}=Y^IK^JY^JK^I`$ and couple it to the supergroup Cartan connection in the same manner as (7). When both $`L^{MN}`$ and $`L^{IJ}`$ have non-zero coupling it is possible to maintain a kappa-type local supersymmetry as well as the Sp$`(2,R)`$ local symmetry coupled to all the coordinates $`(X^M,Y^I)`$ and momenta $`(P^M,K^I).`$ This variation leads to more interesting and intricate 2T models. In fact it is even possible to couple to any a subgroup of the internal group SO$`\left(N\right),`$ SU$`\left(N\right),`$ Sp$`\left(N\right)`$ that appear in (6). If only a subgroup $`H`$ is gauged, then only the corresponding degrees of freedom can be removed. The remaining coset plays the role of harmonic superspace recently discussed in . As mentioned earlier, a motivation for developing these techniques was the study of the quantum system for the toy M-model that will be discussed elsewhere . Having shown that the approach works and establishes connections among previously better understood systems, it may now be used to explore new systems.

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