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# Algebraic Models of Hadron Structure II. Strange Baryons
## 1 Introduction
In the last few years there has been renewed interest in hadron spectroscopy. Especially the development of dedicated experimental facilities to probe the structure of hadrons in the nonperturbative region of QCD with far greater precision than before has generated a considerable amount of experimental and theoretical activity . This has stimulated us to reexamine hadron spectroscopy in a novel approach in which both internal (spin-flavor-color) and space degrees of freedom of hadrons are treated algebraically. The new ingredient is the introduction of a space symmetry or spectrum generating algebra for the radial excitations which for mesons was taken as $`U(4)`$ and for baryons as $`U(7)`$ . The algebraic approach unifies the harmonic oscillator quark model, $`U(4)U(3)`$ for mesons and $`U(7)U(6)`$ for baryons, with collective string-like models of hadrons.
In the first paper of this series we have introduced $`U(7)`$ to study the properties of nonstrange baryons, such as the mass spectrum, electromagnetic couplings and strong decays . In this article we extend these studies to hyperons and present a systematic study of both nonstrange and strange baryons in the framework of a collective string-like $`qqq`$ model in which the orbital excitations are treated as rotations and vibrations of the strings. The algebraic structure of the model enables us to obtain transparent results (mass formula, selection rules and decay widths for strong and electromagnetic couplings) that can be used to analyze and interpret the experimental data, and look for evidence of the existence of unconventional (i.e. non $`qqq`$) configurations of quarks and gluons, such as hybrid quark-gluon states $`qqq`$-$`g`$ or multiquark meson-baryon bound states $`qqq`$-$`q\overline{q}`$.
In particular, we discuss the mass spectrum (Sects. 3-4), the strong (Sects. 5-6) and electromagnetic (Sects. 7-8) decay widths. We do this in a framework in which spin-flavor symmetry is broken in a diagonal way in the masses (i.e. through a dynamic symmetry). This asssumption appears to be sufficient to describe most observables. The breaking of spin-flavor symmetry in hyperon decays can be investigated using a procedure similar to that in Ref. for nonstrange baryons. The results of such a study will be published separately.
## 2 Algebraic models of baryons
We consider baryons to be built of three constituent parts which are characterized by both internal and spatial degrees of freedom.
### 2.1 Degrees of freedom
The internal degrees of freedom of these three parts are taken to be: flavor-triplet $`u,d,s`$ (for the light quark flavors), spin-doublet $`S=1/2`$, and color-triplet. The internal algebraic structure of the constituent parts consists of the usual spin-flavor and color algebras
$`𝒢_\mathrm{i}`$ $`=`$ $`SU_{\mathrm{sf}}(6)SU_\mathrm{c}(3).`$ (2.1)
In we discussed various algebraic models of baryons. These models share a common spin-flavor structure (see Eq. (2.1), but differ in their treatment of radial excitations. Here we consider a collective string-like model with the configuration depicted in Fig. 1. The relevant degrees of freedom for the relative motion of the three constituent parts of this configuration are provided by the relative Jacobi coordinates which we choose as
$`\stackrel{}{\rho }`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}(\stackrel{}{r}_1\stackrel{}{r}_2),`$
$`\stackrel{}{\lambda }`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{m_1^2+m_2^2+(m_1+m_2)^2}}}[m_1\stackrel{}{r}_1+m_2\stackrel{}{r}_2(m_1+m_2)\stackrel{}{r}_3].`$ (2.2)
Here $`m_i`$ and $`\stackrel{}{r}_i`$ ($`i=1,2,3`$) denote the mass and coordinate of the $`i`$-th constituent. When two of the constituents have equal mass ($`m_1=m_2`$), the above choice reduces to
$`\stackrel{}{\rho }`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}(\stackrel{}{r}_1\stackrel{}{r}_2),`$
$`\stackrel{}{\lambda }`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{6}}}(\stackrel{}{r}_1+\stackrel{}{r}_22\stackrel{}{r}_3).`$ (2.3)
Since the quark masses satisfy to a good approximation $`m_u=m_dm_s`$, the Jacobi coordinates of Eq. (2.3) are relevant for all baryons be it with strangeness $`S=0`$, $`1`$, $`2`$ or $`3`$. Instead of a formulation in terms of coordinates and momenta, we use the method of bosonic quantization in which we introduce a dipole boson with $`L^P=1^{}`$ for each independent relative coordinate, and an auxiliary scalar boson with $`L^P=0^+`$
$`b_{\rho ,m}^{},b_{\lambda ,m}^{},s^{}(m=1,0,1).`$ (2.4)
The scalar boson does not represent an independent degree of freedom, but is added under the restriction that the total number of bosons $`N`$ is conserved. This procedure leads to a compact spectrum generating algebra for the radial (or orbital) excitations
$`𝒢_r`$ $`=`$ $`U(7).`$ (2.5)
The $`U(7)`$ algebra enlarges the $`U(6)`$ algebra of the harmonic oscillator quark model , but still describes the dynamics of two vectors. For a system of interacting bosons the model space is spanned by the symmetric irreducible representation $`[N]`$ of $`U(7)`$. This representation contains all oscillator shells with $`n=n_\rho +n_\lambda =0,1,2,\mathrm{},N`$. The value of $`N`$ determines the size of the model space and, in view of confinement, is expected to be large.
### 2.2 Basis states
The full algebraic structure is obtained by combining the spatial part $`𝒢_\mathrm{r}`$ of Eq. (2.5) with the internal spin-flavor-color part $`𝒢_\mathrm{i}`$ of Eq. (2.1)
$`𝒢`$ $`=`$ $`𝒢_\mathrm{r}SU_{\mathrm{sf}}(6)SU_\mathrm{c}(3).`$ (2.6)
The spatial part of the baryon wave function has to be combined with the spin-flavor and color part, in such a way that the total wave function is antisymmetric. Since the color part of the wave function is antisymmetric (color singlet), the remaining part (space-spin-flavor) has to be symmetric. A convenient set of basis states is provided by the case of three identical constituents, for which the spatial and spin-flavor parts of the baryon wave function are in addition labeled by their transformation properties under the permutation group $`S_3`$: $`t=S`$ for the symmetric, $`t=A`$ for the antisymmetric and $`t=M`$ for the two-dimensional $`(M_\rho ,M_\lambda )`$ mixed symmetry representation.
A set of basis states for the spin-flavor part is provided by the decomposition of $`SU_{\mathrm{sf}}(6)`$ into its flavor and spin parts
$`|\begin{array}{ccccccccccc}SU_{\mathrm{sf}}(6)& & SU_\mathrm{f}(3)& & SU_\mathrm{s}(2)& & SU_\mathrm{I}(2)& & U_\mathrm{Y}(1)& & SU_\mathrm{s}(2)\\ & & & & & & & & & & \\ [f_1f_2f_3]& & [g_1g_2]& & S& & I& & Y& & \end{array}.`$ (2.10)
Here $`[f_1f_2f_3]`$ and $`[g_1g_2]`$ represent the Young tableaux, $`S`$ denotes the spin, $`I`$ the isospin and $`Y`$ the hypercharge. The representations of the spin-flavor groups are often labeled by their dimensions (rather than by their Young tableaux)
$`\text{dim}[f_1f_2f_3]`$ $`=`$ $`{\displaystyle \frac{(f_1f_2+1)(f_1f_3+2)(f_2f_3+1)(f_1+5)!(f_2+4)!(f_3+3)!}{3!4!5!(f_1+2)!(f_2+1)!f_3!}},`$
$`\text{dim}[g_1g_2]`$ $`=`$ $`{\displaystyle \frac{1}{2}}(g_1g_2+1)(g_1+2)(g_2+1),`$
$`\text{dim}[S]`$ $`=`$ $`2S+1.`$ (2.11)
For three constituent parts the allowed values of $`[f_1f_2f_3]`$ are $`[300]`$ ($`t=S`$), $`[210]`$ ($`t=M`$) and $`[111]`$ ($`t=A`$) with dimensions 56, 70 and 20, respectively. The flavor part is characterized by $`[g_1g_2]=[30]`$, $`[21]`$ or $`[00]`$ with dimensions 10 (decuplet), 8 (octet) or 1 (singlet), respectively. In the notation of the flavor wave functions are labeled by $`(p,q)=(g_1g_2,g_2)`$. Finally, the total spin of three spin-1/2 objects is $`S=3/2`$ or $`S=1/2`$. The decomposition of representations of $`SU_{\mathrm{sf}}(6)`$ into those of $`SU_\mathrm{f}(3)SU_\mathrm{s}(2)`$ is the standard one
$`S[56]`$ $``$ $`{}_{}{}^{2}8^410,`$
$`M[70]`$ $``$ $`{}_{}{}^{2}8^48^210^21,`$
$`A[20]`$ $``$ $`{}_{}{}^{2}8^41,`$ (2.12)
where we have denoted the irreducible representations by their dimensions. Each flavor multiplet consists of families of baryons which are characterized by their isospin $`I`$ and hypercharge $`Y`$ (see Table I). The electric charge is given by the Gell-Mann and Nishijima relation
$`Q`$ $`=`$ $`I_3+{\displaystyle \frac{Y}{2}}.`$ (2.13)
In Appendix A we present the explicit form for the spin and flavor wave functions in the convention that we have used in this paper.
Since the space-spin-flavor wave function is symmetric ($`t=S`$), the symmetry of the spatial wave function under $`S_3`$ has to be the same as that of the spin-flavor part. Hence it is convenient to label the spatial wave functions by the basis states of a dynamical symmetry of $`U(7)`$ that preserves the $`S_3`$ permutation symmetry. We choose the chain that correponds to the problem of three particles in a common harmonic oscillator potential
$`|\begin{array}{ccccccccccc}U(7)& & U(6)& & 𝒮𝒰(3)& & SU(2)& & 𝒮𝒪(3)& & SO(2)\\ N& ,& n& ,& (n_1,n_2)& ,& F& ,& L& ,& m_F\end{array}.`$ (2.16)
In this decomposition, the behavior in three-dimensional coordinate space $`𝒮𝒰(3)𝒮𝒪(3)`$ is separated from that in index space $`SU(2)SO(2)`$. The allowed values of the quantum numbers can be obtained from the branching rules. For the decomposition of $`U(6)`$ we use the complementarity relationship between the groups $`𝒮𝒰(3)`$ and $`SU(2)`$ within the symmetric irreducible representation $`U(6)`$. As a consequence, the labels of $`𝒮𝒰(3)`$ are determined by those of $`SU(2)`$. The branching rules are
$`n`$ $`=`$ $`0,1,\mathrm{},N,`$
$`F`$ $`=`$ $`n,n2,\mathrm{},1\text{ or }0,`$
$`(n_1,n_2)`$ $`=`$ $`({\displaystyle \frac{n+F}{2}},{\displaystyle \frac{nF}{2}}),`$
$`m_F`$ $`=`$ $`F,F+2,\mathrm{},F.`$ (2.17)
The reduction from the coupled harmonic oscillator group to the rotation group $`𝒮𝒰(3)𝒮𝒪(3)`$ is given by
$`K`$ $`=`$ $`\text{min}\{\lambda ,\mu \},\text{min}\{\lambda ,\mu \}2,\mathrm{},1\text{ or }0,`$
$`K`$ $`=`$ $`0:L=\text{max}\{\lambda ,\mu \},\text{max}\{\lambda ,\mu \}2,\mathrm{},1\text{ or }0,`$
$`K`$ $`>`$ $`0:L=K,K+1,\mathrm{},K+\text{max}\{\lambda ,\mu \}.`$ (2.18)
Here $`(\lambda ,\mu )=(n_1n_2,n_2)=(F,(nF)/2)`$. The label $`K`$ is an extra label that has to be introduced to label the states uniquely . The $`SO(2)`$ group in Eq. (2.16) is related to the permutation symmetry . The states with good $`S_3`$ symmetry are given by the linear combinations
$`|\psi _1`$ $`=`$ $`{\displaystyle \frac{i}{\sqrt{2(1+\delta _{m_F,0})}}}\left[|\varphi _{m_F}|\varphi _{m_F}\right],`$
$`|\psi _2`$ $`=`$ $`{\displaystyle \frac{(1)^\nu }{\sqrt{2(1+\delta _{m_F,0})}}}\left[|\varphi _{m_F}+|\varphi _{m_F}\right].`$ (2.19)
Here we have introduced the label $`\nu `$ by $`m_F=\nu `$ (mod 3). The wave functions $`|\psi _1`$ ($`|\psi _2`$) transform for $`\nu =0`$ as $`t=A`$ ($`S`$), and for $`\nu =1,2`$ as $`t=M_\rho `$ ($`M_\lambda `$). Summarizing, the basis states are characterized uniquely by
$`|N,n,F,m_F,K,L_t^P,`$ (2.20)
where $`P`$ is the parity of the basis states $`P=()^n`$. Finally, the quark orbital angular momentum $`L`$ is coupled with the spin $`S`$ to the total angular momentum $`J`$ of the baryon. In Appendix B we present the space-spin-flavor baryon wave functions with $`S_3`$ symmetry.
## 3 Mass operator
The mass operator depends both on the spatial and the internal degrees of freedom. For the spatial part we adopt a collective model of the nucleon in which the baryons are interpreted as rotational and vibrational excitations of the string configuration of Fig. 1. For two identical constituent parts (as is the case for strange baryons) the vibrations are described by
$`\widehat{M}_{\mathrm{vib}}^2`$ $`=`$ $`AP_1^{}P_1+BP_2^{}P_2+CP_3^{}P_3+D(P_1^{}P_2+P_2^{}P_1),`$ (3.1)
with
$`P_1^{}`$ $`=`$ $`R^2s^{}s^{}b_\rho ^{}b_\rho ^{}b_\lambda ^{}b_\lambda ^{},`$
$`P_2^{}`$ $`=`$ $`(\mathrm{cos}\beta )^2b_\rho ^{}b_\rho ^{}(\mathrm{sin}\beta )^2b_\lambda ^{}b_\lambda ^{},`$
$`P_3^{}`$ $`=`$ $`b_\rho ^{}b_\lambda ^{}.`$ (3.2)
Here $`R`$ is related to the hyperspherical radius $`\sqrt{\rho ^2+\lambda ^2}`$, and $`\beta `$ corresponds to the hyperspherical angle $`\mathrm{tan}\beta =\rho /\lambda `$ with $`\rho =|\stackrel{}{\rho }|`$ and $`\lambda =|\stackrel{}{\lambda }|`$. The mass operator in this case is $`S_2`$ invariant. In the limit of a large model space ($`N\mathrm{}`$) the mass operator of Eqs. (3.1)-(3.2) reduces to leading order in $`N`$ to a harmonic form, and its eigenvalues are given by
$`M_{\mathrm{vib}}^2`$ $`=`$ $`\kappa _1n_u+\kappa _2n_v+\kappa _3n_w,`$ (3.3)
Here $`\kappa _1`$, $`\kappa _2`$ are the eigenvalues of the $`2\times 2`$ symmetric matrix
$`\left(\begin{array}{cc}4ANR^2& 2DN\mathrm{sin}(2\beta )R^2/\sqrt{1+R^2}\\ 2DN\mathrm{sin}(2\beta )R^2/\sqrt{1+R^2}& BN\mathrm{sin}^2(2\beta )R^2/(1+R^2)\end{array}\right).`$ (3.6)
and $`\kappa _3=CNR^2/(1+R^2)`$. The vibrational quantum numbers $`n_u`$, $`n_v`$ and $`n_w`$ denote the number of quanta in the symmetric stretching (or breathing mode), antisymmetric stretching and bending vibrations of the strings, respectively (see Fig. 3 of ). For three identical constituents we obtain the $`S_3`$-invariant mass operator of from Eqs. (3.1)-(3.2) by taking $`D=0`$, $`B=C`$ and $`\beta =\pi /4`$, which leads to $`\kappa _1=4ANR^2`$ and $`\kappa _2=\kappa _3=BNR^2/(1+R^2)`$. In the analysis of the mass spectrum of strange baryons, to be presented below, the $`S_3`$ symmetry of the mass operator is only broken dynamically in the spin-flavor part. Therefore, the baryon wave functions still have good $`S_3`$ symmetry, and the vibrational part of the mass operator only depends on $`\kappa _1`$ and $`\kappa _2`$ for all baryons
$`M_{\mathrm{vib}}^2`$ $`=`$ $`\kappa _1v_1+\kappa _2v_2.`$ (3.7)
Here $`v_1=n_u`$ and $`v_2=n_v+n_w`$ are the vibrational quantum numbers corresponding to the symmetric stretching vibration along the direction of the strings (breathing mode), and two degenerate bending vibrations of the strings. The spectrum consists of a series of vibrational excitations characterized by the labels $`(v_1,v_2)`$, and a tower of rotational excitations built on top of each vibration. The occurrence of linear Regge trajectories suggests to add a term linear in $`L`$ to the mass operator
$`M_{\mathrm{space}}^2`$ $`=`$ $`\kappa _1v_1+\kappa _2v_2+\alpha L.`$ (3.8)
In the application to nonstrange baryons the Roper N$`(1440)`$, the $`\mathrm{\Delta }(1600)`$ and the $`\mathrm{\Delta }(1900)`$ resonances were assigned to the symmetric stretching vibration $`(v_1,v_2)=(1,0)`$, and the N$`(1710)`$ resonance to the $`(v_1,v_2)=(0,1)`$ vibration. The remaining resonances were interpreted as rotational excitations.
For the spin-flavor part of the mass operator we use the Gürsey-Radicati form
$`\widehat{M}_{\mathrm{sf}}^2`$ $`=`$ $`a\left[\widehat{C}_2(SU_{\mathrm{sf}}(6))45\right]+b\left[\widehat{C}_2(SU_\mathrm{f}(3))9\right]+c\left[\widehat{C}_2(SU_\mathrm{s}(2)){\displaystyle \frac{3}{4}}\right]`$ (3.9)
$`+d\left[\widehat{C}_1(U_\mathrm{Y}(1))1\right]+e\left[\widehat{C}_2(U_\mathrm{Y}(1))1\right]+f\left[\widehat{C}_2(SU_\mathrm{I}(2)){\displaystyle \frac{3}{4}}\right].`$
The eigenvalues of the Casimir operators in the basis states of Eq. (2.10) are
$`\widehat{C}_2(SU_{\mathrm{sf}}(6))`$ $`=`$ $`2\left[f_1(f_1+5)+f_2(f_2+3)+f_3(f_3+1){\displaystyle \frac{1}{6}}(f_1+f_2+f_3)^2\right],`$
$`\widehat{C}_2(SU_\mathrm{f}(3))`$ $`=`$ $`{\displaystyle \frac{3}{2}}\left[g_1(g_1+2)+g_2^2{\displaystyle \frac{1}{3}}(g_1+g_2)^2\right],`$
$`\widehat{C}_2(SU_\mathrm{s}(2))`$ $`=`$ $`S(S+1),`$
$`\widehat{C}_1(U_\mathrm{Y}(1))`$ $`=`$ $`Y,`$
$`\widehat{C}_2(U_\mathrm{Y}(1))`$ $`=`$ $`Y^2,`$
$`\widehat{C}_2(SU_\mathrm{I}(2))`$ $`=`$ $`I(I+1).`$ (3.10)
We have defined the operators such that each of the terms vanishes for the ground state of the nucleon. The spin term represents spin-spin interactions, the flavor term denotes the flavor dependence of the interactions, and the $`SU_{\mathrm{sf}}(6)`$ term, which according to Eq. (3.10) depends on the permutation symmetry of the wave functions, represents ‘signature dependent’ interactions. These signature dependent (or exchange) interactions were extensively investigated years ago within the framework of Regge theory . The last two terms represent the isospin and hypercharge dependence of the masses. We do not consider here interaction terms that mix the space and internal degrees of freedom.
## 4 Comparison with experimental mass spectrum
In this section we analyze simultaneously the experimental mass spectrum of strange and nonstrange baryons in terms of the mass formula
$`M^2`$ $`=`$ $`M_0^2+\kappa _1v_1+\kappa _2v_2+\alpha L`$ (4.1)
$`+a\left[2f_1(f_1+5)+2f_2(f_2+3)+2f_3(f_3+1){\displaystyle \frac{1}{3}}(f_1+f_2+f_3)^245\right]`$
$`+b\left[{\displaystyle \frac{3}{2}}\left(g_1(g_1+2)+g_2^2{\displaystyle \frac{1}{3}}(g_1+g_2)^2\right)9\right]+c\left[S(S+1){\displaystyle \frac{3}{4}}\right]`$
$`+d\left[Y1\right]+e\left[Y^21\right]+f\left[I(I+1){\displaystyle \frac{3}{4}}\right].`$
The coefficient $`M_0^2`$ is determined by the nucleon mass $`M_0^2=0.882`$ GeV<sup>2</sup>. The remaining nine coefficients are obtained in a simultaneous fit to the three and four star resonances of Tables III and IV which have been assigned as octet and decuplet states. We find a good overall fit for 48 resonances with an r.m.s. deviation of $`\delta =33`$ MeV. The values of the parameters are given in Table II. In the last column we show for comparison the parameters that were obtained in a fit to 25 nucleon and delta resonances with a r.m.s. deviation to $`\delta =39`$ MeV . In comparison with Table II of the $`\mathrm{\Delta }(1900)S_{31}`$ was left out, since it has been downgraded from a three to a two star resonance . Since for nonstrange resonances $`Y=1`$, the $`d`$ and $`e`$ terms in Eq. (4.1) do not contribute. The flavor and isospin dependent terms that determine the mass splitting between the nucleon and $`\mathrm{\Delta }`$ resonances can be combined into a single $`b`$ term with strength $`b+\frac{1}{3}f=0.031`$ GeV<sup>2</sup>, very close to the fitted value of $`0.030`$ GeV<sup>2</sup> for nonstrange baryons. Thus, the parameter values determined in the present simultaneous study of both strange and nonstrange resonances are almost the same as those found for nonstrange resonances.
Tables III and IV and Figs. 36 show that the mass formula of Eq. (4.1) provides a good overall description of both positive and negative baryon resonances belonging to the $`N`$, $`\mathrm{\Delta }`$, $`\mathrm{\Sigma }`$, $`\mathrm{\Lambda }`$, $`\mathrm{\Xi }`$ and $`\mathrm{\Omega }`$ families. There is no need for an additional energy shift for the positive parity states and another one for the negative parity states, as in the relativized quark model .
### 4.1 Octet and decuplet resonances
The results are presented in Fig. 2 for the ground state baryon octet with $`J^P=1/2^+`$ and the baryon decuplet with $`J^P=3/2^+`$. In Tables III and IV we show a comparison with all three and four star resonances. In our calculation we have assigned the $`N(1440)`$, $`\mathrm{\Delta }(1600)`$, $`\mathrm{\Sigma }(1660)`$ and $`\mathrm{\Lambda }(1600)`$ resonances to the vibration characterized by $`(v_1,v_2)=(1,0)`$, and the $`N(1710)`$, $`\mathrm{\Sigma }(1940)`$ and $`\mathrm{\Lambda }(1810)`$ resonances to the $`(v_1,v_2)=(0,1)`$ vibration. The remaining resonances are assigned as rotational members of the ground band with $`(v_1,v_2)=(0,0)`$.
We have followed the quark model assignments of Table 13.4 of , with the exception of the $`\mathrm{\Sigma }(1750)S_{11}`$ resonance which we have assigned as $`{}_{}{}^{2}8_{1/2}^{}[70,1^{}]`$, the lowest $`S_{11}`$ state with a mass of 1711 MeV. In our calculation the lowest four $`S_{11}`$ $`\mathrm{\Sigma }`$ states occur at 1711, 1755, 1822 and 1974 MeV (for the assignments we refer to Tables V and VI; the second state belongs to the decuplet). In the nucleon and $`\mathrm{\Lambda }`$ families the Roper resonance lies below the first excited negative parity resonance. We expect the same to be true for the $`\mathrm{\Sigma }`$ hyperons. With our assignment, $`\mathrm{\Sigma }(1750)`$ is the octet partner of $`N(1535)`$ and $`\mathrm{\Lambda }(1670)`$, which is also supported by their $`\eta `$ decay properties . In the two star $`\mathrm{\Sigma }(1620)`$ resonance has been assigned as the $`{}_{}{}^{2}8_{1/2}^{}[70,1^{}]`$ state, and the $`\mathrm{\Sigma }(1750)`$ resonance instead as the $`{}_{}{}^{4}8_{1/2}^{}[70,1^{}]`$ state. Our assignment of $`\mathrm{\Sigma }(1750)`$ coincides with that of . In the relatived quark model there are three low-lying $`S_{11}`$ $`\mathrm{\Sigma }`$ resonances at 1630, 1675 and 1695 MeV . The first one was associated with the two star $`\mathrm{\Sigma }(1620)`$ resonance, and the next one with the $`\mathrm{\Sigma }(1750)`$ resonance.
The $`\mathrm{\Sigma }(1940)D_{13}`$ resonance was not assigned in Table 13.4 of , whereas in it was tentatively assigned as $`{}_{}{}^{4}8_{3/2}^{}[70,1^{}]`$, the octet partner of $`N(1700)`$. In our calculation the lowest four $`D_{13}`$ $`\mathrm{\Sigma }`$ states are the spin-orbit partners of the $`S_{11}`$ states at 1711, 1755, 1822 and 1974 MeV, the first one of which has been associated with the $`\mathrm{\Sigma }(1670)`$ resonance. We have assigned the $`\mathrm{\Sigma }(1940)`$ resonance as a member of the $`(v_1,v_2)=(0,1)`$ vibrational band with $`{}_{}{}^{2}8_{3/2}^{}[56,1^{}]`$ which occurs at 1974 MeV. This assignment is supported by its strong decay properties (see Sect. 6). In the relativized quark model there are three low-lying $`D_{13}`$ $`\mathrm{\Sigma }`$ states at 1655, 1750 and 1755 MeV , of which the first two were associated with the $`\mathrm{\Sigma }(1670)`$ and $`\mathrm{\Sigma }(1940)`$ resonances.
### 4.2 Singlet resonances
There are three states which show a deviation of about 100 MeV or more from the data: the $`\mathrm{\Lambda }^{}(1405)`$, $`\mathrm{\Lambda }^{}(1520)`$ and $`\mathrm{\Lambda }^{}(2100)`$ resonances are overpredicted by 236, 121 and 97 MeV, respectively. These three resonances are assigned as singlet states in Table IV (and were not included in the fitting procedure). An additional energy shift for the singlet states (without effecting the masses of the octet and decuplet states) can be obtained by adding to the mass formula of Eq. (4.1) a term that only acts on the singlet states
$`M^2M^2+\mathrm{\Delta }M^2\delta _{g_1,0}\delta _{g_2,0}.`$ (4.2)
This corresponds to a shift in the singlet masses of $`M\sqrt{1+(\mathrm{\Delta }M^2/M^2)}M[1+(\mathrm{\Delta }M^2/2M^2)]`$. Since spin-orbit partners are shifted by the same amount, the mass splitting of 115 MeV between $`\mathrm{\Lambda }^{}(1405)`$ and $`\mathrm{\Lambda }^{}(1520)`$ cannot be reproduced by this mechanism. In principle, this splitting can be obtained from a spin-orbit interaction. However, in the rest of the baryon spectra there is no evidence for such a large spin-orbit coupling. This problem is common to $`qqq`$ models of baryons (e.g. the constituent quark model with chromodynamics, either in its nonrelativistic or its relativized form , and the chiral constituent quark model all overpredict the $`\mathrm{\Lambda }^{}(1405)`$ mass). Another explanation for the mass splitting between $`\mathrm{\Lambda }^{}(1520)`$ and $`\mathrm{\Lambda }^{}(1405)`$ is the proximity of the $`\mathrm{\Lambda }^{}(1405)`$ resonance to the $`N\overline{K}`$ threshold. The inclusion of the coupling to the $`N\overline{K}`$ and $`\mathrm{\Sigma }\pi `$ decay channels produces a downward shift of the $`qqq`$ state toward or even below the $`N\overline{K}`$ threshold . Such an interpretation is supported by the strong and electromagnetic couplings (see Sects. 6 and 8). In a chiral meson-baryon Lagrangian approach with an effective coupled-channel potential the $`\mathrm{\Lambda }^{}(1405)`$ resonance emerges as a quasi-bound state of $`N\overline{K}`$ .
### 4.3 Missing resonances
In Tables III and IV we presented the model states that could be associated with a three or four star resonance. In Tables V, VI and VII we show the masses of all low-lying octet, decuplet and singlet baryons. Since in the present approach no spin-orbit coupling has been taken into account, the states are grouped into multiplets labeled by $`L`$, $`S`$ and $`|LS|JL+S`$. The multiplets for which at least one of its members has been associated with a three or four star resonance in Tables III or IV are labeled by . Tentative assignments of one or two star resonances are indicated by . As in any $`qqq`$ model of baryons there are many more calculated states than have been observed. The lowest socalled ‘missing’ resonances of the octet are associated with the $`{}_{}{}^{2}8_{J}^{}[20,1^+]`$ state. Their calculated mass is given by 1713, 1849, 1826 and 1957 MeV for the $`N`$, $`\mathrm{\Sigma }`$, $`\mathrm{\Lambda }`$ and $`\mathrm{\Xi }`$ resonances. The search for these ‘missing’ resonances is important in order to verify the assignments of the resonances and to distinguish between different models of baryons, such as three quark $`qqq`$ vs. quark-diquark $`qqq`$ models which have less missing states because of the smaller number of degrees of freedom.
In a recent three-channel multi-resonance amplitude analysis by the Zagreb group evidence was found for the existence of a third low-lying $`P_{11}`$ state at $`1740\pm 11`$ MeV. The first two $`P_{11}`$ states at $`1439\pm 19`$ MeV and $`1729\pm 16`$ MeV correspond to the $`N(1440)`$ and $`N(1710)`$ resonances of the PDG . These $`P_{11}`$ states were associated with the states at 1540, 1770 and 1880 MeV in the relativized quark model . In the present calculation, they occur at 1444, 1683 and 1713 MeV, in good agreement with the analysis of the Zagreb group.
A recent analysis of new data on kaon photoproduction has shown evidence for a $`D_{13}`$ resonance at 1895 MeV . In the present calculation, there are several possible assignments (see Table V). The lowest state that can be assigned to this new resonance is a vibrational excitation $`(v_1,v_2)=(0,1)`$ with $`{}_{}{}^{2}8_{3/2}^{}[56,1^{}]`$ and mass 1847 MeV. This state belongs to the same vibrational band as the $`N(1710)`$ resonance, and is the octet partner of $`\mathrm{\Sigma }(1940)`$. Another possible assignment is as a member of the ground state band $`(v_1,v_2)=(0,0)`$ with $`{}_{}{}^{2}8_{3/2}^{}[70,2^{}]`$ and mass 1874 MeV. However, this state is completely decoupled in strong and electromagnetic decays. Finally, there is a state that belongs to the same vibrational band $`(v_1,v_2)=(1,0)`$ as the $`N(1440)`$ Roper resonance with $`{}_{}{}^{2}8_{3/2}^{}[70,1^{}]`$ and mass 1909 MeV. As far as the mass is concerned all three assignments are possible. The strong couplings for these states provide a more sensitive tool to determine the most likely assignment (see Sect. 6). In the relativized quark model a $`D_{13}`$ state has been predicted at 1960 MeV .
## 5 Strong couplings
Strong couplings provide an important test of baryon wave functions, and can be used to distinguish between different models of baryon structure. Here we consider strong decays of baryons by the emission of a pseudoscalar meson
$`BB^{}+M.`$ (5.1)
Several forms have been suggested for the form of the operator inducing the strong transition . We use here the simple form
$`_s`$ $`=`$ $`{\displaystyle \frac{1}{(2\pi )^{3/2}(2k_0)^{1/2}}}{\displaystyle \underset{j=1}{\overset{3}{}}}X_j^M\left[2g(\stackrel{}{s}_j\stackrel{}{k})\text{e}^{i\stackrel{}{k}\stackrel{}{r}_j}+h\stackrel{}{s}_j(\stackrel{}{p}_j\text{e}^{i\stackrel{}{k}\stackrel{}{r}_j}+\text{e}^{i\stackrel{}{k}\stackrel{}{r}_j}\stackrel{}{p}_j)\right],`$ (5.2)
where $`\stackrel{}{r}_j`$, $`\stackrel{}{p}_j`$ and $`\stackrel{}{s}_j`$ are the coordinate, momentum and spin of the $`j`$-th constituent, respectively; $`k_0`$ is the meson energy and $`\stackrel{}{k}=k\widehat{z}`$ denotes the momentum carried by the outgoing meson. The coefficients $`g`$ and $`h`$ denote the strength of the two terms in the transition operator of Eq. (5.2). The flavor operator $`X_j^M`$ corresponds to the emission of an elementary meson by the $`j`$-th constituent: $`q_jq_j^{}+M`$ (see Figure 7).
Using the symmetry of the wave functions, transforming to Jacobi coordinates, integrating over the baryon center of mass coordinate, and adopting the rest frame of the initial baryon, the operator of Eq. (5.2) reduces to
$`_s`$ $`=`$ $`{\displaystyle \frac{1}{(2\pi )^{3/2}(2k_0)^{1/2}}}\mathrm{\hspace{0.17em}6}X_3^M\left[(gk{\displaystyle \frac{1}{6}}hk)s_{3,z}\widehat{U}hs_{3,z}\widehat{T}_z{\displaystyle \frac{1}{2}}h(s_{3,+}\widehat{T}_{}+s_{3,}\widehat{T}_+)\right],`$ (5.3)
with
$`\widehat{U}`$ $`=`$ $`\text{e}^{ik\sqrt{\frac{2}{3}}\lambda _z},`$
$`\widehat{T}_m`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(\sqrt{{\displaystyle \frac{2}{3}}}p_{\lambda ,m}\text{e}^{ik\sqrt{\frac{2}{3}}\lambda _z}+\text{e}^{ik\sqrt{\frac{2}{3}}\lambda _z}\sqrt{{\displaystyle \frac{2}{3}}}p_{\lambda ,m}\right).`$ (5.4)
The calculation of the matrix elements of $`_s`$ can be done in configuration space ($`\stackrel{}{\rho }`$, $`\stackrel{}{\lambda }`$) or in momentum space ($`\stackrel{}{p}_\rho `$, $`\stackrel{}{p}_\lambda `$). The mapping onto the algebraic space of $`U(7)`$ is a convenient way to carry out the calculations, much in the same way as the mapping of coordinates and momenta onto creation and annihilation operators in the harmonic oscillator space. The operators $`\widehat{U}`$ and $`\widehat{T}_m`$ can be expressed algebraically by first making the replacement $`\stackrel{}{p}_\lambda /m_3ik_0\stackrel{}{\lambda }`$ and then mapping the coordinates onto the algebraic operators, $`\sqrt{2/3}\lambda _m\beta \widehat{D}_{\lambda ,m}/X_D`$ . The result is
$`\widehat{U}`$ $`=`$ $`\text{e}^{ik\beta \widehat{D}_{\lambda ,z}/X_D},`$
$`\widehat{T}_m`$ $`=`$ $`{\displaystyle \frac{im_3k_0\beta }{2X_D}}\left(\widehat{D}_{\lambda ,m}\text{e}^{ik\beta \widehat{D}_{\lambda ,z}/X_D}+\text{e}^{ik\beta \widehat{D}_{\lambda ,z}/X_D}\widehat{D}_{\lambda ,m}\right).`$ (5.5)
The dipole operator $`\widehat{D}_{\lambda ,m}`$ is a generator of $`U(7)`$ and $`X_D`$ is its normalization, as discussed in . The spatial matrix elements of $`\widehat{U}`$ and $`\widehat{T}_m`$ are obtained in the collective model of baryons by folding with a distribution function $`g(\beta )`$ of charge and magnetization over the entire volume
$`g(\beta )`$ $`=`$ $`\beta ^2\text{e}^{\beta /a}/2a^3.`$ (5.6)
All spatial matrix elements can be expressed in terms of the collective form factors
$`(k)`$ $`=`$ $`{\displaystyle 𝑑\beta g(\beta )\psi ^{}|\widehat{U}|\psi },`$
$`𝒢_m(k)`$ $`=`$ $`{\displaystyle 𝑑\beta g(\beta )\psi ^{}|\widehat{T}_m|\psi }.`$ (5.7)
Here $`|\psi `$ and $`|\psi ^{}`$ denote the spatial wave functions of the initial and final baryons.
For strong decays in which the initial baryon $`B`$ has angular momentum $`\stackrel{}{J}=\stackrel{}{L}+\stackrel{}{S}`$ and in which the final baryon $`B^{}`$ has $`L^{}=0`$ and thus $`J^{}=S^{}`$, the helicity amplitudes in the collective model are then given by
$`A_\nu (k)`$ $`=`$ $`{\displaystyle 𝑑\beta g(\beta )\mathrm{\Psi }^{}(0,S^{},S^{},\nu )|_s|\mathrm{\Psi }(L,S,J,\nu )},`$ (5.8)
Here $`|\mathrm{\Psi }(L,S,J,M_J)`$ and $`|\mathrm{\Psi }^{}(0,S^{},S^{},\nu )`$ denote the (space-spin-flavor) angular momentum coupled wave functions of the initial and final baryons, respectively (see Appendix B). The final baryon is a ground state baryon belonging either to the octet $`{}_{}{}^{2}8_{1/2}^{}[56,0^+]_{(0,0);0}`$ or the decuplet $`{}_{}{}^{4}10_{3/2}^{}[56,0^+]_{(0,0);0}`$. The helicity amplitudes can be expressed in terms of a spatial matrix element and a spin-flavor matrix element
$`A_\nu (k)`$ $`=`$ $`{\displaystyle \frac{1}{(2\pi )^{3/2}(2k_0)^{1/2}}}[L,0,S,\nu |J,\nu \zeta _0Z_0(k)+{\displaystyle \frac{1}{2}}L,1,S,\nu 1|J,\nu \zeta _+Z_{}(k)`$ (5.9)
$`+{\displaystyle \frac{1}{2}}L,1,S,\nu +1|J,\nu \zeta _{}Z_+(k)],`$
with
$`Z_0(k)`$ $`=`$ $`6(gk{\displaystyle \frac{1}{6}}hk)(k)6h𝒢_z(k),`$
$`Z_\pm (k)`$ $`=`$ $`6h𝒢_\pm (k).`$ (5.10)
In Table VIII we present the collective form factor $`(k)`$ (for details of the derivation we refer to ). The form factors $`𝒢_m(k)`$ are given by
$`𝒢_z(k)`$ $`=`$ $`\delta _{M,0}m_3k_0a{\displaystyle \frac{\text{d}(k)}{\text{d}ka}},`$
$`𝒢_\pm (k)`$ $`=`$ $`\delta _{M,\pm 1}m_3k_0a\sqrt{L(L+1)}{\displaystyle \frac{(k)}{ka}}.`$ (5.11)
For any other model of baryons with the same spin-flavor structure, the corresponding results can be obtained by replacing Table VIII with the appropriate table (for example, by using harmonic oscillator wave functions as discussed in ).
The coefficients $`\zeta _m`$ are the spin-flavor matrix elements of $`X_3^Ms_{3,m}`$ which can either be evaluated for each channel separately, or more conveniently, by using the Wigner-Eckart theorem and isoscalar factors of $`SU_\mathrm{f}(3)`$ . The flavor wave functions are labeled by the quantum numbers $`(p,q),I,Y`$ corresponding to the reduction $`SU_\mathrm{f}(3)SU_\mathrm{I}(2)U_\mathrm{Y}(1)`$. In this notation $`(p,q)=(g_1g_2,g_2)`$, and hence we have $`(p,q)=(1,1)`$, $`(3,0)`$ or $`(0,0)`$ for the baryon flavor octet, decuplet and singlet, respectively, and $`(p,q)=(1,1)`$ or $`(0,0)`$ for the meson flavor octet and singlet, respectively. The spin-flavor matrix elements $`\zeta _m`$ for a given isospin channel can be expressed as
$`\zeta _m`$ $`=`$ $`{\displaystyle \underset{\gamma }{}}\begin{array}{cc}(p_f,q_f)& (p,q)\\ I_f,Y_f& I,Y\end{array}|\begin{array}{c}(p_i,q_i)_\gamma \\ I_i,Y_i\end{array}\alpha _{m,\gamma }.`$ (5.16)
The sum over $`\gamma `$ is over different multiplicities. The $`SU(3)`$ isoscalar factor which appears in Eq. (5.16) depends on the flavor multiplets $`(p,q)`$, the isospin $`I`$ and the hypercharge $`Y`$. A compilation of the $`SU(3)`$ isoscalar factors relevant for strong decays of baryons can be found in . Results for a specific charge channel can be obtained by multiplying $`\zeta _m`$ with the appropriate isospin Clebsch-Gordan coefficient $`I_f,M_{I_f},I,M_I|I_i,M_{I_i}`$.
In Tables IXXII we present the coefficients $`\alpha _{m,\gamma }`$ for strong decays into octet or decuplet baryons emitting a pseudoscalar meson (either octet or singlet). For strong decays of nonstrange baryons into the $`N\pi `$, $`N\eta `$, $`\mathrm{\Delta }\pi `$ and $`\mathrm{\Delta }\eta `$ channels the coefficients $`\zeta _m`$ are given explicitly in Tables III and IV of . Inspection of Tables IXXII and the isoscalar factors on page 184 of yields some interesting selection rules: (i) the $`B_{10}B_{10}+M_8`$ decay
$`\mathrm{\Sigma }^{}`$ $``$ $`\mathrm{\Sigma }^{}+\eta _8,`$ (5.17)
is forbidden since the $`SU(3)`$ isoscalar factor vanishes, and (ii) there is a spin-flavor selection rule for the $`{}_{}{}^{4}8[70,L^P]^28[56,0^+]+M_8`$ decays
$`N`$ $``$ $`\mathrm{\Lambda }+K,`$
$`\mathrm{\Lambda }`$ $``$ $`N+\overline{K},`$
$`\mathrm{\Xi }`$ $``$ $`\mathrm{\Xi }+\eta _8,`$ (5.18)
which is similar to the Moorhouse selection rule in electromagnetic couplings . However, the octet $`\eta _8`$ and singlet $`\eta _1`$ may mix because of $`SU_\mathrm{f}(3)`$ flavor symmetry breaking. The physical mesons $`\eta `$ and $`\eta ^{}`$ are then given in terms of a mixing angle
$`\eta `$ $`=`$ $`\eta _8\mathrm{cos}\theta _P\eta _1\mathrm{sin}\theta _P,`$
$`\eta ^{}`$ $`=`$ $`\eta _8\mathrm{sin}\theta _P+\eta _1\mathrm{cos}\theta _P.`$ (5.19)
The above mentioned forbidden two-body decays into a baryon and the octet meson $`\eta _8`$, are allowed for the physical mesons $`\eta `$ and $`\eta ^{}`$ via the octet-singlet mixing.
## 6 Comparison with experimental strong decays
With the definition of the transition operator in Eq. (5.2) and the helicity amplitudes, the decay widths for a specific isospin channel are given by
$`\mathrm{\Gamma }(BB^{}+M)`$ $`=`$ $`2\pi \rho {\displaystyle \frac{2}{2J+1}}{\displaystyle \underset{\nu >0}{}}|A_\nu (k)|^2.`$ (6.1)
Here we adopt the procedure of , in which the decay widths are calculated in the rest frame of the decaying resonance, and in which the relativistic expression for the phase space factor $`\rho `$ as well as for the momentum $`k`$ of the emitted meson are retained. The expressions for $`k`$ and $`\rho `$ are
$`k^2`$ $`=`$ $`m_M^2+{\displaystyle \frac{(m_B^2m_B^{}^2+m_M^2)^2}{4m_B^2}},`$
$`\rho `$ $`=`$ $`4\pi {\displaystyle \frac{E_B^{}E_Mk}{m_B}}`$ (6.2)
with $`E_B^{}=\sqrt{m_B^{}^2+k^2}`$ and $`E_M=\sqrt{m_M^2+k^2}`$. We consider here the strong decays of baryons in which a pseudoscalar meson (either octet or singlet) is emitted. The present calculation is an extension of in which we only discussed nonstrange decays of nonstrange baryons.
The calculated widths depend on the two parameters $`g`$ and $`h`$ in the transition operator of Eq. (5.3), and on the scale parameter $`a`$ of Eq. (5.6). In accordance with we keep $`g`$, $`h`$ and $`a`$ fixed for all resonances and all decay channels with the values $`g=1.164`$ fm, $`h=0.094`$ fm and $`a=0.232`$ fm (we note here that in the values of $`g`$ and $`h`$ were given in GeV<sup>-1</sup> instead of in fm). The decay widths of resonances that have been interpreted as a vibrational excitation of the string configuration of Fig. 1 depend on the coefficient $`\chi _1`$ for $`(v_1,v_2)=(1,0)`$ (the $`N(1440)`$, $`\mathrm{\Delta }(1600)`$, $`\mathrm{\Sigma }(1660)`$ and $`\mathrm{\Lambda }(1600)`$ resonances), or $`\chi _2`$ for $`(v_1,v_2)=(0,1)`$ (the $`N(1710)`$, $`\mathrm{\Sigma }(1940)`$ and $`\mathrm{\Lambda }(1810)`$ resonances) (see Table VIII). These coefficients are proportional to the intrinsic matrix element for each type of vibration. Here they are taken as constants with the values $`\chi _1=1.0`$ and $`\chi _2=0.7`$. For the pseudoscalar $`\eta `$ mesons we introduce a mixing angle $`\theta _P=23^{}`$ between the octet and singlet mesons . This value is consistent with that determined in a study of meson spectroscopy .
In comparison with other studies, we note that in the calculation in the nonrelativistic quark model of an elementary emission model is used, just as in the present calculation, but with the difference that the decay widths are parametrized by four reduced partial wave amplitudes instead of the two elementary amplitudes $`g`$ and $`h`$. Furthermore, the momentum dependence of these reduced amplitudes is represented by a constant. The calculations in the relativized quark model are done in a pair-creation model for the decay and involve a different assumption on the phase space factor . Both the nonrelativistic and relativized quark model calculations include the effects of mixing induced by the hyperfine interaction, which in the present calculations are not taken into account. It is important to note that we present a comparison of decay widths, rather than of decay amplitudes as was done in and for the nonrelativistic and relativized quark models.
In Tables XIIIXVII we compare the experimental strong decay widths of three and four star baryon resonances from the most recent compilation by the Particle Data Group with the results of our calculation for the nucleon, $`\mathrm{\Delta }`$, $`\mathrm{\Sigma }`$, $`\mathrm{\Lambda }`$ and $`\mathrm{\Xi }`$ families. We have used the experimental value of the mass of the decaying baryon. Strong couplings of missing resonances belonging to the flavor octet, decuplet and singlet are presented in Tables XVIIIXXI, XXIIXXV and XXVI, respectively.
The calculated decay widths are to a large extent a consequence of spin-flavor symmetry and phase space. The use of the collective form factors of Table VIII introduces a power-law dependence on the meson momentum $`k`$, compared to, for example, an exponential dependence for harmonic oscillator form factors. Our results for the strong decay widths are in fair overall agreement with the available data, and show that the combination of a collective string-like $`qqq`$ model of baryons and a simple elementary emission model for the decays can account for the main features of the data. There are a few exceptions which could indicate evidence for the importance of degrees of freedom which are outside the present $`qqq`$ model of baryons.
### 6.1 Nucleon resonances
The $`\pi `$ and $`\eta `$ decays of nucleon resonances have already been discussed in . Whereas the $`\pi `$ decays are in fair agreement with the data, the $`\eta `$ decays of octet baryons show an unusual pattern: the $`S`$-wave states $`N(1535)`$, $`\mathrm{\Sigma }(1750)`$ and $`\mathrm{\Lambda }(1670)`$ all are found experimentally to have a large branching ratio to the $`\eta `$ channel, whereas the corresponding phase space factor is very small . The small calculated $`\eta `$ widths ($`<0.5`$ MeV) for these resonances are due to a combination of spin-flavor symmetry and the size of the phase space factor. The results of our analysis suggest that the observed $`\eta `$ widths are not due to a conventional $`qqq`$ state, but may rather indicate evidence for the presence of a state in the same mass region of a more exotic nature, such as a pentaquark configuration $`qqqq\overline{q}`$ or a quasi-molecular $`S`$-wave resonance $`qqq`$-$`q\overline{q}`$ just below or above threshold, bound by Van der Waals type forces (for example $`N\eta `$, $`\mathrm{\Sigma }K`$ or $`\mathrm{\Lambda }K`$ ). In order to answer this question one has to carry out an analysis, similar to the present one, of the other configurations.
The $`K`$ decays are suppressed with respect to the $`\pi `$ decays because of phase space. In addition, the decay of the $`N(1650)`$, $`N(1675)`$ and $`N(1700)`$ resonances into $`\mathrm{\Lambda }K`$ is forbidden by the spin-flavor selection rule for the decay of $`{}_{}{}^{4}8[70,L^P]`$ nucleon states into this channel (see Sect. 5). For $`N(1675)`$ and $`N(1700)`$ only an upper limit is known, whereas the $`N(1650)`$ resonance has an observed width of $`12\pm 7`$ MeV. However, this resonance is just above the $`\mathrm{\Lambda }K`$ threshold which may lead to a coupling to a quasi-bound meson-baryon $`S`$ wave resonance. A study of the effect of spin-flavor symmetry breaking on these decays is in progress.
### 6.2 Delta resonances
The strong decay widths of the $`\mathrm{\Delta }`$ resonances are in very good agreement with the available experimental data. The same holds for the other resonances that have been assigned as decuplet baryons: $`\mathrm{\Sigma }^{}(1385)`$, $`\mathrm{\Sigma }^{}(2030)`$ and $`\mathrm{\Xi }^{}(1530)`$. For the decuplet baryons there is no $`S`$ state around the threshold of the various decay channels, so therefore there cannot be any coupling to quasi-molecular configurations. Just as for the nucleon resonances, the $`\eta `$ and $`K`$ decays are suppressed by phase space factors.
### 6.3 Sigma resonances
Strange resonances decay predominantly into the $`\pi `$ and $`\overline{K}`$ channel. Phase space factors suppress the $`\eta `$ and $`K`$ decays. The main discrepancy is found for $`\mathrm{\Sigma }(1750)`$. In the discussion of nucleon resonances it was suggested that the $`S`$ wave state $`\mathrm{\Sigma }(1750)`$ is the octet partner of $`N(1535)`$. It has a large observed $`\eta `$ width despite the fact that there is hardly any phase space available for this decay. This may indicate that it has a large quasi-molecular component.
The assignment of $`\mathrm{\Sigma }(1940)D_{13}`$ as a member of the $`(v_1,v_2)=(0,1)`$ vibrational band with $`{}_{}{}^{2}8_{3/2}^{}[56,1^{}]`$ and mass 1974 MeV is based on both its mass and its strong decay properties. If calculated with the observed mass, the other possible states, $`{}_{}{}^{2}10_{3/2}^{}[70,1^{}]`$ and $`{}_{}{}^{4}8_{3/2}^{}[70,1^{}]`$, both have very large widths ($``$ 100 MeV) in the $`\mathrm{\Delta }\overline{K}`$ and $`\mathrm{\Sigma }^{}\pi `$ channels, which is not supported by the data.
### 6.4 Lambda resonances
Also for $`\mathrm{\Lambda }`$ resonances the $`\eta `$ and $`K`$ decays are suppressed with respect to the $`\pi `$ and $`\overline{K}`$ channels because of phase space factors. Table XVI shows that the strong decays of $`\mathrm{\Lambda }`$ resonances show more discrepancies with the data than the other families of resonances.
We have assigned $`\mathrm{\Lambda }(1670)`$ as the octet partner of $`N(1535)`$ and $`\mathrm{\Sigma }(1750)`$. Its decay properties into the $`\eta `$ channel have already been discussed in Sect. 6.1 on the nucleon resonances.
The spin-flavor selection rule that was discussed in Sect. 5 forbids the decay of $`{}_{}{}^{4}8[70,L^P]`$ $`\mathrm{\Lambda }`$ states into the $`N\overline{K}`$ channel. Therefore, the calculated $`N\overline{K}`$ widths of $`\mathrm{\Lambda }(1800)`$, $`\mathrm{\Lambda }(1830)`$ and $`\mathrm{\Lambda }(2110)`$ vanish, whereas all of them have been observed experimentally . The $`\mathrm{\Lambda }(1800)S_{01}`$ state has large decay width into $`N\overline{K}^{}(892)`$ . Since the mass of the resonance is just around the threshold of this channel, this could indicate a coupling with a quasi-molecular $`S`$ wave. The $`N\overline{K}`$ width of $`\mathrm{\Lambda }(1830)`$ is relatively small ($`6\pm 3`$ MeV), and hence in qualitative agreement with the selection rule. The situation for the the $`\mathrm{\Lambda }(2110)`$ resonance is unclear.
The $`\mathrm{\Lambda }^{}(1405)`$ resonance has a anomalously large decay width ($`50\pm 2`$ MeV) into $`\mathrm{\Sigma }\pi `$. This feature emphasizes the quasi-molecular nature of $`\mathrm{\Lambda }^{}(1405)S_{01}`$ due to the proximity of the $`N\overline{K}`$ threshold. It has been shown that the inclusion of the coupling to the $`N\overline{K}`$ and $`\mathrm{\Sigma }\pi `$ decay channels produces a downward shift of the $`qqq`$ state toward or even below the $`N\overline{K}`$ threshold. In a chiral meson-baryon Lagrangian approach with an effective coupled-channel potential the $`\mathrm{\Lambda }^{}(1405)`$ resonance emerges as a quasi-bound state of $`N\overline{K}`$ .
### 6.5 Xi resonances
There is little experimental information available for the strong decays of baryons with strangeness $`2`$ ($`\mathrm{\Xi }`$) and $`3`$ ($`\mathrm{\Omega }`$). We find good agreement with the observed decay widths of the $`\mathrm{\Xi }(1820)`$ (octet) and the $`\mathrm{\Xi }(1530)`$ (decuplet) resonances (see Table XVII).
### 6.6 Missing resonances
For possible use in the analysis of new experimental data and in the search for missing resonances, we present in Tables XVIIIXXVI the strong decay widths of the missing resonances of Tables VVII. The states with $`L^P=1^+`$, $`L^P=2^{}`$ and $`[20,L^P]`$ are decoupled because of spin-flavor symmetry. Most of the strong couplings of the low-lying missing resonances are small, which to a large extent explains their status . Generally speaking, the orbital configurations that have the smallest strong couplings both for the octet, decuplet and singlet resonances are $`[56,2^+]`$ $`(v_1,v_2)=(0,0)`$, $`[70,1^{}](1,0)`$ and $`[70,1^{}](0,1)`$. It is interesting to note that in all cases the resonances associated with the configuration $`[56,1^{}](0,1)`$ exhibit large decay widths. The only resonance that we have assigned as one of these is $`\mathrm{\Sigma }(1940)`$. The majority of the missing resonances with sizeable decay widths belong either to the configuration $`[56,1^{}](0,1)`$ or to $`[70,2^+](0,0)`$.
Tables XVIII and XXII show that the missing nucleon and $`\mathrm{\Delta }`$ resonances that are associated with the configurations $`[70,2^+](0,0)`$, $`[70,0^+](0,1)`$ and $`[56,1^{}](0,1)`$ are predicted to have a large decay width into the $`\pi `$ channel. The $`\eta `$ and $`K`$ decays are suppressed with respect to the $`\pi `$ decays because of phase space. Strange baryons decay predominantly into the $`\pi `$ and $`\overline{K}`$ channels. The $`\eta `$ and $`K`$ widths are small in comparison, due to the available phase space. Inspection of Tables XIXXXI shows that the dominant decay channels of the missing octet baryons are $`N\overline{K}`$, $`\mathrm{\Sigma }\pi `$, $`\mathrm{\Delta }\overline{K}`$ for $`\mathrm{\Sigma }`$ resonances, $`N\overline{K}`$, $`\mathrm{\Sigma }^{}\pi `$ for $`\mathrm{\Lambda }`$ resonances, and $`\mathrm{\Sigma }\overline{K}`$, $`\mathrm{\Xi }\pi `$ for $`\mathrm{\Xi }`$ resonances. Similarly, we see from Tables XXIIIXXV that the missing decuplet baryons are most likely to couple to $`\mathrm{\Lambda }\pi `$, $`\mathrm{\Sigma }^{}\pi `$ for $`\mathrm{\Sigma }^{}`$ resonances, to $`\mathrm{\Sigma }\overline{K}`$, $`\mathrm{\Lambda }\overline{K}`$, $`\mathrm{\Xi }\pi `$, $`\mathrm{\Sigma }^{}\overline{K}`$ for $`\mathrm{\Xi }^{}`$ resonances, and to $`\mathrm{\Xi }\overline{K}`$ for $`\mathrm{\Omega }`$ resonances. Finally, Table XXVI shows that missing singlet baryons are most likely to decay into $`\mathrm{\Lambda }^{}N\overline{K}`$, $`\mathrm{\Sigma }\pi `$.
## 7 Electromagnetic couplings
In constituent models, electromagnetic couplings arise from the coupling of the (point-like) constituent parts to the electromagnetic field . We discuss here the case of the emission of a lefthanded photon
$`BB^{}+\gamma ,`$ (7.1)
for which the nonrelativistic part of the transverse electromagnetic coupling is given by
$`_{em}`$ $`=`$ $`2\sqrt{{\displaystyle \frac{\pi }{k_0}}}{\displaystyle \underset{j=1}{\overset{3}{}}}\mu _je_j\left[ks_{j,}\text{e}^{i\stackrel{}{k}\stackrel{}{r}_j}+{\displaystyle \frac{1}{2g_j}}(p_{j,}\text{e}^{i\stackrel{}{k}\stackrel{}{r}_j}+\text{e}^{i\stackrel{}{k}\stackrel{}{r}_j}p_{j,})\right],`$ (7.2)
where $`\stackrel{}{r}_j`$, $`\stackrel{}{p}_j`$ and $`\stackrel{}{s}_j`$ are the coordinate, momentum and spin of the $`j`$-th constituent, respectively; $`k_0`$ is the photon energy, and $`\stackrel{}{k}=k\widehat{z}`$ denotes the momentum carried by the outgoing photon. The photon is emitted by the $`j`$-th constituent: $`q_jq_j^{}+\gamma `$ (see Figure 8). The transition operator can be simplified by using the symmetry of the baryon wave functions, transforming to Jacobi coordinates and integrating over the baryon center-of-mass coordinate, to obtain
$`_{em}`$ $`=`$ $`6\sqrt{{\displaystyle \frac{\pi }{k_0}}}\mu _3e_3\left[ks_{3,}\widehat{U}{\displaystyle \frac{1}{g_3}}\widehat{T}_{}\right].`$ (7.3)
The operators $`\widehat{U}`$ and $`\widehat{T}_{}`$ are given in Eq. (5.5).
The transverse helicity amplitudes between the final ground state baryon belonging either to the $`J^P=1/2^+`$ octet with $`{}_{}{}^{2}8_{1/2}^{}[56,0^+]_{(0,0);0}`$ or to the $`J^P=3/2^+`$ decuplet with $`{}_{}{}^{4}10_{3/2}^{}[56,0^+]_{(0,0);0}`$, and the initial (excited) state of a baryon resonance are expressed as
$`A_\nu (k)`$ $`=`$ $`{\displaystyle }d\beta g(\beta )\mathrm{\Psi }^{}(0,S^{},S^{},\nu 1))|_{em}|\mathrm{\Psi }(L,S,J,\nu ),`$ (7.4)
$`=`$ $`6\sqrt{{\displaystyle \frac{\pi }{k_0}}}\left[kL,0;S,\nu |J,\nu _\nu L,1;S,\nu 1|J,\nu 𝒜_\nu \right],`$
where $`\nu =1/2`$, $`3/2`$ indicates the helicity. The orbit- and spin-flip amplitudes ($`𝒜_\nu `$ and $`_\nu `$, respectively) are given by
$`_\nu `$ $`=`$ $`{\displaystyle \text{d}\beta g(\beta )\mathrm{\Psi }^{}(0,0;S^{},\nu 1)|\mu _3e_3s_{3,}\widehat{U}|\mathrm{\Psi }(L,0;S,\nu )},`$
$`𝒜_\nu `$ $`=`$ $`{\displaystyle \text{d}\beta g(\beta )\mathrm{\Psi }^{}(0,0;S^{},\nu 1)|\mu _3e_3\widehat{T}_{}/g_3|\mathrm{\Psi }(L,1;S,\nu 1)}.`$ (7.5)
Here $`|\mathrm{\Psi }(L,M_L;S,M_S)`$ denote the (space-spin-flavor) angular momentum uncoupled wave functions of the initial and final baryons (see Appendix B).
In Tables XXVII and XXVIII we show the orbit- and spin-flip amplitudes for some radiative hyperon decays. These results were obtained under the assumption of $`SU_\mathrm{f}(3)`$ flavor symmetry, i.e. $`\mu _3=\mu _p`$ and $`g_3=g`$. The following selection rules apply: (i) the $`{}_{}{}^{4}10[56]^28[56]+\gamma `$ decays
$`\mathrm{\Sigma }^,`$ $``$ $`\mathrm{\Sigma }^{}+\gamma ,`$
$`\mathrm{\Xi }^,`$ $``$ $`\mathrm{\Xi }^{}+\gamma ,`$ (7.6)
are forbidden by $`U`$-spin conservation . These decays can only occur if flavor symmetry is broken ($`m_dm_s`$). Also the $`{}_{}{}^{2}1[70]^410[56]+\gamma `$ decay
$`\mathrm{\Lambda }^{}`$ $``$ $`\mathrm{\Sigma }^{,0}+\gamma ,`$ (7.7)
is forbidden by $`U`$-spin selection rules but, contrary to the previous cases, remains forbidden in the case of flavor symmetry breaking.
## 8 Comparison with experimental electromagnetic couplings
Radiative hyperon decay widths can be calculated from the helicity amplitudes as
$`\mathrm{\Gamma }(BB^{}+\gamma )`$ $`=`$ $`2\pi \rho {\displaystyle \frac{1}{(2\pi )^3}}{\displaystyle \frac{2}{2J+1}}{\displaystyle \underset{\nu >0}{}}|A_\nu (k)|^2.`$ (8.1)
Just as for the strong couplings, the electromagnetic decay widths are calculated assuming $`SU_{\mathrm{sf}}(6)`$ spin-flavor symmetry and using the rest frame of the decaying resonance
$`k`$ $`=`$ $`{\displaystyle \frac{m_B^2m_B^{}^2}{2m_B}},`$
$`\rho `$ $`=`$ $`4\pi {\displaystyle \frac{E_B^{}k^2}{m_B}}`$ (8.2)
with $`E_B^{}=\sqrt{m_B^{}^2+k^2}`$. The scale parameter $`a`$ was determined in a simultaneous fit to the proton charge radius, the proton electric and magnetic form factors and the neutron magnetic form factor to be $`a=0.232`$ fm . This is the same value as has been determined independently in a fit of the $`N\pi `$ decay widths of nucleon and delta resonances . For all cases we take the quark $`g`$ factors $`g=1`$. The quark scale magnetic moment is equal to the proton magnetic moment $`\mu _p`$, which corresponds to a constituent quark mass $`m=0.336`$ GeV.
Recently, the SELEX collaboration has measured the charge radius of the $`\mathrm{\Sigma }^{}`$ hyperon. The preliminary value is $`r^2_\mathrm{\Sigma }^{}=0.60\pm 0.08\pm 0.08`$ fm<sup>2</sup> . This value is in good agreement with our predicted value of $`r^2_\mathrm{\Sigma }^{}=r^2_p=12a^2=0.65`$ fm<sup>2</sup>.
The radiative decays between baryons with $`L=0`$ and $`L^{}=0`$ only involve the magnetic transitions. The corresponding widths can be expressed in terms of the transition magnetic moments $`\mu _{BB^{}}(k)`$ via
$`{\displaystyle \underset{\nu >0}{}}|A_\nu (k)|^2`$ $`=`$ $`4\pi k\mu _{BB^{}}^2(k).`$ (8.3)
In Table XXIX we show the transition moments for the $`{}_{}{}^{2}8[56]^28[56]+\gamma `$ and $`{}_{}{}^{4}10[56]^28[56]+\gamma `$ transitions. In the absence of the form factor (i.e. $`(k)=1`$), we recover the symmetry relations between the decuplet to octet transitions . For the conventions used in Appendices A and B we obtain the relations
$`\mu _{\mathrm{\Sigma }^0\mathrm{\Lambda }}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{3}}}\mu _p,`$
$`\mu _{\mathrm{\Delta }^+p}`$ $`=`$ $`\mu _{\mathrm{\Delta }^0n}=\mu _{\mathrm{\Sigma }^{,+}\mathrm{\Sigma }^+}=2\mu _{\mathrm{\Sigma }^{,0}\mathrm{\Sigma }^0}`$
$`=`$ $`{\displaystyle \frac{2}{\sqrt{3}}}\mu _{\mathrm{\Sigma }^{,0}\mathrm{\Lambda }}=\mu _{\mathrm{\Xi }^{,0}\mathrm{\Xi }^0}={\displaystyle \frac{2\sqrt{2}}{3}}\mu _p,`$
$`\mu _{\mathrm{\Sigma }^,\mathrm{\Sigma }^{}}`$ $`=`$ $`\mu _{\mathrm{\Xi }^,\mathrm{\Xi }^{}}=\mathrm{\hspace{0.33em}0}.`$ (8.4)
The numerical values are given in the third column of Table XXIX. A comparison with the last column shows the reduction of the transition magnetic moments due to the form factor $`(k)=1/(1+k^2a^2)^2`$.
The experimental information on radiative decays of hyperons is very limited. In Table XXX we present the radiative decay widths of low-lying hyperon resonances, and compare wherever possible with the data. The $`\mathrm{\Delta }^+p+\gamma `$ decay width is underpredicted by 35 $`\%`$, a common feature of all $`qqq`$ constituent quark models. This discrepancy has been shown to be due to nonresonant meson-exchange mechanisms . Just as for the energies and the strong decays, the $`\mathrm{\Lambda }^{}(1405)`$ resonance shows large deviations for the radiative decay widths, which once agains confirms its uncertain nature as a $`qqq`$ state. The forbidden decays $`\mathrm{\Lambda }^{}(1405)\mathrm{\Sigma }^{,0}+\gamma `$ and $`\mathrm{\Lambda }^{}(1520)\mathrm{\Sigma }^{,0}+\gamma `$ have not been observed. For comparison we also present the radiative decay widths of decuplet hyperons as obtained from lattice calculations and from a chiral constituent quark model with electromagnetic exchange currents between quarks . The negative parity hyperon decay widths for $`\mathrm{\Sigma }^,\mathrm{\Sigma }^{}+\gamma `$ and $`\mathrm{\Xi }^,\mathrm{\Xi }^{}+\gamma `$ have a small nonvanishing value in , whereas in the present calculation they are forbidden by flavor symmetry selection rules. In a subsequent publication we plan to investigate the effects of $`SU_\mathrm{f}(3)`$ flavor symmetry breaking on the electromagnetic couplings.
## 9 Summary and conclusions
We have presented in this article a systematic analysis of spectra and transition rates of strange baryons in the framework of a collective string-like $`qqq`$ model, in which the orbitally excited baryons are interpreted as collective rotations and vibrations of the strings. The algebraic structure of the model, both for the internal degrees of freedom of spin-flavor-color and for the spatial degrees of freedom, has been used to derive transparent results, such as a mass formula, selection rules and closed expressions for strong and electromagnetic couplings.
The situation is similar to that encountered for nonstrange baryons. While spectra are reasonably well described, transition rates, especially strong decay widths are only qualitatively described. The combination of a collective string-like $`qqq`$ model of baryons and a simple elementary emission model for the decays can account for the main features of the data. The main discrepancies are found for the low-lying $`S`$-wave states, specifically $`N(1535)`$, $`\mathrm{\Sigma }(1750)`$, $`\mathrm{\Lambda }^{}(1405)`$, $`\mathrm{\Lambda }(1670)`$ and $`\mathrm{\Lambda }(1800)`$. All of these resonances have masses which are close to the threshold of a meson-baryon decay channel, and hence they could mix with a quasi-molecular $`S`$ wave resonance of the form $`qqqq\overline{q}`$. In contrary, decuplet baryons have no low-lying $`S`$ states with masses close to the threshold of a particular decay channel, and their spectroscopy is described very well. The results of our analysis suggest that in future experiments particular attention be paid to the resonances mentioned above in order to elucidate their structure, and to look for evidence of the existence of exotic (non $`qqq`$) configurations of quarks and gluons.
In our calculations we have included only a diagonal breaking of the spin-flavor symmetry. This seems to be a good approximation to the actual situation and no major discrepancy appear to be related to non-diagonal breakings. A study of the effects of $`SU_\mathrm{f}(3)`$ flavor symmetry breaking due to different quark masses on the radiative decays and strong couplings is in progress, and will be published separately.
This paper concludes our analysis of $`q^3`$ configurations in baryons. The next step is the study of more complex configurations of quarks and gluons, such as hybrid quark-gluon states $`qqqg`$, pentaquark states $`q^4\overline{q}`$ and multiquark meson-baryon bound states $`qqqq\overline{q}`$.
## Acknowledgements
This work is supported in part by DGAPA-UNAM under project IN101997, by CONACyT under project 32416-E and by D.O.E. Grant DE-FG02-91ER40608.
## Appendix A Spin-flavor wave functions
Here we list the conventions used for the spin and flavor wave functions which are consistent with the choice of Jacobi coordinates of Eq. (2.3). They coincide with the conventions of .
### A.1 Spin wave functions
The spin wave functions $`|S,M_S`$ are given by :
$`|1/2,1/2`$ $`:`$ $`\chi _\rho =[||]/\sqrt{2},`$
$`:`$ $`\chi _\lambda =[2|||]/\sqrt{6},`$
$`|3/2,3/2`$ $`:`$ $`\chi _S=|.`$ (A.1)
We only show the state with the largest component of the projection $`M_S=S`$. The other states are obtained by applying the lowering operator in spin space.
### A.2 Flavor wave functions
For the flavor wave functions $`|(p,q),I,M_I,Y`$ we adopt the convention of with $`(p,q)=(g_1g_2,g_2)`$.
(i) The octet baryons $`(p,q)=(1,1)`$:
$`|(1,1),1/2,1/2,1`$ $`:`$ $`\varphi _\rho (p)=[|udu|duu]/\sqrt{2},`$
$`:`$ $`\varphi _\lambda (p)=[2|uud|udu|duu]/\sqrt{6},`$
$`|(1,1),1,1,0`$ $`:`$ $`\varphi _\rho (\mathrm{\Sigma }^+)=[|suu|usu]/\sqrt{2},`$
$`:`$ $`\varphi _\lambda (\mathrm{\Sigma }^+)=[|suu+|usu2|uus]/\sqrt{6},`$
$`|(1,1),0,0,0`$ $`:`$ $`\varphi _\rho (\mathrm{\Lambda })=[2|uds2|dus|dsu+|sdu|sud+|usd]/\sqrt{12},`$
$`:`$ $`\varphi _\lambda (\mathrm{\Lambda })=[|dsu|sdu+|sud+|usd]/2,`$
$`|(1,1),1/2,1/2,1`$ $`:`$ $`\varphi _\rho (\mathrm{\Xi }^0)=[|sus|uss]/\sqrt{2},`$ (A.2)
$`:`$ $`\varphi _\lambda (\mathrm{\Xi }^0)=[2|ssu|sus|uss]/\sqrt{6}.`$
(ii) The decuplet baryons $`(p,q)=(3,0)`$:
$`|(3,0),3/2,3/2,1`$ $`:`$ $`\varphi _S(\mathrm{\Delta }^{++})=|uuu,`$
$`|(3,0),1,1,0`$ $`:`$ $`\varphi _S(\mathrm{\Sigma }^+)=[|suu+|usu+|uus]/\sqrt{3},`$
$`|(3,0),1/2,1/2,1`$ $`:`$ $`\varphi _S(\mathrm{\Xi }^0)=[|ssu+|sus+|uss]/\sqrt{3},`$
$`|(3,0),0,0,2`$ $`:`$ $`\varphi _S(\mathrm{\Omega }^{})=|sss.`$ (A.3)
(iii) The singlet baryons $`(p,q)=(0,0)`$:
$`|(0,0),0,0,0`$ $`:`$ $`\varphi _A(\mathrm{\Lambda })=[|uds|dus+|dsu|sdu+|sud|usd]/\sqrt{6}.`$ (A.4)
We only show the highest charge state $`M_I=I`$ with $`Q=I+Y/2`$. The other charge states are obtained by applying the lowering operator in isospin space.
## Appendix B Baryon wave functions
The $`S_3`$ invariant space-spin-flavor ($`\mathrm{\Psi }=\psi \chi \varphi `$) baryon wave functions are given by
$`{}_{}{}^{2}8[56,L^P]`$ $`:`$ $`\psi _S(\chi _\rho \varphi _\rho +\chi _\lambda \varphi _\lambda )/\sqrt{2},`$
$`{}_{}{}^{2}8[70,L^P]`$ $`:`$ $`[\psi _\rho (\chi _\rho \varphi _\lambda +\chi _\lambda \varphi _\rho )+\psi _\lambda (\chi _\rho \varphi _\rho \chi _\lambda \varphi _\lambda )]/2,`$
$`{}_{}{}^{4}8[70,L^P]`$ $`:`$ $`(\psi _\rho \varphi _\rho +\psi _\lambda \varphi _\lambda )\chi _S/\sqrt{2},`$
$`{}_{}{}^{2}8[20,L^P]`$ $`:`$ $`\psi _A(\chi _\rho \varphi _\lambda \chi _\lambda \varphi _\rho )/\sqrt{2},`$
$`{}_{}{}^{4}10[56,L^P]`$ $`:`$ $`\psi _S\chi _S\varphi _S,`$
$`{}_{}{}^{2}10[70,L^P]`$ $`:`$ $`(\psi _\rho \chi _\rho +\psi _\lambda \chi _\lambda )\varphi _S/\sqrt{2},`$
$`{}_{}{}^{2}1[70,L^P]`$ $`:`$ $`(\psi _\rho \chi _\lambda \psi _\lambda \chi _\rho )\varphi _A/\sqrt{2},`$
$`{}_{}{}^{4}1[20,L^P]`$ $`:`$ $`\psi _A\chi _S\varphi _A.`$ (B.1)
The quark orbital angular momentum $`L`$ is coupled with the spin $`S`$ to the total angular momentum $`J`$ of the baryon.
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# Magnetic Foehn Effect in Nonadiabatic Transition
## Abstract
The magnetization curves as a response of sweeping magnetic field in the thermal environment are investigated using the quantum master equation. In a slow velocity region where the system almost behaves adiabatically, the magnetic plateau appears which has been observed in the recent experiment of V<sub>15</sub> \[ Phys. Rev. Lett. (2000)\]. We investigate its mechanism and propose that this phenomenon is quite universal in the quasi-adiabatic transition with small inflow of the heat, and we call it ’Magnetic Foehn Effect’. We observe the crossover between this mechanism and the Landau-Zener-Stückelberg mechanism changing the velocity. Some experiment is proposed to clarify the inherent mechanism of this effect.
PACS number: 75.40.Gb,76.20.+q
Properties of quantum dynamics have been studied extensively for nanoscale magnets and also in microscopic system with nanostructure . There the nonadiabatic transition plays important roles. Landau, Zener, and Stükelberg (LZS) derived the well-known nonadiabatic transition probability for a two-level system with a sweeping field. It depends on the sweeping velocity and energy gap. Although the LZS formula is derived in two-level systems, it can also describe nonadiabatic transitions at the avoided crossings in the uniaxial magnets . Therefore it can be widely applied to analyze phenomena related to the nonadiabatic transition in various materials . However real experiments are always done in a thermal environment. Hence the studies on the nonadiabatic transition with an environment are quite crucial for further understanding of time-dependent phenomena in magnetic systems . There universal aspects of thermal effect independent of details of reservoirs are very important. As an example of such universal aspects, we have studied the magnetization process in uniaxial molecule magnets such as Mn<sub>12</sub> and Fe<sub>8</sub> at very low temperatures in thermal environment . There the step-wise magnetization process is observed, which is due to the nonadiabatic transition and a fast damping process to the ground state. We call it apparent (or deceptive) nonadiabatic process. We also found that the transition probability of purely quantum mechanics is deducible from this process. This mechanism does not depend on the detailed structure of reservoir apid fluctuation of noise at the resonant field
In this Letter, we propose another universal qualitative aspect of the thermal effect on the quasi-adiabatic transition where the LZS probability is almost one. We understand that the effect is quite general in systems which behave almost adiabatically in the dissipative environment. The present study is directly related to the recent experiment for the molecule V<sub>15</sub> which is effectively regarded as two level system . This molecule has $`15`$ atoms of V but they divided into subclusters of $`6`$,$`3`$, and $`6`$ spins. The subcluster of $`6`$ spins forms a singlet state and contributes little to the magnetization and only the cluster of $`3`$ spins mainly contributes to the magnetization. Therefore the model can be regarded as a two level system. This molecule is a very simple system and we may expect to see the LZS process clearly . However it was observed that the scattered population (i.e., that at the excited level) decreases when the sweeping velocity becomes faster . This is opposite to what we expect in view point of the LZS mechanism where the probability of the adiabatic transition should decrease for faster change of the field and the population of the excited state should increase.
Chiorescu et al. explained this behavior in the view of the phonon bottleneck effect which means a lack of phonon number which contributes to the excitation from the ground state at near resonant magnetic fields. In the experiment, the heat reserver has a double-structure, i.e., the spin system is attached to the phonon system of the crystal, and the phonon system is attached to the external reservoir which is the liquid He. The contact between the phonon system and the external reservoir is so week that in short time scale, thermal effect in the spin system is caused by only phonons. Due to this effect, the population of the excited state does not increase enough, and saturates at some value, which causes a magnetic plateau.
In this Letter, we show that plateau in the magnetization curve is also observed in the case that the spin system is connected with a single heat reserver with slow relaxation rate, and propose that this qualitative property is universal with regardless to the detailed structure of environment. We investigate a magnetization process for a sweeping field by making use of the quantum master equation which we have used in our studies of quantum dynamics in dissipative environments . Thereby, we investigate the magnetization plateau for various sweeping velocities and temperatures.
The Hamiltonian we shall consider is given by
$$=H(t)S^z+\mathrm{\Gamma }S^x,$$
(1)
where $`H(t)`$ is the sweeping field, $`H(t)=vt`$ and $`\mathrm{\Gamma }`$ is the transverse field. Transverse field represents a term causing quantum fluctuation and does not commute with the magnetization $`S_z`$. This simple system is realized in many cases, e.g., the isotropic anti-ferromagnetic Heisenberg chain with odd number of spins has the doublet in the ground state. Actually V<sub>15</sub> is in this situation .
For this system (1), the LZS transition probability is given by,
$$P_{\mathrm{LZS}}(v)=1\mathrm{exp}\left[\frac{\pi \mathrm{\Gamma }^2}{v}\right],$$
(2)
in the case of $`S=1/2`$ spin system . Thus the normalized magnetization at $`t=\mathrm{}`$ is given by,
$`M_{\mathrm{out}}=12P_{\mathrm{LZS}}(v).`$ (3)
We should note that this expression of $`M_{\mathrm{out}}`$ is also exact for any values for $`S`$ although (2) is derived for $`S=\frac{1}{2}`$.
We introduce a thermal environment taking the phonon system as the bath, $`_\mathrm{B}=_\omega \omega b_\omega ^{}b_\omega `$, where $`b_\omega `$ and $`b_\omega ^{}`$ are the annihilation and creation boson operators of the frequency $`\omega `$. We adopt the spectral density of the boson bath $`I(\omega )`$ in the form $`I(\omega )=I_0\omega ^\alpha (\omega 0),\mathrm{\hspace{0.17em}\hspace{0.17em}0}(\omega <0)`$, with $`\alpha =2`$. In the experimental situation in the magnetic molecules such as $`S=10`$ in Mn<sub>12</sub> and Fe<sub>8</sub>, the hyperfine interaction and the dipole interaction are not negligible at very low temperatures . In the case of V<sub>15</sub>, the phonon gives a dominant contribution .
In the case of phonon bath, we can derive an equation of motion of the reduced density matrix $`\rho `$ tracing out the degree of freedom of the bath in the following form (the quantum master equation ):
$$\frac{\rho (t)}{t}=\frac{1}{i\mathrm{}}[,\rho (t)]\lambda \left([X,R\rho (t)]+[X,R\rho (t)]^{}\right),$$
(4)
where $`X`$ is a system operator through which the system and the bath couple with the constant $`\lambda `$. The first term of the right-hand side describes the pure quantum dynamics of the system while the second term represents effects of environments at a temperature $`T(=\beta ^1)`$. There $`R`$ is defined as follows:
$`k|R|m`$ $`=`$ $`\zeta (E_kE_m)n_\beta (E_kE_m)k|X|m,`$ (5)
$`\zeta (\omega )`$ $`=`$ $`I(\omega )I(\omega ),\mathrm{and}n_\beta (\omega )=(e^{\beta \omega }1)^1,`$ (6)
where $`|k`$ and $`|m`$ represent the eigenstates of $``$ with the eigenenergies $`E_k`$ and $`E_m`$, respectively.
We simulate the evolution given by Eq. (4) for various sweeping velocities in $`1/2`$ spin system of (1). Throughout this Letter, we set $`\mathrm{}`$ to be unity. From now on, we set parameters as $`\mathrm{\Gamma }=0.5,T=1.0`$, and $`\lambda =0.001`$. In Fig.1(a), we present the magnetization curves for fast sweeping rates, $`v=0.1,0.2`$, and $`0.4`$. Here we clearly find that the magnetic plateau decreases when $`v`$ increases, which is consistent with (3). On the other hand, in the case of slow sweeping rates we also find the magnetic plateau as shown in Fig.$`1`$(b) although in these sweeping rates the LZS transition probabilities (2) are almost one. Here we should note that the magnetic plateau increases when when $`v`$ increases, which is an opposite property to the fast sweeping case. This is the same phenomenon as the experiment .
In Fig.2 we show the nonmonotonic dependence of the plateau height $`M_{\mathrm{out}}`$ on velocities. Trivially when $`v`$ is very large, $`M_{\mathrm{out}}`$ shows linear dependence on $`1/v`$ from Eq. (3). In the slow sweeping rate region with $`P_{\mathrm{LZS}}1`$, $`M_{\mathrm{out}}`$ goes down.
Here we discuss the population in the excited level $`\rho _{22}(H)`$ for various sweeping velocities. Time dependences of $`\rho _{22}(H)`$ are shown in Fig.$`3`$, where the population of $`\rho _{22}(H)`$ changes as a function of the sweeping velocity. Because the matrix element of the operator $`R`$ is proportional to $`\lambda `$ and $`\mathrm{\Delta }^\alpha ,(\mathrm{\Delta }=\sqrt{H^2+\mathrm{\Gamma }^2})`$ , the thermalization rate $`\gamma (\mathrm{\Delta },\lambda )`$ is proportional to these values as $`\gamma (\mathrm{\Delta },\lambda )\lambda \mathrm{\Delta }^\alpha `$.
Let us investigate the relation between the thermalization rate $`\gamma (\mathrm{\Delta },\lambda )`$ and the sweeping velocity $`v`$. If $`v`$ is much larger, i.e., $`v\gamma (H,\lambda )`$, then no thermal relaxation occurs. In this case, abrupt change in the distribution $`\rho _{22}`$ due to the nonadiabatic transition takes place at $`H=0`$. This behavior is demonstrated in the case of $`v=2\times 10^1`$, and $`4\times 10^1`$ in Fig.$`3`$. This sudden change causes the magnetic jump as shown in Fig.1(a). Here it should be noted that we see the precession which means that the state is still highly coherent. In this region of the velocity, when the velocity becomes small, the plateau goes up because $`P_{\mathrm{LZS}}`$ monotonically increases until a simple adiabatic magnetization curve appears due to $`P_{\mathrm{LZS}}1`$.
In the further slow velocity, thermalization process begins to take place, that is, $`\rho _{22}(H)`$ tends to relax to its equilibrium value,
$`\rho _{22}^{\mathrm{eq}}(H)={\displaystyle \frac{e^{\beta \mathrm{\Delta }}}{1+e^{\beta \mathrm{\Delta }}}}.`$ (7)
This increase of $`\rho _{22}(H)`$ causes a magnetic plateau as shown for $`v=2\times 10^3`$, and $`8\times 10^3`$ in Fig.$`3`$. In order to realize a visible plateau, thermalization rate should be small comparing with $`v`$. Actually, because of $`\gamma (H,\lambda )\lambda \mathrm{\Delta }^2`$, thermalization is very slow at around the resonant point due to small $`\mathrm{\Delta }`$. We also studied the system with $`\alpha =0`$, where we found that the plateau sustain for large values of $`H(t)`$ and the shape of the magnetization process seems very different.
We associate the mechanism of magnetic plateau with the well-known Foehn phenomenon in the meteorology. The air with the vapor climbs up the mountain getting colder adiabatically and next the rain ensues. At this moment the vapor gives the heat to the air as the latent heat. Then, the air alone goes down the mountain and the temperature of the air increases higher than the original one due to the inflow of the latent heat at the mountain. In the present magnetic system, ’climbing up the mountain’ corresponds to ’sweeping of magnetic field to $`H=0`$’, ’the latent heat’ to ’the heat from the phonon’, and ’the increase of temperature of the air’ to ’the increase of temperature of the spin’. After the plateau the magnetization relaxes to the equilibrium value due to cooling by the thermal bath, which corresponds to the cooling of the air by the land after hot air reaches to the ground. Thus these similarities lead us to call this magnetic phenomenon ’Magnetic Foehn effect’. We confirmed that the Magnetic Foehn effect is also observed in the Hamiltonian (1) with larger $`S`$. The essential mechanism of this phenomenon is inflow of heat during adiabatic process. Because the mechanism is quite simple, we can say that the ’Magnetic Foehn effect’ is a universal phenomenon in the magnetic systems which behave almost adiabatically.
Because the LZS transition occurs only in the vicinity of $`H=0`$ and at other points only relaxations due to the dissipation term occur, the time evolution process may be divided into three regions:$`\rho (t)𝒰_\mathrm{d}(t:0)𝒰_{\mathrm{LZS}}𝒰_\mathrm{d}(0:\mathrm{})\rho (\mathrm{})`$, where $`𝒰_\mathrm{d}(t:s)`$ is the time evolution operator due to dissipation part from a time $`s`$ to $`t`$, and $`𝒰_{\mathrm{LZS}}`$ represents the LZS scattering matrix which corresponds to the pure quantum dynamics in the first term in (4) . For slow velocity case $`v1`$, we can consider that the system behaves adiabatically and thus we put the evolution $`𝒰_{\mathrm{LZS}}`$ to be unity. Thus the transfer between the levels occurs only by the dissipation $`𝒰_\mathrm{d}(t:s)`$. If we write down the second term of (4) explicitly, we have the following equations for the diagonal term
$`\dot{\rho }_{22}={\displaystyle \frac{1}{2}}X_{12}\lambda n_\beta \left(\mathrm{\Delta }\right)\zeta \left(\mathrm{\Delta }\right)\left[(e^{\beta \mathrm{\Delta }}+1)\rho _{22}1\right],`$ (8)
where $`X_{12}`$ is the matrix elements of the operator $`X`$. In our simulation, it is expressed as $`X_{12}=\frac{\left|H\mathrm{\Gamma }\right|}{\sqrt{H^2+\mathrm{\Gamma }^2}}`$ for $`X=S_x+S_z`$ . The term $`X_{12}`$ depends on the choice of $`X`$, and is not universal. This equation (8) leads the Magnetic Foehn effect.
At the end of this Letter, let us discuss some experimental situations. As for the observation of the crossover from the Magnetic Foehn region to that of LZS region is difficult in V<sub>15</sub>, because the crossover from the Magnetic Foehn region to the LZS region locates at a very fast sweeping rate. For example, if $`\mathrm{\Gamma }=0.1`$K, the crossover sweeping rate $`v`$ is of order $`10^7`$ H$`/`$sec. Such fast change of the field is not easy to realize. However when the gap become small, the field at the crossover sweeping rate $`v`$ becomes small. In Fig.$`4`$, we show the dependence of the gap on the size and anisotropy for the system:
$$=J\underset{i}{}\left(S_i^xS_{i+1}^x+S_i^yS_{i+1}^y+AS_i^zS_{i+1}^z\right),$$
(9)
where $`𝐒_i`$ is an $`S=1/2`$ spin and the open boundary condition is adopted. The gap becomes small with $`A`$, and decreases with $`N`$ exponentially. Thus we expect that the crossover would be observed in some system and there we can obtain many informations of effects due to combination of the quantum process and the thermal effect.
In V<sub>15</sub>, Chiorescu et al. attributes the Magnetic Foehn phenomenon to the lack of phonon as a mechanism of the insufficient supply of heat. In such a situation, we would propose an experiment to realize the adiabatic transition of an isolated spin system. That is, we first sweep the field from the large negative value to the field ($`H_\mathrm{p}`$) where the magnetic plateau $`(M_\mathrm{p})`$ is observed. At this point the number of phonon is very small and supply of phonons from out side is expected to be slow. In this circumstance, if the field is swept in the opposite direction from $`H_\mathrm{p}`$ to $`H_\mathrm{p}`$, the spins at the lower spin can not be excited by the phonon because no phonon is available, and behaves pure quantum mechanically. In the experiment the LZS probability is almost one. The spins at the higher level may emit the phonon and relax to the lower level. But the emitted phonon will be used to excite the spin again because the population of the upper level is much smaller than that of equilibrium. Thus we expect that the magnetization simply changes the sign when $`H_\mathrm{p}H_\mathrm{p}`$. In the iteration of this process $`H_\mathrm{p}H_\mathrm{p}H_\mathrm{p}H_\mathrm{p}\mathrm{}`$, the magnetization would maintain the same amplitude for a while, $`1M_\mathrm{p}M_\mathrm{p}M_\mathrm{p}M_\mathrm{p}\mathrm{}`$, before the heat flows in from the external environment and equilibrates the system. From this slow relaxation of magnetization we could know the relaxation rate between the phonon and the external bath. On the other hand, in the Magnetic Foehn phenomenon due to slow relaxation but not short of phonon number, $`M_\mathrm{p}`$ relaxes with the thermalization rate. From Eq.(8), we can derive the relation of $`\rho _{22}`$ for the iteration: $`\rho _{22}^{(n)}=p_1+p_2\rho _{22}^{(n1)}`$ with $`\rho _{22}^{(0)}=0`$ and $`p_1=(1M_\mathrm{p})/2`$. The sequence $`M_\mathrm{p}^{(n)}`$ is given by $`|M^{(n)}|=12\rho _{22}^{(n)}`$ where $`\rho _{22}^{(n)}=(1M_\mathrm{p})(1p_2^n)/(1p_2)`$. Here $`p_2`$ is given by
$`p_2=\mathrm{exp}\left(\lambda ^{}{\displaystyle _{H_\mathrm{p}/v}^{H_\mathrm{p}/v}}𝑑\tau \mathrm{\Delta }^\alpha (\tau )\mathrm{coth}(\beta \mathrm{\Delta }(\tau )/2)\right).`$ (10)
We would like to thank Professor B. Barbara and Dr. I. Chiorescu for valuable communications for their work of the reference. The present work is partially supported by Grant-in-Aid for Scientific Research from Ministry of Education, Science, Sports and Culture of Japan.
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# Random Words, Toeplitz Determinants and Integrable Systems. II.
## 1 Introduction
This paper, a continuation of , connects the analysis of the length of the longest weakly increasing subsequence of inhomogeneous random words to a Riemann-Hilbert problem and an associated system of integrable PDEs. That such a connection exists is not so surprising given the fundamental work of Baik, Deift and Johansson connecting the related problem involving random permutations to a Riemann-Hilbert problem. For the reader’s convenience we first summarize some of the results of before presenting our new results.
A word is a string of symbols, called letters, which belong to an ordered alphabet $`𝒜`$ of fixed size $`k`$. The set of all such words of length $`N`$, $`𝒲(𝒜,N)`$, forms the sample space in our statistical analysis. We equip the space $`𝒲(𝒜,N)`$ with a natural inhomogeneous measure by assigning to each letter $`i𝒜`$ a probability $`p_i`$ and defining the probability measure on words by the product measure. We also order the $`p_i`$ so that
$$p_1p_2\mathrm{}p_k$$
and decompose our alphabet $`𝒜`$ into subsets $`𝒜_1`$, $`𝒜_2`$, …, $`𝒜_M`$, $`Mk`$, such that $`p_i=p_j`$ if and only if $`i`$ and $`j`$ belong to the same $`𝒜_\alpha `$.
Let
$$w=\alpha _1\alpha _2\mathrm{}\alpha _N𝒲,\alpha _i𝒜,$$
be a word. A weakly increasing subsequence of the word $`w`$ is a subsequence $`\alpha _{i_1}\alpha _{i_2}\mathrm{}\alpha _{i_m}`$ such that $`i_1<i_2<\mathrm{}<i_m`$ and $`\alpha _{i_1}\alpha _{i_2}\mathrm{}\alpha _{i_m}`$. The positive integer $`m`$ is called the length of this weakly increasing subsequence. For each word $`w𝒲`$ we define $`\mathrm{}_N(w)`$ to equal the length of the longest weakly increasing subsequence in $`w`$.<sup>1</sup><sup>1</sup>1There may be many subsequences of $`w`$ that have the identical length $`\mathrm{}_N(w)`$. The function
$$\mathrm{}_N:𝒲(𝒜,N)\text{R}$$
is the principal random variable in our analysis, and the corresponding distribution function,
$$F_N(n):=\text{Prob}\left(\mathrm{}_N(w)n\right),$$
is our principal object. (Prob is the inhomogeneous measure on random words; it depends upon $`N`$ and the probabilities $`p_i`$.)
To formulate the basic result of , define
$$k_\alpha =|𝒜_\alpha |,$$
where
$$𝒜=\underset{\alpha =1}{\overset{M}{}}𝒜_\alpha $$
is the decomposition of the alphabet $`𝒜`$ introduced above, then
$`\underset{N\mathrm{}}{lim}\text{Prob}\left({\displaystyle \frac{\mathrm{}_NNp_1}{\sqrt{Np_1}}}s\right)`$ $`=`$ $`(2\pi )^{(k1)/2}{\displaystyle \underset{\alpha }{}}(1!\mathrm{\hspace{0.17em}2}!\mathrm{}(k_\alpha 1)!)^1\times `$
$`{\displaystyle \genfrac{}{}{0pt}{}{{\displaystyle \mathrm{}}}{\genfrac{}{}{0pt}{}{\xi _i\mathrm{\Xi }}{\xi _1s}}}{\displaystyle \underset{\alpha }{}}\mathrm{\Delta }_\alpha (\xi )^2e^{{\scriptscriptstyle \xi _i^2/2}}\delta ({\displaystyle \sqrt{p_i}\xi _i})d\xi _1\mathrm{}d\xi _k,`$
where $`\mathrm{\Delta }_\alpha (\xi )`$ is the Vandermonde determinant of those $`\xi _i`$ with $`i𝒜_\alpha `$, and $`\mathrm{\Xi }`$ denotes the set of those $`\xi _i`$ that $`\xi _{i+1}\xi _i`$ whenever $`i`$ and $`i+1`$ belong to the same $`𝒜_\alpha `$.
This result has the following random matrix interpretation. The limiting distribution function (as $`N\mathrm{}`$) for the appropriately centered and normalized random variable $`\mathrm{}_N`$ is related to the distribution function for the eigenvalues $`\xi _i`$ in the direct sum of mutually independent $`k_\alpha \times k_\alpha `$ Gaussian unitary ensembles,<sup>2</sup><sup>2</sup>2A basic reference for random matrices is Mehta’s book . conditional on the eigenvalues $`\xi _i`$ satisfying $`\sqrt{p_i}\xi _i=0`$. In the case when one letter occurs with greater probability than the others, this result implies that the limiting distribution of $`(\mathrm{}_NNp_1)/\sqrt{N}`$ is Gaussian with variance equal to $`p_1(1p_1)`$. In the case when all the probabilities are distinct, we proved the refined asymptotic result
$$\text{E}(\mathrm{}_N)=Np_1+\underset{j>1}{}\frac{p_j}{p_1p_j}+O(\frac{1}{\sqrt{N}}),N\mathrm{}.$$
The derivation of the above asymptotic formulae follows from a direct asymptotic analysis of the right hand side of the basic combinatorial equation,
$$\text{Prob}\left(\mathrm{}_N(w)n\right)=\underset{\genfrac{}{}{0pt}{}{\lambda N}{\lambda _1n}}{}s_\lambda (p)f^\lambda .$$
Here $`\lambda N`$ denotes a partition of $`N`$, $`s_\lambda (p)`$ is the Schur function of shape $`\lambda `$ evaluated at $`p:=(p_1,p_2,\mathrm{},p_k,0,0,\mathrm{})`$, and $`f^\lambda `$ equals the number of standard Young tableaux of shape $`\lambda `$, see, e.g. . After was written, Stanley showed that the measure $`\text{Prob}(\{\lambda \}):=s_\lambda (p)f^\lambda `$ also underlies the analysis of certain (generalized) riffle shuffles of Bayer and Diaconis . Stanley relates this measure to quasisymmetric functions and does not require that $`p`$ have finite support. (Many of our results generalize to the case when $`p`$ does not have finite support, but we do not consider this here.) The measure considered here and in is a specialization of the Schur measure $`\text{Prob}(\{\lambda \}):=s_\lambda (x)s_\lambda (y)`$ . For the Schur measure, Okounkov has shown that the associated correlation functions satisfy an infinite hierarchy of PDEs; namely, the Toda lattice hierarchy of Ueno and Takasaki . Similar results were also obtained by Adler and van Moerbeke .
Gessel’s theorem (see also ) implies that the (exponential) generating function of $`\text{Prob}(\mathrm{}_Nn)`$ is a Toeplitz determinant<sup>3</sup><sup>3</sup>3If $`\varphi `$ is a function on the unit circle with Fourier coefficients $`\varphi _j:=1/2\pi _\pi ^\pi e^{ij\theta }\varphi (e^{i\theta })𝑑\theta `$ then $`T_n(\varphi )`$ denotes the Toeplitz matrix $`\left(\varphi _{ij}\right)_{i,j=0,1,\mathrm{},n1}`$ and $`D_n(\varphi )`$ its determinant.
$$G_I(n;\{p_i\},t):=\underset{N=0}{\overset{\mathrm{}}{}}\text{Prob}\left(\mathrm{}_N(w)n\right)\frac{t^N}{N!}=D_n(f_I),$$
(1.1)
where
$$f_I(z)=e^{t/z}\underset{j=1}{\overset{k}{}}(1+p_jz).$$
Probabilistically, $`G_I(n;\{p_i\},t)`$ is the Poissonization of $`\mathrm{}_N`$. Similar Poissonizations have proved crucial in the analysis of the length of the longest increasing subsequences in random permutations (see also and references therein).
In the present paper we use (1.1) to express $`G_I(n;\{p_i\},t)`$ in terms of the solution of a certain integrable system of nonlinear PDEs. Indeed, we show that $`G_I(n;\{p_i\},t)`$ can be identified as the Jimbo-Miwa-Ueno $`\tau `$-function corresponding to the (generalized) Schlesinger isomonodromy deformation equations of the $`2\times 2`$ matrix linear ODE which has $`M+1`$ simple poles in the finite complex plane and one Poincaré index 1 irregular singular point at infinity. Recall that the number $`M`$ is the total number of the subsets $`𝒜_\alpha 𝒜`$. The poles are located at $`0`$ and $`p_{i_\alpha }`$ ($`i_\alpha `$ = max $`𝒜_\alpha `$). The integers $`k_\alpha `$ appear as the formal monodromy exponents at the respective points $`p_{i_\alpha }`$. We also evaluate the remaining monodromy data and formulate a $`2\times 2`$ matrix Riemann-Hilbert problem which provides yet another analytic representation for the function $`G_I(n;\{p_i\},t)`$. Similar to the problems considered in and , the Riemann-Hilbert representation of $`G_I(n;\{p_i\},t)`$ can be used for the further asymptotic analysis of the random variable $`\mathrm{}_N(w)`$ via the Deift-Zhou method . In the homogeneous case, i.e. when $`M=1`$, the system of Schlesinger equations we obtain reduces to a special case of Painlevé V equation. This result was obtained earlier in . The exact formulation of the results indicated above is presented in Theorem 1 in §4.
Our derivation of the differential equations for the function $`G_I(n;\{p_i\},t)`$ follows a scheme well known in soliton theory (see e.g. ) called the Zakharov-Shabat dressing method. We are able to apply this scheme since there exists a matrix Riemann-Hilbert problem associated to any Toeplitz determinant as was shown by Deift . For the reader’s convenience this is derived in §2.
The basic idea of the Riemann-Hilbert approach to Toeplitz determinants suggested in is a representation of a Toeplitz determinant $`D_n(\varphi )`$ as a Fredholm determinant of an integral operator acting on $`L_2(C)`$, $`C`$=unit circle, and belonging to a special integrable class which admits a Riemann-Hilbert representation . Borodin and Okounkov (see also for a simplified derivation and , for a particular case of $`\varphi `$) found a different Fredholm determinant representation for $`D_n(\varphi )`$. The Fredholm operator in this representation acts on $`l_2(\{n,n+1,\mathrm{}\})`$ which makes the representation quite suitable for the analysis of the large $`n`$ asymptotics of $`D_n(\varphi )`$ (see , ). Borodin subsequently observed that the discrete Fredholm representation of involves a discrete analog of the integrable kernels and can be supplemented by a discrete analog of the Riemann-Hilbert problem. (This is similar to the pure soliton constructions in the theory of integrable PDEs .)
We conclude this introduction by noting that our derivation of integrable PDEs for the Toeplitz determinant $`D_n(f_I)`$ can be applied to any Toeplitz determinant whose symbol $`\varphi `$ satisfies the condition,
$$\frac{d}{dz}\mathrm{log}\varphi (z)=\text{rational function of }z.$$
This is one place where the finite support of $`p`$ is crucial. It is an interesting open problem, particularly in light of , to remove this restriction.
## 2 Fredholm Determinant Representation of the Toeplitz Determinant and the Riemann-Hilbert Problem
Let $`\varphi (z)`$ be a continuous function on the unit circle $`C=\{|z|=1\}`$ oriented in the counterclockwise direction. Let $`n\text{N}`$ and denote by $`K_n(\varphi )`$ the integral operator acting on $`L_2(C)`$ with kernel
$$K_n(z,z^{})=\frac{z^n(z^{})^n1}{zz^{}}\frac{1\varphi (z^{})}{2\pi i}.$$
(2.1)
It was shown in that
$$D_n(\varphi )=det(1K_n(\varphi )),$$
(2.2)
where the determinant on the right is a Fredholm determinant taken in $`L_2(C)`$. (Note that $`K_n(z,z^{})`$ has no singularity at $`z=z^{}`$.) Equation (2.2) follows from the “geometric sum form” of the kernel $`K_n`$,
$$K_n(z,z^{})=\underset{k=0}{\overset{n1}{}}z^k\frac{1\varphi (z^{})}{2\pi i}(z^{})^{k1},$$
which shows that the Toeplitz matrix $`T_n(\varphi )`$ is essentially the matrix representation of the operator $`1K_n(\varphi )`$ in the basis $`\{z^k\}_{\mathrm{}<k<\mathrm{}}`$. (For more details see .)
The integral operator $`K_n(\varphi )`$ belongs to the class of integrable Fredholm operators , i.e., its kernel is of the form
$$K_n(z,z^{})=\frac{f^T(z)g(z^{})}{zz^{}},$$
where
$$f(z)=(f_1,f_2)^T=(z^n,1)^T$$
(2.3)
and
$$g(z)=(g_1,g_2)^T=(z^n,1)^T\frac{1\varphi (z)}{2\pi i}.$$
(2.4)
We require, so that there is no singularity on the diagonal of the kernel,
$$f^T(z)g(z)=0.$$
(2.5)
An important property of these operators is that the resolvent $`R_n=(1K)^11`$ also belongs to the same class (see again , , ). Precisely,
$$R_n(z,z^{})=\frac{F^T(z)G(z^{})}{zz^{}},$$
(2.6)
where
$$F_j=(1K_n)^1f_j,G_j=(1K_n^T)^1g_j,j=1,2.$$
The vector functions $`F`$ and $`G`$ can be in turn computed in terms of a certain matrix Riemann-Hilbert problem . Indeed, let us define (cf. ) the $`2\times 2`$ matrix valued function
$$Y(z)=I_CF(z^{})g^T(z^{})\frac{dz^{}}{z^{}z}zC.$$
(2.7)
Let $`Y_\pm (z)`$ denote the boundary values of the function $`Y(z)`$ on the contour $`C`$,
$$Y_\pm (z)=\underset{\genfrac{}{}{0pt}{}{z^{}z}{z^{}\pm \text{side}}}{lim}Y_+(z^{}).$$
From (2.7) it follows that
$$Y_+(z)Y_{}(z)=2\pi iF(z)g^T(z)$$
(2.8)
and hence (recall (2.5))
$$Y_+(z)f(z)=Y_{}(z)f(z).$$
Using this, the matrix identity
$$F(z^{})g^T(z^{})f(z)=f^T(z)g(z^{})F(z^{}),(\text{associativity of the matrix product})$$
and (2.7), we have
$$Y_\pm (z)f(z)=f(z)_Cf^T(z)g(z^{})F(z^{})\frac{dz^{}}{z^{}z}=f(z)+_CK(z,z^{})F(z^{})𝑑z^{}.$$
From the definition of $`F`$ it follows that
$$F(z)=Y_\pm (z)f(z).$$
(2.9)
This and (2.8) imply the jump equation,
$$Y_{}(z)=Y_+(z)(I+2\pi if^T(z)g(z)),zC.$$
(2.10)
This equation, supplemented by the obvious analytic properties of the Cauchy integral in (2.7), shows that the function $`Y(z)`$ solves the following $`2\times 2`$ matrix Riemann-Hilbert problem:
* $`Y(z)`$ is holomorphic for all $`zC`$,
* $`Y(\mathrm{})=I`$,
* $`Y_{}(z)=Y_+(z)H(z),zC`$,
where the jump matrix $`H`$ is
$`H(z)`$ $`=`$ $`I+2\pi if^T(z)g(z)`$ (2.11)
$`=`$ $`\left(\begin{array}{cc}2\varphi (z)& (\varphi (z)1)z^n\\ (1\varphi (z))z^n& \varphi (z)\end{array}\right).`$ (2.14)
These analytic properties determine $`Y`$ uniquely. To see this, we first observe that $`detH(z)1`$ implies that the scalar function $`detY(z)`$ has no jump on $`C`$; and hence, it is holomorphic and bounded on the whole complex plane. This together with the normalization condition at $`z=\mathrm{}`$ implies that $`detY(z)1`$. Suppose that $`\stackrel{~}{Y}(z)`$ is another solution. Since both the functions $`\stackrel{~}{Y}(z)`$ and $`Y(z)`$ satisfy the same jump condition across the contour $`C`$, the matrix ratio $`\stackrel{~}{Y}(z)Y^1(z)`$ has no jump across $`C`$. This means that $`\stackrel{~}{Y}(z)Y^1(z)\text{constant}`$, and from the condition at $`z=\mathrm{}`$ we actually have that $`\stackrel{~}{Y}(z)Y^1(z)I`$. The uniqueness now follows.
Since $`Y`$ is the unique solution of the Riemann-Hilbert problem, one can now reconstruct the resolvent $`R`$ using (2.9) and the similarly derived identity
$$G(z)=(Y_\pm ^T)^1(z)g(z),$$
(2.15)
for $`G`$. We shall refer to this Riemann-Hilbert problem as the $`Y`$-RH problem.
Following the $`Y`$-RH problem can be transformed to an equivalent Riemann-Hilbert problem which is directly connected with the polynomials on the circle $`C`$ orthogonal with respect to the (generally complex) weight $`\varphi (e^{i\theta })`$. To this end we first note that since the entries of $`f`$ are polynomials in $`z`$, (2.9) implies that $`F`$ is an entire function of $`z`$. Since $`Y(z)I`$ as $`z\mathrm{}`$, it follows in fact that $`F`$ is polynomial,
$$F(z)=\left(\begin{array}{c}P_n(z)\\ Q_{n1}(z)\end{array}\right),P_n(z)=z^n+\mathrm{},Q_{n1}(z)=q_{n1}z^{n1}+\mathrm{},$$
(2.16)
for some constant $`q_{n1}`$. On the other hand, denoting by $`Y_j`$ the $`j`$-th column of the matrix $`Y`$, we obtain from the jump equation (2.10) (or, more precisely, from the equation $`Y_+=Y_{}H^1`$) that
$`Y_{1+}(z)`$ $`=`$ $`Y_{}(z)\left(\begin{array}{c}\varphi (z)\\ (\varphi (z)1)z^n\end{array}\right)`$ (2.19)
$`=`$ $`z^nY_{}(z)\left(\begin{array}{c}0\\ 1\end{array}\right)+\varphi (z)z^nY_{}(z)\left(\begin{array}{c}z^n\\ 1\end{array}\right)`$
$`=`$ $`z^nY_2(z)+\varphi (z)z^nY_{}(z)f(z)`$
$`=`$ $`z^nY_2(z)+\varphi (z)z^nF(z).`$ (2.25)
Define
$$J(z)=\{\begin{array}{cc}Y_1(z),\hfill & |z|<1,\hfill \\ z^nY_2(z),\hfill & |z|>1,\hfill \end{array}$$
and consider the $`2\times 2`$ matrix function
$$Z(z)=\sigma _3(F(z),J(z))\sigma _3,\sigma _3=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right).$$
(2.26)
The function $`Z`$ is analytic outside of $`C`$, and it has the following asymptotic behavior as $`z\mathrm{}`$:
$$Z(z)=\left(I+O\left(\frac{1}{z}\right)\right)\left(\begin{array}{cc}z^n& 0\\ 0& z^n\end{array}\right).$$
(2.27)
For the jump relation on the contour $`C`$ we have from (2.25),
$$Z_+(z)=\sigma _3(F(z),Y_{1+}(z))\sigma _3=\sigma _3(F(z),z^nY_2(z)\varphi (z)z^nF(z))\sigma _3$$
$$=\sigma _3(F(z),z^nY_2(z))\left(\begin{array}{cc}1& \varphi (z)z^n\\ 0& 1\end{array}\right)\sigma _3=\sigma _3(F(z),z^nY_2(z))\sigma _3\sigma _3\left(\begin{array}{cc}1& \varphi (z)z^n\\ 0& 1\end{array}\right)\sigma _3$$
$$=Z_{}(z)\left(\begin{array}{cc}1& \varphi (z)z^n\\ 0& 1\end{array}\right).$$
Summarizing the analytic properties of $`Z`$ we conclude that it solves the following Riemann-Hilbert problem:
* $`Z`$ is holomorphic for all $`zC`$,
* $`Z(z)z^{n\sigma _3}I,z\mathrm{}`$,
* $`Z_+(z)=Z_{}(z)S(z),zC`$,
where the jump matrix $`S(z)`$ is
$$S(z)=\left(\begin{array}{cc}1& z^n\varphi (z)\\ 0& 1\end{array}\right).$$
(2.28)
We shall refer to this Riemann-Hilbert problem as $`Z`$-RH problem. As in the $`Y`$-RH problem, the solution of the $`Z`$-RH problem is unique. Indeed, assuming that $`\stackrel{~}{Z}`$ is another solution, we introduce the matrix ratio $`X:=\stackrel{~}{Z}Z^1`$. By the same reasoning as in the case of the $`Y`$-RH problem, we conclude that $`X`$ is entire. Since
$$X(z)=(\stackrel{~}{Z}(z)z^{n\sigma _3})(z^{n\sigma _3}Z^1(z))I\text{as}z\mathrm{},$$
it follows that $`XI`$; and hence, that $`Z`$ is unique. We note that $`Y`$ (and hence the resolvent $`R`$) can be reconstructed from $`Z`$ using (2.26). It also should be pointed out that the existence of the solution of the $`Z`$-RH problem (as well as of the $`Y`$-RH problem) is equivalent to the nondegeneracy of the Toeplitz matrix $`T_n(\varphi )`$, i.e. to the inequality
$$D_n(\varphi )0,$$
which we always assume.
Remark. There is a more direct and elegant way to pass to the $`Z`$-RH problem which was pointed out by the referee of this paper. One first notes that the jump matrix $`H`$ admits the factorization,
$$H(z)=\left(\begin{array}{cc}z^n& 1\\ 1& 0\end{array}\right)\left(\begin{array}{cc}1& z^n\varphi (z)\\ 0& 1\end{array}\right)\left(\begin{array}{cc}z^n& 0\\ 1& z^n\end{array}\right),$$
which then suggests the definition
$$\stackrel{~}{Y}(z)=\{\begin{array}{cc}Y(z)\left(\begin{array}{cc}z^n& 1\\ 1& 0\end{array}\right),\hfill & |z|<1,\hfill \\ & \\ Y(z)\left(\begin{array}{cc}z^n& 0\\ 1& z^n\end{array}\right)Y(z)\left(\begin{array}{cc}1& 0\\ z^n& 1\end{array}\right)z^{n\sigma _3},\hfill & |z|>1,\hfill \end{array}$$
(2.29)
so that the function $`\stackrel{~}{Y}`$ would satisfy the Riemann-Hilbert problem,
* $`\stackrel{~}{Y}`$ is holomorphic for all $`zC`$,
* $`\stackrel{~}{Y}(z)z^{n\sigma _3}I,z\mathrm{}`$,
* $`\stackrel{~}{Y}_{}(z)=\stackrel{~}{Y}_+(z)\left(\begin{array}{cc}1& z^n\varphi (z)\\ 0& 1\end{array}\right),zC`$.
The function $`Z(z)`$ is related to $`\stackrel{~}{Y}(z)`$ by
$$Z(z)=\sigma _3\stackrel{~}{Y}(z)\sigma _3.$$
(2.30)
We conclude this section by summarizing the relation of the $`Z`$-RH problem to the orthogonal polynomials on $`C`$ with respect to the (generally complex) weight $`\varphi `$. This relation is due to Deift (see also ).<sup>4</sup><sup>4</sup>4The $`Z`$-RH problem is the analog for polynomials on the circle of the Riemann-Hilbert problem derived in for polynomials which are orthogonal with respect to an exponential weight on the line (see also ). Let $`\{P_k(z)\}_{k=0,1,\mathrm{}}`$ denote the system of the monic polynomials defined by
$$P_k(z)=z^k+\mathrm{},$$
$$_CP_n(z)\overline{P}_m(z)\varphi (z)\frac{dz}{iz}=h_n\delta _{nm},nm,$$
where bar denotes complex conjugation. Similarly, introduce a second system of polynomials, $`\{P_k^{}(z)\}_{k=0,1,\mathrm{}}`$, by replacing $`\varphi `$ with $`\overline{\varphi }`$ in the definition of $`P_n`$. Suppose now that
$$D_k(\varphi )0,k=1,\mathrm{},n+1.$$
(2.31)
Then (see ) both the sets of polynomials $`\{P_k\}_{k=0,1,\mathrm{},n}`$ and $`\{P_k^{}\}_{k=0,1,\mathrm{},n}`$ exist, and the normalization constants $`h_k`$ and $`h_k^{}`$, $`k=1,\mathrm{},n`$, are all nonzero. In fact, we have the explicit representations
$$P_n(z)=\frac{D_{n+1}(\varphi |z)}{D_n(\varphi )},h_n=2\pi \frac{D_{n+1}(\varphi )}{D_n(\varphi )},P_n^{}(z)=\frac{D_{n+1}(\overline{\varphi }|z)}{D_n(\overline{\varphi })},h_n^{}=2\pi \frac{D_{n+1}(\overline{\varphi })}{D_n(\overline{\varphi })},$$
(2.32)
where $`D_{n+1}(\varphi |z)`$ denotes the Toeplitz determinant $`D_{n+1}(\varphi )`$ whose last row is replaced by the row $`(1,z,z^2,\mathrm{},z^n)`$. If we define
$$Q_k=\frac{2\pi }{\overline{h_k^{}}}\overline{P}_k^{}(1/\overline{z})z^k,$$
(2.33)
and
$$Z(z)=\left(\begin{array}{cc}P_n(z)& \frac{1}{2\pi i}_CP_n(z^{})(z^{})^n\varphi (z^{})\frac{dz^{}}{z^{}z}\\ & \\ Q_{n1}(z)& \frac{1}{2\pi i}_CQ_{n1}(z^{})(z^{})^n\varphi (z^{})\frac{dz^{}}{z^{}z}\end{array}\right),$$
(2.34)
then it is a calculation to show that this $`Z`$ defines a (unique) solution of the $`Z`$-RH problem (cf. , , and ). Indeed, the analyticity in $`\text{C}C`$ and the jump condition follow from the basic properties of Cauchy integral, and the asymptotic condition at $`z=\mathrm{}`$ is equivalent to the fact that the polynomials $`P_n`$ and $`P_{n1}^{}`$ are monic orthogonal polynomials with the weights $`\varphi (z)dz`$ and $`\overline{\varphi }(z)dz`$, respectively.
## 3 Toeplitz Determinants as Integrable Systems
### 3.1 Universal recursion relation
In this section $`\varphi `$ will be an arbitrary continuous function with Fourier coefficients $`\varphi _j`$. We assume that the associated Toeplitz matrix $`T_n(\varphi )`$ is invertible. Then the corresponding matrix RH problem is uniquely solvable, and the following equation connects the Toeplitz determinant $`D_n(\varphi )`$ with the solution $`Z`$ of the Riemann-Hilbert problem,
$$\frac{D_{n+1}}{D_n}=Z_{12}(0),$$
(3.1)
where $`Z_{ij}`$, $`i,j=1,2`$, denotes the entries of matrix $`Z`$. Indeed, using (2.32) we have that
$$\frac{D_{n+1}}{D_n}=\frac{1}{2\pi }h_n,$$
On the other hand, (2.34) gives
$$Z_{12}(0)=\frac{1}{2\pi }h_n,$$
and (3.1) follows.
Remark. One can prove (3.1) using only the connection with the integrable operator $`K_n(\varphi )`$ introduced in (2.1). To see this first note
$`K_{n+1}(z,z^{})`$ $`=`$ $`{\displaystyle \frac{(z/z^{})^{n+1}1}{zz^{}}}{\displaystyle \frac{1\varphi (z^{})}{2\pi i}}`$
$`=`$ $`{\displaystyle \frac{1}{z^{}}}{\displaystyle \underset{k=0}{\overset{n}{}}}\left({\displaystyle \frac{z}{z^{}}}\right)^k{\displaystyle \frac{1\varphi (z^{})}{2\pi i}}`$
$`=`$ $`K_n(z,z^{})+f_1(z)g_1(z^{}){\displaystyle \frac{1}{z^{}}},`$
where $`f`$ and $`g`$ are defined in (2.3) and (2.4), respectively. Attaching to the functions $`f`$, $`g`$ superscript “$`n`$” to denote their $`n`$ dependence, we can rewrite the last equation as an operator equation
$$K_{n+1}=K_n+f_1^ng_1^{n+1},$$
(3.2)
where the symbol $`ab`$ denotes the integral operator with kernel $`a(z)b(z^{})`$. Recalling the definition of $`F`$, (2.6), it follows from (3.2) (cf. ) that
$`det(1K_{n+1})`$ $`=`$ $`det(1K_n)det(1[(1K_n)^1f_1^n]g_1^{n+1})`$
$`=`$ $`det(1K_n)det(1F_1^ng_1^{n+1})`$
$`=`$ $`det(1K_n)(1\text{trace}F_1^ng_1^{n+1})`$
$`=`$ $`det(1K_n)\left(1{\displaystyle _C}F_1(z)g_1(z){\displaystyle \frac{dz}{z}}\right)`$
where $`F_1:=F_1^n`$ and $`g_1:=g_1^n`$. Thus (see also (2.2))
$$\frac{D_{n+1}}{D_n}=\frac{detK_{n+1}}{detK_n}=1_CF_1(z)g_1(z)\frac{dz}{z}.$$
(3.3)
Recalling (2.7), we rewrite (3.3) as
$$\frac{D_{n+1}}{D_n}=Y_{11}(0),$$
which together with (2.26) yields (3.1).
### 3.2 Differentiation formulas
Here we restrict to the symbol
$$\varphi (z)=e^{tz}\underset{\alpha =1}{\overset{M}{}}\left(\frac{zr_\alpha }{z}\right)^{k_\alpha },$$
(3.4)
where $`r_\alpha :=p_{i_\alpha }`$, and we recall (see §1) that $`i_\alpha =\mathrm{max}𝒜_\alpha `$, $`k_\alpha =|𝒜_\alpha |`$, and
$$𝒜=\underset{\alpha =1}{\overset{M}{}}𝒜_\alpha $$
is the decomposition of the alphabet $`𝒜`$ into subsets $`𝒜_1,𝒜_2,\mathrm{},𝒜_M`$ such that $`p_i=p_j`$ if and only if $`i`$ and $`j`$ belong the same $`𝒜_\alpha `$. We also recall that
$$\underset{\alpha =1}{\overset{M}{}}k_\alpha =k,$$
and
$$1>p_1p_2\mathrm{}p_k>0,\underset{j=1}{\overset{k}{}}p_j=1$$
(3.5)
denote the probabilities assigned to the letters $`i=1,2,\mathrm{},k`$, in the alphabet $`𝒜`$. Note that from the probabilistic conditions (3.5) it follows that
$$1<r_\alpha <0,\alpha =1,\mathrm{},M,r_\alpha r_\beta ,\alpha \beta ,$$
(3.6)
and
$$\underset{\alpha =1}{\overset{M}{}}k_\alpha r_\alpha =1.$$
(3.7)
The symbols $`f_I`$ and $`\varphi `$ are related by
$$f_I(z)=\varphi (1/z),$$
and therefore; the corresponding Toeplitz matrices are mutually transpose. Thus
$$G_I(n;\{p_i\},t)=D_n(\varphi ).$$
In what follows, we will write $`T_n(t)`$, $`K_n(t)`$ and $`D_n(t)`$ for $`T_n(\varphi )`$, $`K_n(\varphi )`$ and $`D_n(\varphi )`$, respectively; or $`T_n(\{p_i\},t)`$, $`K_n(\{p_i\},t)`$ and $`D_n(\{p_i\},t)`$ if the dependence on $`p_1,\mathrm{},p_k`$ is of interest.
We shall derive the differential formulas for the Toeplitz determinant $`D_n(t)`$ with respect to the variables $`t`$ and $`r_\alpha ,\alpha =1,\mathrm{},M`$ assuming that the latter are subject to restriction (3.6) only, i.e. we only will assume that
$$1<r_\alpha <0,\alpha =1,\mathrm{},M,r_\alpha r_\beta ,\alpha \beta .$$
The integers $`k_\alpha `$ will be kept constant. This means that when vary the points $`r_\alpha `$ we do not assume restriction (3.7) to hold. We will begin with the $`t`$ \- derivative.
Since $`\varphi /t=z\varphi `$,
$$\frac{}{t}K_n(z,z^{})=\frac{(z/z^{})^n1}{zz^{}}(z^{})\frac{\varphi (z^{})}{2\pi i}$$
$$=\frac{1\varphi (z^{})}{2\pi i}+z\frac{(z/z^{})^{n1}1}{zz^{}}\frac{1\varphi (z^{})}{2\pi i}\frac{(z/z^{})^n1}{zz^{}}\frac{z^{}}{2\pi i}.$$
(3.8)
Let $`\mathrm{\Lambda }`$ be the integral operator with the kernel
$$\mathrm{\Lambda }(z,z^{})=\frac{(z/z^{})^n1}{zz^{}}\frac{z^{}}{2\pi i}.$$
Consider the operator product $`\mathrm{\Lambda }K_n`$:
$$(\mathrm{\Lambda }K_n)(z,z^{})=\frac{1\varphi (z^{})}{2\pi i}_C\frac{(z/w)^n1}{zw}w\frac{(w/z^{})^n1}{wz^{}}\frac{dw}{2\pi i}$$
$$=\frac{1\varphi (z^{})}{2\pi i}_C\underset{j,l=0}{\overset{n1}{}}\left(\frac{z}{w}\right)^l\left(\frac{w}{z^{}}\right)^j(z^{})^1\frac{dw}{2\pi i}=\frac{1\varphi (z^{})}{2\pi i}\underset{jl=1}{}z^l(z^{})^{j1}$$
$$=\frac{1\varphi (z^{})}{2\pi i}\underset{j=0}{\overset{n2}{}}z^{j+1}(z^{})^{j1}=\frac{1\varphi (z^{})}{2\pi i}z\frac{(z/z^{})^{n1}1}{zz^{}}.$$
Recalling the definitions of $`f`$ and $`g`$, (2.3) and(2.4), (3.8) can be written compactly as
$$\frac{}{t}K_n=f_2g_2\mathrm{\Lambda }(1K_n).$$
(3.9)
From this formula we see that (cf. the derivation of (3.3))
$`{\displaystyle \frac{}{t}}\mathrm{log}D_n(t)`$ $`=`$ $`\text{trace}\left((1K_n)^1{\displaystyle \frac{}{t}}K_n\right)`$ (3.10)
$`=`$ $`\text{trace}F_2g_2+\text{trace}\mathrm{\Lambda }`$
$`=`$ $`{\displaystyle _C}F_2(z)g_2(z)𝑑z,`$
where we used the fact that
$$\text{trace}\mathrm{\Lambda }=\frac{n}{2\pi i}_C𝑑z=0.$$
Recalling (2.7) we convert (3.10) into the identity
$$\frac{}{t}\mathrm{log}D_n(t)=\text{res}_{z=\mathrm{}}(Y_{22}(z)),$$
which in terms of the $`Z`$-function is
$$\frac{}{t}\mathrm{log}D_n(t)=\text{res}_{z=\mathrm{}}(z^nZ_{22}(z)),$$
or equivalently,
$$\frac{}{t}\mathrm{log}D_n(t)=(\mathrm{\Gamma }_1)_{22},$$
(3.11)
where the matrix $`\mathrm{\Gamma }_1=\mathrm{\Gamma }_1(\{p_i\},t)`$ is defined by the expansion,
$$Z(z)=\left(I+\underset{j=1}{\overset{\mathrm{}}{}}\frac{\mathrm{\Gamma }_j}{z^j}\right)z^{n\sigma _3},|z|>1.$$
(3.12)
Remark. In the basis $`\{z^n\}_{n=\mathrm{}}^{\mathrm{}}`$, (3.9) coincides with (3.22) of .
Equation (3.11) is the $`t`$ \- differentiation formula, i.e. it gives an expression of the $`t`$ \- derivative of $`\mathrm{log}D_n`$ in terms of the solution $`Z`$ of the Riemann-Hilbert problem. We shall proceed now with the derivation of the $`r_\alpha `$ \- differentiation formula. Since (we recall that $`r_\alpha `$ are assumed independent and that the $`k_\alpha `$ are kept constant)
$$\frac{}{r_\alpha }\varphi =\frac{k_\alpha }{zr_\alpha }\varphi ,$$
the $`r_\alpha `$ \- analog of (3.8) reads
$$\frac{}{r_\alpha }K_n(z,z^{})=k_\alpha \frac{(z/z^{})^n1}{zz^{}}\frac{1}{z^{}r_\alpha }\frac{\varphi (z^{})}{2\pi i}$$
$$=\frac{k_\alpha }{2\pi i}\frac{(z/z^{})^n1}{zz^{}}\frac{1}{z^{}r_\alpha }k_\alpha \frac{(z/z^{})^n1}{zz^{}}\frac{1}{z^{}r_\alpha }\frac{1\varphi (z^{})}{2\pi i}.$$
(3.13)
Introducing the integral operator $`\mathrm{\Lambda }_\alpha `$ with the kernel,
$$\mathrm{\Lambda }_\alpha (z,z^{})=\frac{k_\alpha }{2\pi i}\frac{(z/z^{})^n1}{zz^{}}\frac{1}{z^{}r_\alpha },$$
we consider again the operator product $`\mathrm{\Lambda }_\alpha K_n`$. The residue type calculations, similar to the ones used in the $`t`$ \- case, yield the equation
$$(\mathrm{\Lambda }_\alpha K_n)(z,z^{})=\frac{k_\alpha }{2\pi i}(1\varphi (z^{}))\left[\frac{(z/z^{})^n1}{(z^{}z)(r_\alpha z)}+\frac{(r_\alpha /z^{})^n1}{(r_\alpha z)(r_\alpha z^{})}\right],$$
which in turn impies that (3.13) can be rewritten as
$$\frac{}{r_\alpha }K_n(z,z^{})=\frac{k_\alpha }{2\pi i}\frac{1\varphi (z^{})}{(r_\alpha z)(r_\alpha z^{})}\left[\left(\frac{r_\alpha }{z^{}}\right)^n\left(\frac{z}{z^{}}\right)^n\right]+[\mathrm{\Lambda }_\alpha (1K_n)](z,z^{}).$$
With the help of the vector functions,
$$\stackrel{~}{f}(z):=\frac{1}{zr_\alpha }f(z),\stackrel{~}{g}(z):=\frac{1}{zr_\alpha }g(z),$$
the last equation can be transformed into the following compact form (cf. (3.9))
$$\frac{}{r_\alpha }K_n=k_\alpha r_\alpha ^n\stackrel{~}{f}_2\stackrel{~}{g}_1k_\alpha \stackrel{~}{f}_1\stackrel{~}{g}_1+\mathrm{\Lambda }_\alpha (1K_n).$$
(3.14)
Let the vector function $`\stackrel{~}{F}(z)=(\stackrel{~}{F}_1(z),\stackrel{~}{F}_2(z))^T`$ be defined by the equation
$$\stackrel{~}{F}_j:=(1K_n)^1\stackrel{~}{f}_j,j=1,2.$$
We observe that
$$\stackrel{~}{F}(z)=\frac{1}{zr_\alpha }Y^1(r_\alpha )F(z),$$
(3.15)
where the matrix function $`Y(z)`$ is the solution of the $`Y`$-RH problem corresponding to $`D_n(t)`$. Indeed by the definition of the vector function $`F(z)`$ (see (2.6)) its component $`F_j(z)`$ satisfies the integral equation,
$$F_j(z)_CK_n(z,z^{})F_j(z^{})𝑑z^{}=f_j(z).$$
(3.16)
Dividing both sides of this equation by $`(zr_\alpha )`$, using the formula
$$K_n(z,z^{})=\frac{f^T(z)g(z^{})}{zz^{}},$$
and simple algebra, we can rewrite (3.16) as an equation for the ratio $`F_j/(zr_\alpha )`$:
$$\left[\frac{F_j(z)}{zr_\alpha }\right]_CK_n(z,z^{})\left[\frac{F_j(z^{})}{z^{}r_\alpha }\right]𝑑z^{}+_C\stackrel{~}{f}^T(z)g(z^{})\left[\frac{F_j(z^{})}{z^{}r_\alpha }\right]𝑑z^{}=\stackrel{~}{f}_j(z).$$
By applying the operator $`(1K_n)^1`$ to the both sides of this equation it can be transformed into the relation,
$$\left[\frac{F_j(z)}{zr_\alpha }\right]+_C\stackrel{~}{F}^T(z)g(z^{})\left[\frac{F_j(z^{})}{z^{}r_\alpha }\right]𝑑z^{}=\stackrel{~}{F}_j(z),$$
or
$$\frac{1}{zr_\alpha }F_j(z)+\underset{i=1}{\overset{2}{}}\stackrel{~}{F}_i(z)_Cg_i(z^{})F_j(z^{})\frac{dz^{}}{z^{}r_\alpha }=\stackrel{~}{F}_j(z).$$
The last equation in turn can be viewed as the linear algebraic system for the vector $`\stackrel{~}{F}(z)`$,
$$\stackrel{~}{F}_j(z)\underset{i=1}{\overset{2}{}}A_{ji}\stackrel{~}{F}_i(z)=\frac{1}{zr_\alpha }F_j(z),j=1,2,$$
(3.17)
where the matrix $`A`$ is given by the formula,
$$A_{ji}=_CF_j(z^{})g_i(z^{})\frac{dz^{}}{z^{}r_\alpha }.$$
Equation (3.15) follows directly from (3.17) in virtue of definition (2.7) of the matrix function $`Y(z)`$.
We now able to finish the derivation of the $`r_\alpha `$ \- differentiation formula for the Toeplitz determinant $`D_n(t)`$. In fact from (3.14) it follows that (cf. the derivation of (3.10))
$`{\displaystyle \frac{}{r_\alpha }}\mathrm{log}D_n(t)`$ $`=`$ $`\text{trace}\left((1K_n)^1{\displaystyle \frac{}{r_\alpha }}K_n\right)`$ (3.18)
$`=`$ $`k_\alpha r_\alpha ^n\text{trace}\stackrel{~}{F}_2\stackrel{~}{g}_1+k_\alpha \text{trace}\stackrel{~}{F}_1\stackrel{~}{g}_1\text{trace}\mathrm{\Lambda }_\alpha `$
$`=`$ $`k_\alpha r_\alpha ^n{\displaystyle _C}\stackrel{~}{F}_2(z)g_1(z){\displaystyle \frac{dz}{zr_\alpha }}+k_\alpha {\displaystyle _C}\stackrel{~}{F}_1(z)g_1(z){\displaystyle \frac{dz}{zr_\alpha }},`$
where, similar to the $`t`$ \- derivative case, we used the fact that
$$\text{trace}\mathrm{\Lambda }_\alpha =\frac{nk_\alpha }{2\pi i}_C\frac{dz}{z(zr_\alpha )}=0.$$
Using now (3.15) and the fact that $`detY(z)1`$ we derive from (3.18) that
$`{\displaystyle \frac{}{r_\alpha }}\mathrm{log}D_n(t)`$ $`=`$ $`k_\alpha r_\alpha ^nY_{21}(r_\alpha ){\displaystyle _C}F_1(z)g_1(z){\displaystyle \frac{dz}{(zr_\alpha )^2}}k_\alpha r_\alpha ^nY_{11}(r_\alpha ){\displaystyle _C}F_2(z)g_1(z){\displaystyle \frac{dz}{(zr_\alpha )^2}}`$
$`+k_\alpha Y_{22}(r_\alpha ){\displaystyle _C}F_1(z)g_1(z){\displaystyle \frac{dz}{(zr_\alpha )^2}}k_\alpha Y_{12}(r_\alpha ){\displaystyle _C}F_2(z)g_1(z){\displaystyle \frac{dz}{(zr_\alpha )^2}},`$
or
$`{\displaystyle \frac{}{r_\alpha }}\mathrm{log}D_n(t)`$ $`=`$ $`k_\alpha r_\alpha ^nY_{21}(r_\alpha )Y_{11}^{}(r_\alpha )+k_\alpha r_\alpha ^nY_{11}(r_\alpha )Y_{21}^{}(r_\alpha )`$ (3.19)
$`k_\alpha Y_{22}(r_\alpha )Y_{11}^{}(r_\alpha )+k_\alpha Y_{12}(r_\alpha )Y_{21}^{}(r_\alpha ),`$
where we use the notation,
$$Y_{ij}^{}(r_\alpha ):=\frac{Y_{ij}(z)}{z}|_{z=r_\alpha },i,j=1,2.$$
Equation (3.19) can be also rewritten as
$`{\displaystyle \frac{}{r_\alpha }}\mathrm{log}D_n(t)`$ $`=`$ $`k_\alpha \left(r_\alpha ^nY_{21}(r_\alpha )+Y_{22}(r_\alpha )\right)Y_{11}^{}(r_\alpha )`$ (3.20)
$`+k_\alpha \left(r_\alpha ^nY_{11}(r_\alpha )+Y_{12}(r_\alpha )\right)Y_{21}^{}(r_\alpha ),`$
which in turn can be transformed into an expression of the $`\mathrm{log}D_n(t)/r_\alpha `$ in terms of the $`Z`$-function. Indeed recalling formulae (2.26), (2.9), and (2.3), we see that inside the unit circle $`C`$ the following equation takes place,
$$Z(z)=\left(\begin{array}{cc}z^nY_{11}(z)+Y_{12}(z)& Y_{11}(z)\\ z^nY_{21}(z)Y_{22}(z)& Y_{21}(z)\end{array}\right),|z|<1,$$
so that (3.20) can be converted into the $`r_\alpha `$-differentiation formula,
$$\frac{}{r_\alpha }\mathrm{log}D_n(t)=k_\alpha \left(Z_{11}(r_\alpha )Z_{22}^{}(r_\alpha )Z_{21}(r_\alpha )Z_{12}^{}(r_\alpha )\right).$$
(3.21)
### 3.3 Schlesinger equations
In this section we show that $`D_n(t)`$ is the Jimbo-Miwa-Ueno $`\tau `$-function of the generalized Schlesinger system of nonlinear differential equations describing the isomonodromy deformations of the $`2\times 2`$ matrix linear ODE which has $`M+1`$ simple poles in the finite complex plane and one Poincaré index 1 irregular singular point at infinity. We will also evaluate the relevant monodromy data that single out the $`D_n(t)`$ from all the other solutions of the Schlesinger system. In the uniform case, when all $`p_i`$ are equal, the system reduces to the particular case of Painlevé V equation, i.e. we are back to the uniform result of .
Define
$$\mathrm{\Phi }^0(z)=e^{\frac{tz}{2}\sigma _3}\left(\begin{array}{cc}1& 0\\ 0& z^n\psi ^1(z)\end{array}\right)$$
(3.22)
where
$$\psi (z)=\underset{\alpha =1}{\overset{M}{}}\left(\frac{zr_\alpha }{z}\right)^{k_\alpha }.$$
(3.23)
We note that that the product
$$e^{tz}\psi (z):=\varphi (z)$$
is our symbol, i.e. the function defined in (3.4). We also note that $`\mathrm{\Phi }^0`$ is analytic and invertible in $`\text{C}\{0,r_1,\mathrm{},r_M\}`$, and that it satisfies the linear differential equations
$`\mathrm{\Phi }_z^0(z)`$ $`=`$ $`\mathrm{\Omega }(z)\mathrm{\Phi }^0(z),`$ (3.24)
$`\mathrm{\Phi }_t^0(z)`$ $`=`$ $`{\displaystyle \frac{z}{2}}\sigma _3\mathrm{\Phi }^0(z),`$ (3.25)
$`\mathrm{\Phi }_{r_\alpha }^0(z)`$ $`=`$ $`{\displaystyle \frac{k_\alpha }{zr_\alpha }}\left(\begin{array}{cc}0& 0\\ 0& 1\end{array}\right)\mathrm{\Phi }^0(z),`$ (3.28)
where $`\mathrm{\Omega }`$ is the rational matrix function
$$\mathrm{\Omega }(z)=\frac{t}{2}\sigma _3+\frac{n+k}{z}\left(\begin{array}{cc}0& 0\\ 0& 1\end{array}\right)\underset{\alpha =1}{\overset{M}{}}\frac{k_\alpha }{zr_\alpha }\left(\begin{array}{cc}0& 0\\ 0& 1\end{array}\right).$$
(3.29)
(Subscripts on $`\mathrm{\Phi }^0`$ denote differentiation.) Introduce
$$\mathrm{\Phi }(z)=Z(z)\mathrm{\Phi }^0(z),$$
(3.30)
where $`Z`$ is the solution of the $`Z`$-RH problem corresponding to our symbol $`\varphi `$, and consider the logarithmic derivative
$$B(z):=\mathrm{\Phi }_z(z)\mathrm{\Phi }^1(z).$$
(3.31)
The key observation is that $`B`$ is continuous across the contour $`C`$. Indeed, the $`Z`$-jump matrix $`S`$ (see (2.28)) admits the following factorization,
$$S(z)=\mathrm{\Phi }^0(z)\left(\begin{array}{cc}1& 1\\ 0& 1\end{array}\right)(\mathrm{\Phi }^0(z))^1,$$
(3.32)
so that the $`\mathrm{\Phi }`$-jump matrix does not depend on $`z`$. In fact we have
$$\mathrm{\Phi }_+(z)=\mathrm{\Phi }_{}(z)\left(\begin{array}{cc}1& 1\\ 0& 1\end{array}\right),zC.$$
(3.33)
This implies that
$$B_+(z)=B_{}(z),zC,$$
and hence the function $`B(z)`$ is an analytic function on $`\text{C}\{0,r_1,\mathrm{},r_M\}`$. We also recall that the only conditions which we impose on the points $`r_\alpha `$ are the inequalities (3.6), i.e.,
$$1<r_\alpha <0,\alpha =1,\mathrm{},M,r_\alpha r_\beta ,\alpha \beta .$$
(3.34)
We now calculate the principal part of $`B`$ at each of its singular points. Since $`Z`$ is holomorphic and invertible inside of $`C`$, it follows from
$$B(z)=Z(z)\mathrm{\Omega }(z)Z^1(z)+Z_z(z)Z^1(z)$$
(3.35)
that in a neighborhood of $`z=r_\alpha `$,
$$B(z)=\frac{k_\alpha }{zr_\alpha }Z(r_\alpha )\left(\begin{array}{cc}0& 0\\ 0& 1\end{array}\right)Z^1(r_\alpha )+\underset{j=0}{\overset{\mathrm{}}{}}b_j^\alpha (zr_\alpha )^j.$$
(3.36)
Likewise in a neighborhood of $`z=0`$, (3.35) and (3.29) imply that
$$B(z)=\frac{n+k}{z}Z(0)\left(\begin{array}{cc}0& 0\\ 0& 1\end{array}\right)Z^1(0)+\underset{j=0}{\overset{\mathrm{}}{}}b_j^0z^j.$$
(3.37)
Finally, from (3.35) and the Laurent expansion (3.12) we obtain the power series of $`B`$ at $`\mathrm{}`$,
$$B(z)=\frac{t}{2}\sigma _3+\underset{j=1}{\overset{\mathrm{}}{}}b_j^{\mathrm{}}z^j.$$
(3.38)
Equations (3.36)–(3.38) imply that $`B`$ is a rational function,
$$B(z)=\frac{t}{2}\sigma _3+\frac{B_0}{z}+\underset{\alpha =1}{\overset{M}{}}\frac{B_\alpha }{zr_\alpha },$$
(3.39)
with the matrix residues given by
$$B_0=Z(0)\left(\begin{array}{cc}0& 0\\ 0& n+k\end{array}\right)Z^1(0),$$
(3.40)
$$B_\alpha =Z(r_\alpha )\left(\begin{array}{cc}0& 0\\ 0& k_\alpha \end{array}\right)Z^1(r_\alpha ),\alpha =1,\mathrm{},M.$$
(3.41)
Thus from (3.31) we conclude that $`\mathrm{\Phi }`$ satisfies the linear differential equation,
$$\mathrm{\Phi }_z(z)=B(z)\mathrm{\Phi }(z),$$
(3.42)
with the coefficient matrix $`B`$ determined by (3.39)–(3.41).
Remark. In soliton theory (see ), the method that we used to derive (3.42) is called Zakharov-Shabat dressing of the vacuum equation (3.24). We also note that, as is common in the analysis of soliton equations, we have moved the exponential factor $`e^{\pm z/2}`$ to the asymptotic condition at $`z=\mathrm{}`$.
Let us now dress the $`t`$-vacuum equation (3.25), i.e. consider the $`t`$-logarithmic derivative of $`\mathrm{\Phi }`$
$$V(z):=\mathrm{\Phi }_t(z)\mathrm{\Phi }^1(z).$$
(3.43)
The $`\mathrm{\Phi }`$-jump matrix (3.33) does not depend on $`t`$ as well. Hence
$$V_+(z)=V_{}(z),zC,$$
and $`V`$ is analytic on $`\text{C}\{0,r_1,\mathrm{},r_m\}`$. In fact, since
$$V(z)=Z(z)\frac{z}{2}\sigma _3Z^1(z)+Z_t(z)Z^1(z),$$
(3.44)
(cf. (3.35)) and $`Z`$ is holomorphic at the points $`\{0,r_1,\mathrm{},r_M\}`$, we conclude that $`V`$ is entire. Moreover, from the expansion (3.12) we have that
$$V(z)=\frac{z}{2}\sigma _3+\frac{1}{2}[\sigma _3,\mathrm{\Gamma }_1]+\underset{j=1}{\overset{\mathrm{}}{}}v_jz^j,|z|>1,$$
and hence
$$V(z)=\frac{z}{2}\sigma _3+\frac{1}{2}[\sigma _3,\mathrm{\Gamma }_1].$$
(3.45)
($`[L,M]:=LMML`$.)
This in turn yields the $`t`$-equation for $`\mathrm{\Phi }`$,
$$\mathrm{\Phi }_t(z)=V(z)\mathrm{\Phi }(z),$$
(3.46)
where the coefficient matrix $`V`$ is defined by the equations,
$`V(z)`$ $`=`$ $`{\displaystyle \frac{z}{2}}\sigma _3+V_0,`$ (3.47)
$`V_0`$ $`=`$ $`{\displaystyle \frac{1}{2}}[\sigma _3,\mathrm{\Gamma }_1],`$ (3.48)
$`\mathrm{\Gamma }_1`$ $`=`$ $`\text{res}_{z=\mathrm{}}(Z(z)z^{n\sigma _3})(\text{see also (}\text{3.12}\text{)}).`$ (3.49)
Equations (3.42) and (3.46) form an overdetermined system for the function $`\mathrm{\Phi }`$ in the variables $`z`$ and $`t`$. From the compatibility condition,
$$\mathrm{\Phi }_{zt}=\mathrm{\Phi }_{tz},$$
we derive the following equation for the coefficient matrices $`B`$ and $`V`$,
$$B_t(z)V_z(z)=[V(z),B(z)],$$
(3.50)
or, taking into account (3.47),
$$B_t(z)\frac{\sigma _3}{2}=\frac{z}{2}[\sigma _3,B(z)]+[V_0,B(z)].$$
(3.51)
Since this equation is satisfied identically in $`z`$, a comparison of the principal parts of both the sides at $`z=0,r_1,\mathrm{},r_M`$, then leads to the differential relations,
$$\frac{B_\alpha }{t}=[\frac{r_\alpha }{2}\sigma _3+V_0,B_\alpha ].\alpha =0,1,\mathrm{},M,$$
(3.52)
(It is notationally convenient to define $`r_0=0`$.)
The important point now is that the matrix $`V_0`$ can be expressed in terms of the matrices $`B_\alpha `$, so that relations (3.52) form a closed system of nonlinear ODEs for the matrix residues $`B_\alpha `$. In fact, expanding both sides of (3.35) in a Laurent series at $`z=\mathrm{}`$, using (3.29), (3.39), and (3.12), and equating the terms of order $`z^1`$ we have
$`{\displaystyle \underset{\alpha =0}{\overset{m}{}}}B_\alpha `$ $`=`$ $`{\displaystyle \frac{t}{2}}[\mathrm{\Gamma }_1,\sigma _3]+(n+k)\left(\begin{array}{cc}0& 0\\ 0& 1\end{array}\right){\displaystyle \underset{\alpha =1}{\overset{m}{}}}k_\alpha \left(\begin{array}{cc}0& 0\\ 0& 1\end{array}\right)+n\sigma _3`$ (3.57)
$`=`$ $`{\displaystyle \frac{t}{2}}[\mathrm{\Gamma }_1,\sigma _3]+n\left(\begin{array}{cc}0& 0\\ 0& 1\end{array}\right)+n\sigma _3`$ (3.60)
$`=`$ $`{\displaystyle \frac{t}{2}}[\mathrm{\Gamma }_1,\sigma _3]+{\displaystyle \frac{n}{2}}\sigma _3+{\displaystyle \frac{n}{2}}I.`$ (3.61)
Comparing the last equation with (3.48) we obtain
$$V_0=\frac{1}{t}\underset{\alpha =0}{\overset{m}{}}B_\alpha \frac{n}{2t}\sigma _3\frac{n}{2t}I,$$
(3.62)
so that (3.52) becomes
$$\frac{B_\alpha }{t}=\frac{ntr_\alpha }{2t}[B_\alpha ,\sigma _3]+\underset{\gamma =0}{\overset{m}{}}\frac{[B_\gamma ,B_\alpha ]}{t},\alpha =0,1,\mathrm{},M.$$
(3.63)
If we vary the points $`r_\alpha `$, then we obtain $`M`$ additional linear differential equations for $`\mathrm{\Phi }(z)=\mathrm{\Phi }(z,t,r_1,\mathrm{},r_M)`$,
$$\mathrm{\Phi }_{r_\alpha }(z)=\frac{B_\alpha }{zr_\alpha }\mathrm{\Phi }(z),\alpha =1,\mathrm{},M.$$
(3.64)
Indeed, introducing the $`r_\alpha `$ \- logarithmic derivative,
$$U_\alpha (z):=\mathrm{\Phi }_{r_\alpha }(z)\mathrm{\Phi }^1(z),$$
and using exactly the same line of arguments as before, we conclude that $`U_\alpha `$ is analytic on $`\text{C}\{0,r_1,\mathrm{},r_m\}`$. Simultaneously, the $`r_\alpha `$ \- vacuum equation (3.28) implies the identity,
$$U_\alpha (z)=\frac{k_\alpha }{zr_\alpha }Z(z)\left(\begin{array}{cc}0& 0\\ 0& 1\end{array}\right)Z^1(z)+Z_{r_\alpha }(z)Z^1(z),$$
(3.65)
(cf. (3.35) and (3.44)) which indicates that the only singularity of $`U_\alpha `$ is a simple pole at $`z=r_\alpha `$ with
$$\frac{k_\alpha }{zr_\alpha }Z(r_\alpha )\left(\begin{array}{cc}0& 0\\ 0& 1\end{array}\right)Z^1(r_\alpha )\frac{B_\alpha }{zr_\alpha }(\text{see also (}\text{3.41}\text{)})$$
as the corresponding principal part. Moreover, taking into account that the asymptotics of $`Z(z)`$ as $`z\mathrm{}`$ does not depend on $`r_\alpha `$ we conclude that
$$U_\alpha (z)0,z\mathrm{},$$
and hence
$$U_\alpha (z)=\frac{B_\alpha }{zr_\alpha }.$$
Equation (3.64) now follows.
The compatibility conditions of equations (3.64) with (3.42) lead to the nonlinear $`r`$-differential equations for the matrices $`B_\alpha =B_\alpha (t,r_1,\mathrm{},r_M)`$,
$`{\displaystyle \frac{B_\alpha }{r_\gamma }}`$ $`=`$ $`{\displaystyle \frac{[B_\alpha ,B_\gamma ]}{r_\alpha r_\gamma }},\alpha \gamma =1,\mathrm{},M,`$ (3.66)
$`{\displaystyle \frac{B_0}{r_\alpha }}`$ $`=`$ $`{\displaystyle \frac{[B_0,B_\alpha ]}{r_0r_\alpha }},\alpha =1,\mathrm{},M,`$ (3.67)
$`{\displaystyle \frac{B_\alpha }{r_\alpha }}`$ $`=`$ $`{\displaystyle \underset{\gamma \alpha }{}}{\displaystyle \frac{[B_\gamma ,B_\alpha ]}{r_\gamma r_\alpha }},\alpha =1,\mathrm{},M.`$ (3.68)
which supplement $`t`$-equation (3.63).
The total system ((3.63), (3.66)–(3.68)) of nonlinear PDEs is the (generalized) system of Schlesinger equations which describes the isomonodromy deformations (see e.g. ) of the coefficients of the $`2\times 2`$ system of linear ODEs having $`M+1`$ regular singularities at the points $`z=r_\alpha ,\alpha =0,\mathrm{},M`$ and an irregualr singular point of Poincaré index 1 at infinity (see (3.42), (3.39)),
$$\frac{d\mathrm{\Phi }(z)}{dz}=B(z)\mathrm{\Phi }(z),B(z)=\frac{t}{2}\sigma _3+\underset{\alpha =0}{\overset{M}{}}\frac{B_\alpha }{zr_\alpha }.$$
(3.69)
The monodromy data of equation (3.69) which single out the solution of ((3.63), (3.66)–(3.68)), which we are interested in, coincide, after the proper normalization, with the data of the $`Z`$-RH problem. More precisely, let us denote $`\mathrm{\Phi }^{\mathrm{}}(z)`$ the analytic continuation of $`\mathrm{\Phi }(z)`$ from $`|z|>1`$ to the whole complex $`z`$-plane. Then, the $`Z`$-RH problem and equation (3.30) imply the following representations of the function $`\mathrm{\Phi }^{\mathrm{}}(z)`$ in the neighborhoods of its singular points,
$`\mathrm{\Phi }^{\mathrm{}}(z)`$ $`=`$ $`\widehat{\mathrm{\Phi }}_\alpha (z)\left(\begin{array}{cc}1& 0\\ 0& (zr_\alpha )^{k_\alpha }\end{array}\right)\left(\begin{array}{cc}1& 1\\ 0& 1\end{array}\right),zU_{r_\alpha },`$ (3.74)
$`\mathrm{\Phi }^{\mathrm{}}(z)`$ $`=`$ $`\widehat{\mathrm{\Phi }}_0(z)\left(\begin{array}{cc}1& 0\\ 0& z^{n+k}\end{array}\right)\left(\begin{array}{cc}1& 1\\ 0& 1\end{array}\right),zU_0,`$ (3.79)
$`\mathrm{\Phi }^{\mathrm{}}(z)`$ $`=`$ $`\widehat{\mathrm{\Phi }}_{\mathrm{}}(z)e^{\frac{tz}{2}\sigma _3}\left(\begin{array}{cc}z^n& 0\\ 0& 1\end{array}\right),zU_{\mathrm{}},\widehat{\mathrm{\Phi }}_{\mathrm{}}(\mathrm{})=I.`$ (3.82)
Here $`\widehat{\mathrm{\Phi }}_\alpha (z)`$, $`\widehat{\mathrm{\Phi }}_0(z)`$, and $`\widehat{\mathrm{\Phi }}_{\mathrm{}}(z)`$ denote the matrix functions which are holomorphic and invertible in the neighborhoods $`U_{r_\alpha }`$, $`U_0`$, and $`U_{\mathrm{}}`$ respectively. Formulae (3.74)–(3.82) allow us to identify the diagonal matrices,
$$E_\alpha =\left(\begin{array}{cc}0& 0\\ 0& k_\alpha \end{array}\right),E_0=\left(\begin{array}{cc}0& 0\\ 0& n+k\end{array}\right),\text{and}E_{\mathrm{}}=\left(\begin{array}{cc}n& 0\\ 0& 0\end{array}\right),$$
(3.83)
as the formal monodromy exponents (cf. ) of $`\mathrm{\Phi }^{\mathrm{}}(z)`$ at the points $`r_\alpha `$, $`0`$, and $`\mathrm{}`$ respectively. The corresponding connection matrices, i.e. the matrices $`C_\alpha `$ in the representations,
$$\mathrm{\Phi }^{\mathrm{}}(z)=\widehat{\mathrm{\Phi }}_\alpha (z)(zr_\alpha )^{E_\alpha }C_\alpha ,\alpha =0,\mathrm{},M,$$
all are given by
$$C_\alpha =\left(\begin{array}{cc}1& 1\\ 0& 1\end{array}\right),\alpha =0,\mathrm{},M.$$
(3.84)
Since the numbers $`k_\alpha `$, $`k`$, and $`n`$ are integers, all the monodromy matrices of $`\mathrm{\Phi }^{\mathrm{}}(z)`$ are trivial. There are also no Stokes’ matrices at the irregular singular point $`z=\mathrm{}`$ since the asymptotic series (3.12), as a Laurent series, converges in a disk centered at infinity. Therefore the complete monodromy data of the linear system (3.69) for our random word problem, consists of (i) (3.83), the formal monodromy exponents at the singular points, and (ii) (3.84), the corresponding connection matrices.
### 3.4 Toeplitz determinant as a $`\tau `$-function
In this section we shall derive the exact formulae for the logarithmic derivatives of the Toeplitz determinant $`D_n(t,r_1,\mathrm{},r_M)`$ in terms of the matrices $`B_\alpha `$ which, as we saw in the previous section, satisfy the Schlesinger equations ((3.63), (3.66)–(3.68)). To this end we will exploit (3.11) and (3.21) whose right hand sides we will express via $`B_\alpha `$ using a technique similar to the one that led to (3.62). We begin with (3.11).
Equation (3.62) was obtained by expanding both sides of (3.35) about $`\mathrm{}`$ and then equating the terms of order $`z^1`$. Let us now analyze the terms of order $`z^2`$. From (3.29) it follows that
$`\mathrm{\Omega }(z)`$ $`=`$ $`{\displaystyle \frac{t}{2}}\sigma _3+{\displaystyle \frac{n}{z}}\left(\begin{array}{cc}0& 0\\ 0& 1\end{array}\right){\displaystyle \frac{1}{z^2}}{\displaystyle \underset{\alpha =1}{\overset{M}{}}}r_\alpha k_\alpha \left(\begin{array}{cc}0& 0\\ 0& 1\end{array}\right)+O\left({\displaystyle \frac{1}{z^3}}\right)`$
$`:=`$ $`{\displaystyle \frac{t}{2}}\sigma _3+{\displaystyle \frac{\mathrm{\Omega }_1}{z}}+{\displaystyle \frac{\mathrm{\Omega }_2}{z^2}}+O\left({\displaystyle \frac{1}{z^3}}\right).`$
Combining this with expansion (3.12) of $`Z`$, we get the following expression for the order $`z^2`$ term of the right hand side of (3.35):
$$\frac{t}{2}[\sigma _3,\mathrm{\Gamma }_1]\mathrm{\Gamma }_1+\frac{t}{2}[\mathrm{\Gamma }_2,\sigma _3]+[\mathrm{\Gamma }_1,\mathrm{\Omega }_1]+\mathrm{\Omega }_2\mathrm{\Gamma }_1+\frac{n}{2}[\mathrm{\Gamma }_1,\sigma _3].$$
The order $`z^2`$ term of the left hand side follows directly from (3.39):
$$\underset{\alpha =1}{\overset{M}{}}r_\alpha B_\alpha .$$
Equating the two expressions we arrive at
$$\underset{\alpha =1}{\overset{M}{}}r_\alpha B_\alpha =\frac{t}{2}[\sigma _3,\mathrm{\Gamma }_1]\mathrm{\Gamma }_1+\frac{t}{2}[\mathrm{\Gamma }_2,\sigma _3]+[\mathrm{\Gamma }_1,\mathrm{\Omega }_1]+\mathrm{\Omega }_2\mathrm{\Gamma }_1+\frac{n}{2}[\mathrm{\Gamma }_1,\sigma _3].$$
(3.86)
This equation together with (3.61) determines $`\mathrm{\Gamma }_1`$ in terms of the matrices $`B_\alpha `$. Indeed, (3.61) gives the off diagonal part of $`\mathrm{\Gamma }_1`$. Using that for $`L`$ diagonal
$$\text{diag}[P,L]=0,$$
we have from (3.86) that
$$\text{diag}\mathrm{\Gamma }_1=\text{diag}\underset{\alpha =1}{\overset{m}{}}r_\alpha B_\alpha +\frac{t}{2}\text{diag}([\sigma _3,\mathrm{\Gamma }_1]\mathrm{\Gamma }_1)+\mathrm{\Omega }_2.$$
(3.87)
Using the identity that for any $`2\times 2`$ matrix $`P`$,
$$\text{diag}([\sigma _3,P]P)=\frac{1}{2}[\sigma _3,P]^2\sigma _3,$$
we obtain from (3.87) and (3.61) the final expression for the diagonal part of $`\mathrm{\Gamma }_1`$,
$`\text{diag}\mathrm{\Gamma }_1`$ $`=`$ $`\text{diag}{\displaystyle \underset{\alpha =1}{\overset{M}{}}}r_\alpha B_\alpha {\displaystyle \frac{1}{t}}\left({\displaystyle \underset{\alpha =0}{\overset{M}{}}}B_\alpha {\displaystyle \frac{n}{2}}\sigma _3{\displaystyle \frac{n}{2}}I\right)\left({\displaystyle \underset{\alpha =0}{\overset{M}{}}}B_\alpha {\displaystyle \frac{n}{2}}\sigma _3{\displaystyle \frac{n}{2}}I\right)\sigma _3+\mathrm{\Omega }_2`$ (3.88)
$`=`$ $`\text{diag}{\displaystyle \underset{\alpha =1}{\overset{M}{}}}r_\alpha B_\alpha {\displaystyle \frac{1}{t}}\text{diag}\left(\left({\displaystyle \underset{\alpha =0}{\overset{M}{}}}B_\alpha {\displaystyle \frac{n}{2}}\sigma _3{\displaystyle \frac{n}{2}}I\right){\displaystyle \underset{\alpha =0}{\overset{M}{}}}B_\alpha \sigma _3\right)+\mathrm{\Omega }_2`$
$`=`$ $`\text{diag}{\displaystyle \underset{\alpha =1}{\overset{M}{}}}r_\alpha B_\alpha {\displaystyle \frac{1}{t}}\text{diag}\left({\displaystyle \underset{\alpha ,\gamma =0}{\overset{M}{}}}B_\alpha B_\gamma \sigma _3{\displaystyle \frac{n}{2}}{\displaystyle \underset{\alpha =0}{\overset{M}{}}}(\sigma _3+I)B_\alpha \sigma _3\right)+\mathrm{\Omega }_2`$
$`=`$ $`\text{diag}{\displaystyle \underset{\alpha =1}{\overset{M}{}}}r_\alpha B_\alpha {\displaystyle \frac{1}{t}}\text{diag}\left({\displaystyle \underset{\alpha ,\gamma =0}{\overset{M}{}}}B_\alpha B_\gamma \sigma _3{\displaystyle \frac{n}{2}}{\displaystyle \underset{\alpha =0}{\overset{M}{}}}B_\alpha (\sigma _3+I)\right)+\mathrm{\Omega }_2.`$
We also made use of the identities,
$$\text{diag}(PL)=0\text{if}\text{diag}P=0\text{and}\text{diag}L=L,$$
and
$$\text{diag}\left(\underset{\alpha =0}{\overset{M}{}}B_\alpha \frac{n}{2}\sigma _3\frac{n}{2}I\right)=0,$$
(3.89)
(The latter follows from (3.61).)
We are at last ready to evaluate $`\mathrm{log}D_n/t`$ in terms of $`B_\alpha `$. To this end it is convenient to use
$$\text{trace}\mathrm{\Gamma }_1=0(\text{which follows from}detZ1)$$
to rewrite (3.11) in the form
$$\frac{}{t}\mathrm{log}D_n(t)=\frac{1}{2}\text{trace}\left(\mathrm{\Gamma }_1\sigma _3\right),$$
(3.90)
and then use (3.88) and (3.89) to obtain
$$\frac{}{t}\mathrm{log}D_n(t)=\frac{1}{2}\underset{\alpha =0}{\overset{M}{}}r_\alpha \text{trace}(B_\alpha \sigma _3)+\frac{1}{2t}\underset{\alpha ,\gamma =0}{\overset{M}{}}\text{trace}B_\alpha B_\gamma \frac{n^2}{2t}\frac{1}{2}\underset{\alpha =1}{\overset{M}{}}r_\alpha k_\alpha .$$
(3.91)
For future comparison with the Jimbo-Miwa-Ueno $`\tau `$-function it is convenient to use (3.89) one more time and rewrite (3.91) as
$$\frac{}{t}\mathrm{log}D_n(t)=\frac{1}{2}\underset{\alpha =0}{\overset{M}{}}r_\alpha \text{trace}(B_\alpha \sigma _3)+\frac{1}{2t}\underset{ji=1,2}{}(\underset{\alpha =0}{\overset{M}{}}B_\alpha )_{ij}(\underset{\gamma =0}{\overset{M}{}}B_\gamma )_{ji}\frac{1}{2}\underset{\alpha =1}{\overset{M}{}}r_\alpha k_\alpha .$$
(3.92)
Let us now perform the similar transformations with the right hand side of equation (3.21). We first notice that its subscripts - free form can be written down as
$$\frac{}{r_\alpha }\mathrm{log}D_n(t)=\text{trace}\left(Z^1(r_\alpha )Z^{}(r_\alpha )E_\alpha \right),$$
(3.93)
where
$$E_\alpha =\left(\begin{array}{cc}0& 0\\ 0& k_\alpha \end{array}\right)$$
is the formal monodromy exponent at $`r_\alpha `$ (see (3.83)) and in transforming (3.21) into (3.93) we took into account that $`detZ(z)1`$. Secondly, by rewriting equation (3.35) as the equation
$$Z^1(z)Z^{}(z)=Z^1(z)B(z)Z(z)\mathrm{\Omega }(z),$$
we get the following representation of the product $`Z^1(r_\alpha )Z^{}(r_\alpha )E_\alpha `$,
$`Z^1(r_\alpha )Z^{}(r_\alpha )E_\alpha `$ $`=`$ $`{\displaystyle \frac{t}{2}}Z^1(r_\alpha )\sigma _3Z(r_\alpha )E_\alpha +{\displaystyle \underset{\genfrac{}{}{0pt}{}{\gamma =0}{\gamma \alpha }}{\overset{M}{}}}{\displaystyle \frac{Z^1(r_\alpha )B_\gamma Z(r_\alpha )E_\alpha }{r_\alpha r_\gamma }}`$ (3.94)
$``$ $`{\displaystyle \frac{t}{2}}\sigma _3E_\alpha {\displaystyle \underset{\genfrac{}{}{0pt}{}{\gamma =0}{\gamma \alpha }}{\overset{M}{}}}{\displaystyle \frac{E_\gamma E_\alpha }{r_\alpha r_\gamma }}+[Z^1(r_\alpha )Z^{}(r_\alpha )E_\alpha ,E_\alpha ].`$
(For notational convenience we set, as before, $`r_0:=0`$ and $`k_0:=nk`$.)
Using equation (3.94) in the right hand side of equation (3.93) and taking into account that
$$B_\alpha =Z(r_\alpha )E_\alpha Z^1(r_\alpha ),$$
we arrive to the following $`r_\alpha `$-analog of (3.91)
$`{\displaystyle \frac{}{r_\alpha }}\mathrm{log}D_n(t)`$ $`=`$ $`{\displaystyle \frac{t}{2}}\text{trace}(B_\alpha \sigma _3)+{\displaystyle \underset{\genfrac{}{}{0pt}{}{\gamma =0}{\gamma \alpha }}{\overset{M}{}}}{\displaystyle \frac{\text{trace}\left(B_\alpha B_\gamma \right)}{r_\alpha r_\gamma }}`$ (3.95)
$`{\displaystyle \frac{k_\alpha t}{2}}{\displaystyle \underset{\genfrac{}{}{0pt}{}{\gamma =0}{\gamma \alpha }}{\overset{M}{}}}{\displaystyle \frac{k_\alpha k_\gamma }{r_\alpha r_\gamma }},`$
Combining equations (3.92) and (3.95) we obtain the main result of this section which is the following equation for the total differential of the function $`\mathrm{log}D_n(t,r_1,\mathrm{},r_M)`$,
$`d\mathrm{log}D_n`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{\genfrac{}{}{0pt}{}{\alpha ,\gamma =0}{\alpha \gamma }}{\overset{M}{}}}\text{trace}\left(B_\alpha B_\gamma \right){\displaystyle \frac{dr_\alpha dr_\gamma }{r_\alpha r_\gamma }}+{\displaystyle \frac{1}{2}}{\displaystyle \underset{\alpha =0}{\overset{M}{}}}\text{trace}(B_\alpha \sigma _3)d(r_\alpha t)`$ (3.96)
$`+{\displaystyle \frac{1}{2}}{\displaystyle \underset{\genfrac{}{}{0pt}{}{1i,j2}{ij}}{}}({\displaystyle \underset{\alpha =0}{\overset{M}{}}}B_\alpha )_{ij}({\displaystyle \underset{\gamma =0}{\overset{M}{}}}B_\gamma )_{ji}{\displaystyle \frac{dt}{t}}`$
$`{\displaystyle \frac{1}{2}}{\displaystyle \underset{\alpha =1}{\overset{M}{}}}k_\alpha d(r_\alpha t){\displaystyle \frac{1}{2}}{\displaystyle \underset{\genfrac{}{}{0pt}{}{\alpha ,\gamma =0}{\alpha \gamma }}{\overset{M}{}}}k_\alpha k_\gamma {\displaystyle \frac{dr_\alpha dr_\gamma }{r_\alpha r_\gamma }}.`$
Equation (3.96) describes the Toeplitz determinant $`D_n(t)`$ in terms of the solution of the Schlesinger system ((3.63), (3.66)–(3.68)) up to a multiplicative constant (depending on $`n`$ and $`k_\alpha `$). Simulteneously, this equation shows, upon comparison with the expression (5.17) in for the logarithmic derivative of the $`\tau `$-function, that
$$D_n(t)=e^{\frac{t}{2}_\alpha r_\alpha k_\alpha }\underset{\genfrac{}{}{0pt}{}{\alpha ,\gamma =0}{\alpha \gamma }}{\overset{M}{}}|r_\alpha r_\gamma |^{\frac{k_\alpha k_\gamma }{2}}\tau _{\text{JMU}},$$
(3.97)
where we use the notation $`\tau _{\text{JMU}}\tau _{\text{JMU}}(t,r_0,r_1,\mathrm{},r_M)`$ for the Jimbo-Miwa-Ueno $`\tau `$-function corresponding to the linear system (3.69) and evaluated for the monodromy data given in (3.83) and (3.84).
Remark. It follows from (3.97) that $`\tau _{\text{JMU}}`$ vanishes as $`r_\alpha r_\gamma `$ for some pair $`(\alpha ,\gamma )`$. This fact, of course, can be established directly from the definition of the $`\tau `$-function.
## 4 Summary of the results
Recall that $`D_n(\varphi )`$ denotes the Toeplitz determinant associated with the symbol
$$\varphi (z)=e^{tz}\underset{\alpha =1}{\overset{M}{}}\left(\frac{zr_\alpha }{z}\right)^{k_\alpha },1<r_\alpha <0,\alpha =1,\mathrm{},M,r_\alpha r_\beta ,\alpha \beta ,$$
$$k_\alpha \text{N},k_\alpha =k,t\text{R},$$
and that the generating function $`G_I(n;\{p_i\},t)`$ is given by the formula
$$G_I(n;\{p_i\},t)=D_n(\varphi ),r_\alpha =p_{i_\alpha }.$$
We also denote by $`D_{n+1}(\varphi |z)`$ the Toeplitz determinant $`D_{n+1}(\varphi )`$ whose last row is replaced by the row $`(1,z,z^2,\mathrm{},z^n)`$, and we shall assume that $`D_n(\varphi )0`$.
The following theorem identifies $`D_n(\varphi )`$ as an object of the theory of integrable systems; more specifically, as an object of the theory of generalized Schlesinger equations developed in .
Theorem 1. Let $`Z`$ denote the $`2\times 2`$ matrix function defined by
$$Z(z)=\left(\begin{array}{cc}\frac{D_{n+1}(\varphi |z)}{D_n(\varphi )}& \frac{i}{2\pi }_C\frac{D_{n+1}(\varphi |z^{})}{D_n(\varphi )}(z^{})^n\varphi (z^{})\frac{dz^{}}{zz^{}}\\ & \\ \frac{\overline{D}_n(\overline{\varphi }|1/\overline{z})}{\overline{D}_n(\overline{\varphi })}z^{n1}& \frac{i}{2\pi }_C\frac{\overline{D}_n(\overline{\varphi }|1/\overline{z^{}})}{\overline{D}_n(\overline{\varphi })}(z^{})^1\varphi (z^{})\frac{dz^{}}{zz^{}}\end{array}\right),$$
(4.1)
where $`C`$ is the unit circle $`|z|=1`$ oriented counterclockwise. Introduce the $`2\times 2`$ matrices $`B_\alpha :=B_\alpha (t):=B_\alpha (\{r_\alpha \},t),\alpha =0,1,\mathrm{},M`$, by the equations,
$$B_\alpha =Z(r_\alpha )\left(\begin{array}{cc}0& 0\\ 0& k_\alpha \end{array}\right)Z^1(r_\alpha ),\alpha =0,\mathrm{},M,$$
(4.2)
where
$$r_0:=0,\text{and}k_0:=nk.$$
(The invertibility of $`Z`$ follows from statement 4 below.) Then the following statements hold:
1. $`d\mathrm{log}D_n(t,r_1,\mathrm{},r_M)`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{\genfrac{}{}{0pt}{}{\alpha ,\gamma =0}{\alpha \gamma }}{\overset{M}{}}}\text{trace}\left(B_\alpha B_\gamma \right){\displaystyle \frac{dr_\alpha dr_\gamma }{r_\alpha r_\gamma }}+{\displaystyle \frac{1}{2}}{\displaystyle \underset{\alpha =0}{\overset{M}{}}}\text{trace}(B_\alpha \sigma _3)d(r_\alpha t)`$ (4.3)
$`+{\displaystyle \frac{1}{2}}{\displaystyle \underset{\genfrac{}{}{0pt}{}{1i,j2}{ij}}{}}({\displaystyle \underset{\alpha =0}{\overset{M}{}}}B_\alpha )_{ij}({\displaystyle \underset{\gamma =0}{\overset{M}{}}}B_\gamma )_{ji}{\displaystyle \frac{dt}{t}}`$
$`{\displaystyle \frac{1}{2}}{\displaystyle \underset{\alpha =1}{\overset{M}{}}}k_\alpha d(r_\alpha t){\displaystyle \frac{1}{2}}{\displaystyle \underset{\genfrac{}{}{0pt}{}{\alpha ,\gamma =0}{\alpha \gamma }}{\overset{M}{}}}k_\alpha k_\gamma {\displaystyle \frac{dr_\alpha dr_\gamma }{r_\alpha r_\gamma }}.`$
2. The matrices $`B_\alpha `$ satisfy the system of nonlinear PDEs (generalized Schlesinger equations),
$`{\displaystyle \frac{B_\alpha }{t}}`$ $`=`$ $`{\displaystyle \frac{ntr_\alpha }{2t}}[B_\alpha ,\sigma _3]+{\displaystyle \underset{\gamma =0}{\overset{M}{}}}{\displaystyle \frac{[B_\gamma ,B_\alpha ]}{t}},\alpha =0,1,\mathrm{},M.`$ (4.4)
$`{\displaystyle \frac{B_\alpha }{r_\gamma }}`$ $`=`$ $`{\displaystyle \frac{[B_\alpha ,B_\gamma ]}{r_\alpha r_\gamma }},\alpha \gamma =1,\mathrm{},M,`$ (4.5)
$`{\displaystyle \frac{B_0}{r_\alpha }}`$ $`=`$ $`{\displaystyle \frac{[B_0,B_\alpha ]}{r_0r_\alpha }},`$ (4.6)
$`{\displaystyle \frac{B_\alpha }{r_\alpha }}`$ $`=`$ $`{\displaystyle \underset{\gamma \alpha }{}}{\displaystyle \frac{[B_\gamma ,B_\alpha ]}{r_\gamma r_\alpha }},\alpha =1,\mathrm{},M.`$ (4.7)
3. Equations (4.4)–(4.7) are the compatibility conditions for the system of linear equations,
$`{\displaystyle \frac{\mathrm{\Phi }(z)}{z}}`$ $`=`$ $`\left({\displaystyle \frac{t}{2}}\sigma _3+{\displaystyle \underset{\alpha =0}{\overset{M}{}}}{\displaystyle \frac{B_\alpha }{zr_\alpha }}\right)\mathrm{\Phi }(z),`$ (4.8)
$`{\displaystyle \frac{\mathrm{\Phi }(z)}{t}}`$ $`=`$ $`\left({\displaystyle \frac{z}{2}}\sigma _3{\displaystyle \frac{n}{2t}}\sigma _3{\displaystyle \frac{n}{2t}}I+{\displaystyle \frac{1}{t}}{\displaystyle \underset{\alpha =0}{\overset{M}{}}}B_\alpha \right)\mathrm{\Phi }(z),`$ (4.9)
$`{\displaystyle \frac{\mathrm{\Phi }(z)}{r_\alpha }}`$ $`=`$ $`{\displaystyle \frac{B_\alpha }{zr_\alpha }}\mathrm{\Phi }(z),\alpha =1,\mathrm{},M,`$ (4.10)
which in turn implies that the system (4.4)–(4.7) describes the isomonodromy deformations of the $`z`$– equation (4.8).
4. The function $`Z`$ is alternatively defined as an unique solution of the matrix Riemann-Hilbert problem,
* $`Z`$ is holomorphic for all $`zC`$,
* $`Z(z)z^{n\sigma _3}I,z\mathrm{}`$,
* $`Z_+(z)=Z_{}(z)\left(\begin{array}{cc}1& z^n\varphi (z)\\ 0& 1\end{array}\right),zC`$.
(In particular, we have that $`detZ1`$.) Equation (4.3) can be rewritten in terms of $`Z`$ as
$$d\mathrm{log}D_n=\left(\text{res}_{z=\mathrm{}}(z^nZ_{22}(z))\right)dt\underset{\alpha =1}{\overset{M}{}}k_\alpha \left(Z_{11}(r_\alpha )Z_{22}^{}(r_\alpha )Z_{21}(r_\alpha )Z_{12}^{}(r_\alpha )\right)dr_\alpha .$$
(4.11)
Also,
$$\frac{D_{n+1}}{D_n}=Z_{12}(0).$$
(4.12)
5. The function
$$\mathrm{\Phi }(z):=Z(z)e^{\frac{tz}{2}\sigma _3}\left(\begin{array}{cc}1& 0\\ 0& z^n\psi ^1(z)\end{array}\right),\psi (z)=\underset{\alpha =1}{\overset{M}{}}\left(\frac{zr_\alpha }{z}\right)^{k_\alpha },$$
satisfies the linear system (4.8)–(4.10) with the matrices $`B_\alpha `$ given by (4.2).
6. The matrices $`B_\alpha `$ are alternatively defined as the solution of the inverse monodromy problem for the linear equation (4.8) characterized by the following monodromy data:
* the formal monodromy exponents at the singular points $`r_\alpha `$, $`\mathrm{}`$ are given by the equations
$$E_\alpha =\left(\begin{array}{cc}0& 0\\ 0& k_\alpha \end{array}\right),E_{\mathrm{}}=\left(\begin{array}{cc}n& 0\\ 0& 0\end{array}\right),$$
* the corresponding connection matrices are
$$C_\alpha =\left(\begin{array}{cc}1& 1\\ 0& 1\end{array}\right),C_{\mathrm{}}=I.$$
* the Stokes matrices at the irregular singular point, $`z=\mathrm{}`$, are trivial.
7. $$D_n(\varphi )=e^{\frac{t}{2}_\alpha r_\alpha k_\alpha }\underset{\genfrac{}{}{0pt}{}{\alpha ,\gamma =0}{\alpha \gamma }}{\overset{M}{}}|r_\alpha r_\gamma |^{\frac{k_\alpha k_\gamma }{2}}\tau _{\text{JMU}},$$
(4.13)
where $`\tau _{\text{JMU}}\tau _{\text{JMU}}(t,r_0,r_1,\mathrm{},r_M)`$ denotes the Jimbo-Miwa-Ueno $`\tau `$-function corresponding to the linear system (4.8) and evaluated for the monodromy data indicated. Equation (4.13) in turn implies the following representation for the generating function $`G_I(n;\{p_i\},t)`$,
$$G_I(n;\{p_i\},t)=e^{\frac{t}{2}}\underset{\genfrac{}{}{0pt}{}{\alpha ,\gamma =0}{\alpha \gamma }}{\overset{M}{}}|p_{i_\alpha }p_{i_\gamma }|^{\frac{k_\alpha k_\gamma }{2}}c\tau _{\text{JMU}}(t,0,p_{i_1},\mathrm{},p_{i_M}),p_{i_0}:=0.$$
(4.14)
Remark 1. In the uniform case, i.e. when $`M=1`$ and $`k_1=k`$, the linear system (4.8) reduces to the $`2\times 2`$ system of linear ODEs which has two regular singular points and one irregular point of Poincaré index 1. In this case, as it is shown in , the isomonodromy equations (4.4)–(4.7) reduce to the special case of the fifth Painlevé equation. Consequently this suggests that the uniform generating function $`G_I(n;t)`$ can be expressed in terms of a solution of the fifth Painlevé equation. That this is so was obtained earlier in via a direct analysis of the Toeplitz determinant $`D_n(t)`$.
Remark 2. The methods developed in this paper can be easily generalized to any symbol $`\varphi (z):=\varphi (z,t)`$ such that
$$_z\mathrm{log}\varphi ,_t\mathrm{log}\varphi $$
are rational in $`z`$.
Remark 3. In virtue of the Fredholm determinant formula (2.2) for the Toeplitz determinant $`D_n(\varphi )`$, equation (4.13) can be interpreted as an example of the general relation between the Jimbo-Miwa-Ueno isomonodromy $`\tau `$-function and the Sato-Segal-Wilson $`\tau `$-function defined via an appropriate determinant bundle (see also for another example of this relation).
Remark 4. The generalized Schlesinger system (4.4)–(4.7) appeared earlier in in connection with the sine kernel Fredholm determinant considered on a union of intervals. The corresponding monodromy data, and hence the solution, are different from the ones related to the Toeplitz determinant $`D_n(t)`$. For instance, the sine kernel monodromy matrices are not trivial (see ; see also for higher matrix dimensional generalizations); in fact, each of them equals the identity matrix plus a one dimensional projection.
Remark 5. From the point of view of the asymptotic analysis of the Toeplitz determinant, the most important statement of Theorem 1 is Statement 4. It allows one to apply the Riemann-Hilbert asymptotic methods of .
Remark 6. This paper has been primarily concerned with the isomondromy/Riemann-Hilbert aspect of our integrable system. Presumably an analysis of the additional compatibility conditions, which arise if one extends (4.8)–(4.10) by a relevant $`n`$-difference equation, would lead to a Toda like system, see Okounkov and Adler and van Moerbeke .
Acknowledgments
This work was begun during the MSRI Semester Random Matrix Models and Their Applications. We wish to thank D. Eisenbud and H. Rossi for their support during this semester. This work was supported in part by the National Science Foundation through grants DMS–9801608, DMS–9802122 and DMS–9732687. The last two authors thank Y. Chen for his kind hospitality at Imperial College where part of this work was done as well as the EPSRC for the award of a Visiting Fellowship, GR/M16580, that made this visit possible. We are also grateful to the referee for several valuable suggestions.
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# Mott transition, antiferromagnetism, and unconventional superconductivity in layered organic superconductors
## Abstract
The phase diagram of the layered organic superconductor $`\kappa `$-(ET)<sub>2</sub>Cu\[N(CN)<sub>2</sub>\]Cl has been accurately measured from a combination of <sup>1</sup>H NMR and AC susceptibility techniques under helium gas pressure. The domains of stability of antiferromagnetic and superconducting long-range orders in the pressure vs temperature plane have been determined. Both phases overlap through a first-order boundary that separates two regions of inhomogeneous phase coexistence. The boundary curve is found to merge with another first order line related to the metal-insulator transition in the paramagnetic region. This transition is found to evolve into a crossover regime above a critical point at higher temperature. The whole phase diagram features a point-like region where metallic, insulating, antiferromagnetic and non s-wave superconducting phases all meet.
The determination of the conditions giving rise to superconductivity (SC) in layered organic conductors constitute one of the chief objectives in understanding the physics of these strongly correlated electronic materials . Closely bound to the now classical issue of proximity of antiferromagnetism (AF) in the emergence of superconductivity stands the problem of the ‘normal’ phase which, depending on pressure conditions in these systems, is either a Mott insulator (MI) or an unconventional metal . A pressure-driven metal-insulator transition can be thus revealing of the strong coupling conditions for electrons that are responsible for broken symmetry states .
In this matter, the phase diagram of the series of layered organic superconductors $`\kappa `$-(BEDT-TTF)<sub>2</sub>X as a function of both hydrostatic and chemical (or anion X substitution) pressures is set to stand out of the debate. By chemical means, the study of anion substituted compounds has allowed few discrete shifts of the pressure scale. Thus for X= Cu\[N(CN)<sub>2</sub>\]Br and X= Cu(NCS)<sub>2</sub>, experiments adduce growing evidence for an unconventional metal and a non s-wave SC state , whereas AF order is shown to become in turn stable on the deuterated X= d<sup>n</sup>-Cu\[N(CN)<sub>2</sub>\]Br compound .
Among all members of the series $`\kappa `$(BEDT-TTF)<sub>2</sub>Cu\[N(CN)<sub>2</sub>\]Cl, denoted as $`\kappa `$ Cl , is the prototype compound of the series showing the complete sequence of states namely, the Mott-insulating, antiferromagnetic, metallic and superconducting states, within a pressure interval of few hundred bars . Despite the numerous experimental efforts recently expended on the properties of this salt, the information collected from experiments done under pressure remained until now scattered and limited by the selectivity of the experimental probe used. Regions of stability of the metallic and superconducting phases have been investigated whereas the information about the pressure profile of AF critical point is missing so far . Our knowledge on the multicritical structure of opposing phases and the nature of the MI transition under pressure is also partial so that a major part of the phase diagram remained until now grounded on a conjectural rather than an empirical basis .
The experiments that are presented in this work were undertaken in order to yield an accurate phase diagram of $`\kappa `$Cl, which is shown in Figure 1. An hydrostatic helium gas pressure technique has been used in order to cover the $`PT`$ phase diagram from both isothermal and isobar sweeps. <sup>1</sup>H NMR and AC suceptibility techniques were simultaneously employed and separated sectors of the phase diagram where either AF or SC state is stable have been unraveled. Both phases meet at $`(P^{},T^{})`$ (282 bar, 13.1K), a point that ends a first order AF-SC boundary, which in turn separates two regions of inhomogeneous coexistence of AF and SC phases. As pressure is swept in the high-temperature paramagnetic domain, one crosses a first-order line associated with the Mott transition and which evolves towards a crossover above a critical point where the MI transition line ends.
All measurements under pressure were performed on single crystals of approximately 0.85 x 0.75 x 0.075 mm<sup>3</sup>, synthesized and grown by standard electrochemical methods . Temperature and pressure sweeps were sufficiently slow to show no dependance on the sweeping rate ($`0.07`$ K/min, $`1`$ bar/min). In order to ensure hydrostatic pressure conditions, our measurements are restricted to the region above the Helium solidification line. <sup>1</sup>H NMR measurements were done in a low static field of 0.5 T oriented perpendicularly to the conducting layers, which essentially allows to obtain a zero field phase diagram. For AC susceptibility, we used a zero-phase feedback loop to track the rf resonance frequency of the NMR circuit with a sensitivity of few ppm. No static field was applied and the AC magnetic field was approximately 0.05 Oe in a direction parallel to the layers $``$ in order to also statisfy NMR requirements.
Different field and temperature conditions, namely 1 T and 3 K in the vortex $``$lock-in $``$ orientation, have been applied at about 300 bar in a clamp pressure cell in an attempt to detect the presence of AF vortex cores and to establish the stability of the coexistence region down to 3 K.
The results for the temperature dependence of the nuclear relaxation rate $`T_1^1`$ are shown in Figure 2 for various pressures. The singular peak of $`T_1^1`$, which marks the onset of critical AF ordering at the Néel temperature $`T_N`$ is gradually suppressed under pressure with the value $`dT_N/dP0.025`$K/bar for the pressure coefficient. Above 275 bar, there is no peak indicating the absence of a magnetic transition. A plateau of $`T_1^1`$, however, persists up to 40 K or so, which marks a sizeable enhancement due to short range AF correlations.
An accurate determination of superconducting order in the phase diagram can be obtained from the AC susceptibility measurements as shown in Figure 3 for selected temperature and pressure sweeps respectively. In the high pressure domain above 400 bar, the $`T_c(P)`$ line below which there is a finite density of superconducting condensate slowly decreases in agreement with previous results ($`dT_c/dP`$ -3.8 K/kbar, squares in Figure 1). As pressure is reduced below 400 bar, lower saturation levels are recorded indicating a gradual suppression of the superconducting order down to 200 bar where it vanishes; $`T_c(P)`$ thus crosses $`T_N(P)`$ and the regions of stability for SC and AF overlap.
When the variation of the superconducting condensate is analyzed under pressure between 8.5K and 12.8 K (Figure 3, bottom), there is a rapid change of the diamagnetic signal centered around a transition pressure denoted $`P_1`$ $``$ defined at the inflexion point. At 8.5 K, $`P_1290`$ bar, and $`P_1`$ slightly moves downward as temperature is raised to finally reach $`P^{}`$ at $`T^{}`$ (double dashed line in Figure 1). According to Figure 3, an hysterisis can be found in a given interval of pressure and temperature. This region of coexistence which conveys some metastability of the transition is delimited by the shaded area in Fig.1.
The nature of this region of the phase diagram can be further sharpened by first looking at the pressure dependence of the <sup>1</sup>H NMR line shape as shown in Figure 4 for 8.5 K. In the AF phase at $`P=15`$ bar, the signal is split into a number of discrete peaks, a characteristic of commensurate AF order ; this structure persists up to 200 bar, above which a narrow peak close to the origin grows in importance concomitant with a reduction of the AF line shape structures. At high pressure, where the system is completely in the SC state, this peak dominates the spectrum. A simulation of the line shape then allows to determine the fraction of the sample that becomes either AF or SC. From Figure 4 (right inset), there is a gradual suppression of AF order which becomes steepest at $`P_1`$290 bar and is tied to a concomitant increase of the SC order. The 10 bar hysteresis at $`P_1`$ is in accordance with the AC susceptibility measurements of Figure 3.
A further, more direct, confirmation of the first order character of the AF-SC transition line is provided by the pressure variation of the AF order parameter as obtained from the NMR spectral line shape. This is exhibited in the left inset of Figure 4 where the AF order parameter at $`T=8.5`$ K shows a slow decrease under pressure followed by an abrupt drop at $`P_1`$=290 bar.
At this point, a few remarks are in order. Although our results do show evidence of coexisting phases, they do not allow to determine if a macroscopic or mesoscopic (e.g stripes, ) type of coexistence is taking place. Given here the constant band filling (half-filling) under pressure, the stripe formation, if it exists, would be of different nature than for high-T<sub>c</sub> materials. The nature of the point $`(P^{},T^{})`$ in the phase diagram is also of interest. Since it exhibits hysteresis (Figs. 3), it can hardly be classified as a bicritical point, which is second order in character. As we will see shortly, however, the point $`(P^{},T^{})`$ also belongs to another transition line associated with the MI transition between two unbroken symmetry states that is, the Mott insulating and the metallic states.
We have used the NMR spectra to check if superconducting vortices with AF cores may be part of the SC-rich sector $``$ AF vortices are predicted in the SO(5) scenario for unification of magnetism and superconductivity . By looking at the modification of internal field distribution in NMR line shape as function of the static field orientation at $`P300`$ bar and down to 3K, however, we failed to detect any decrease of the AF part of the spectra when the static field is oriented along the so-called ‘lock-in’ direction. This orientation corresponds to the situation of intrinsic pinning of vortices between the conducting planes and where vortex cores should sustain a sizeable reduction of their magnetic component. Our results then indicate that the vortex cores are non magnetic, at the very least in this region of the phase diagram.
When AC susceptibility measurements are performed under pressure in the paramagnetic temperature domain, a jump in the diamagnetic signal, albeit small in amplitude, is clearly found (Figure 5). It reveals an increase of diamagnetism when the system enters in the metallic phase $``$ through skin depth effect. These observations corroborate previous electrical transport measurements by Ito et al., who located the MI transition in the same pressure and temperature domain .
When the temperature is increased, the jump in the diamagnetic susceptibility evolves towards a smooth concave profile above some point $`(P_0,T_0)(220\mathrm{bar},32.5\mathrm{K})`$. Well below this point, the diamagnetic susceptibility shows a small but detectable hysteresis that decreases in amplitude as the temperature is raised from ($`P^{},T^{})`$, indicating that the MI transition is first order in character (inset of Figure 5). Within experimental accuracy, the MI line also starts from ($`P^{},T^{})`$ where all other phases meet. The end point $`(P_0,T_0)`$ can then be conjectured to be a critical point.
Refering to Figure 1, over most part of the MI-metal equilibrium curve $`dP/dT`$ is negative. According to the Clausius-Clapeyron relation a negative $`dP/dT`$ would indicate a reduction of the spin entropy on the insulator side below the metallic level. A reasonable explanation for this could come from low dimensional short-range AF correlations, which extend relatively deep in the paramagnetic domain (Figure 2) and would quench a sizeable part of the spin entropy in the vicinity of $`T_N(P)`$. Sufficently far from the AF transition, however, entropies nearly balance so that $`dP/dT`$ is close to zero around 30 K or so.
As for global phase diagram, important conclusions may be drawn about the description of this series of layered organic superconductors. It is noticeable in the first place that the MI transition is discontinuous and clearly evolves toward a mere crossover above a critical point ($`P_0,T_0`$). To our knowledge, it seems to be the first time that a genuine electronic transition that combines all these characteristics is discovered in a quasi-two-dimensional system at half-filling . In the second place, the joining of the MI line with $`T_N(P)`$ at ($`P^{},T^{})`$ is of great interest since it shows within experimental accuracy the absence of boundary between the metallic and a complete AF phases. This confirms previous inferences made about the absence of itinerant antiferromagnetism in $`\kappa `$(BEDT-TTF)<sub>2</sub>X and the relevance of a description of magnetic ordering in terms of interacting spins localized on dimers . The fact that SC and AF phases overlap below $`P^{}`$ indicates that superconductivity can be directly stabilized from the insulating phase. An inescapable outcome of this result is the obvious exclusion of a weak coupling scenario for the emergence of unconventional pairing in layered organic superconductors as function of pressure.
The existence of the point-like region at $`(P^{},T^{})`$ where metal, Mott insulator, antiferromagnet and non s-wave superconductor all meet demonstrates that strong electron correlations and broken symmetry variables are equally important for a unified theory of antiferromagnetism and non s-wave superconductivity in these compounds .
Acknowledgements.$``$ The authors thank K. Behnia, L. G. Caron, V. J. Emery, A. Georges, C. Kallin, S. Kivelson, R. B. Laughlin, D. Sénéchal, A.-M. Tremblay and D. Zanchi for useful discussions. We are grateful to N. Nardone for his crucial technical support and precious advices. P. W. would like thank the NSF for financial support during his stay at UCLA . C.B thanks the NSERC and the Canadian Institute for Advanced Research for financial support.
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# Measurement of the flux of atmospheric muons with the CAPRICE94 apparatus
## I Introduction
Recent atmospheric neutrino observations by Super-Kamiokande and, with lower statistics, Soudan-2 and MACRO collaborations have been used as evidence of neutrino oscillation. The observed rate of neutrino interactions were compared to the rates calculated using the neutrino fluxes derived from atmospheric cascade simulations . It was found that the observed number of events induced by muon neutrinos is too few compared to that from simulation. While several aspects of the measurements, such as the zenith angular dependence of the observed atmospheric neutrinos , strongly point toward the hypothesis of neutrino oscillations, the precise determination of the allowed and excluded regions in the oscillation parameter space rely heavily on the comparison of measurements with calculations .
Recently Gaisser et al. compared different calculations of the atmospheric neutrino flux and concluded that the differences can be attributed to three main effects. First, the different parametrizations used of the particle production and pion yield in hadronic interactions of the primary cosmic rays with air nuclei. Second, the absolute energy spectra of the primary cosmic rays (protons and helium nuclei). Third, the solar modulation and geomagnetic effects. These differences reflect the existing experimental uncertainties on these topics. The available accelerator data on the pion yield in interactions of hadrons with air nuclei are limited and the results used in the simulations do not cover the whole phase space particularly at low values of the Feynman x. Older measurements of the primary cosmic ray spectra differ by as much as 40% from more recent observations which are in agreement at the level of 10%–20%, compatible with the statistical and systematic uncertainties of the measurements. Besides the uncertainties on the absolute values, the primary cosmic ray fluxes, particularly below 10 GeV/n, are affected by the periodic solar activity and the geomagnetic field. Furthermore, the geomagnetic field affects the distribution of the cascade in the atmosphere and these effects are not accounted for in most of the simulations which use a 1-dimensional approach, that is all secondaries are assumed to be collinear with the direction of the primary particles. Gaisser et al. concluded that the calculated neutrino interaction rates are uncertain at the level of $`\pm 30\%`$. They also concluded that because of the cancellation of errors, the ratio of $`\nu _\mu `$ to $`\nu _e`$ has a smaller uncertainty, on the order of $`\pm 5\%`$.
The majority of the Super-Kamiokande events are in the Sub-GeV region where the above uncertainties are the most critical. Measurements of the flux of atmospheric muons provide a powerful tool for cross-checking the cascade simulations. Nearly all the sub-GeV neutrinos events originate from the $`\pi \mu e`$ decay chain. Neutrinos with an energy of 1 GeV are generated, on average, by muons with an energy of about 3 GeV. This interaction typically takes place at altitudes between 12 and 26 km (200 to 20 $`\mathrm{g}/\mathrm{cm}^2`$). Hence, muon measurements must be performed over a broad energy range, from a few hundred MeV to tens of GeV, and over an extended range of atmospheric depths in order to be suitable for use in cross-checking.
We present in this paper results on the muon spectra in the atmosphere obtained with the CAPRICE94 instrument from ground (360 m) to float (36–38 km) altitude. First results on this analysis were reported earlier . In we compared our measured fluxes to results of simulations and concluded that the calculations overestimated our measured muon fluxes; the differences depending on momentum and atmospheric depth. Theoreticians are introducing changes in the simulations procedures, e.g. changing from one-dimensional to three dimensional interaction models, to account for these discrepancies. Therefore, in this paper we do not make any comparisons with simulation results. In this paper we provide details of the data analysis and present our final results. We describe the detector system in this experiment in section 2, the data analysis in section 3 and the results and discussion comprise section 4.
It is worth mentioning that the primary cosmic ray hydrogen and helium spectra also were measured with the CAPRICE94 instrument. These can be used as the input spectra for the cascade simulations in order to reduce the overall systematic uncertainties associated with the comparison of observed and calculated muon fluxes.
## II Detector system
Figure 1 shows the NMSU-WiZard/CAPRICE spectrometer that was operated on ground at Lynn Lake, Manitoba, Canada (56.5 North Latitude, 101.0 West Longitude, 360 m altitude), in summer 1994. The payload was flown by balloon from Lynn Lake to Peace River, Alberta, Canada (56.15 North Latitude, 117.2 West Longitude), on August 8–9, 1994. The balloon reached a altitude of 36.0 km, corresponding to an atmospheric pressure of 4.5 mbar, in about three hours of ascent and it floated at altitudes ranging from 38.1 to 36 km, i.e. residual atmosphere of 3.3–4.6 $`\mathrm{g}/\mathrm{cm}^2`$, for about 23 hours. The apparatus included from top to bottom: a Ring Imaging Cherenkov (RICH) detector, a time-of-flight (ToF) system, a superconducting magnet spectrometer equipped with multiwire proportional chambers (MWPC) and drift chambers (DC) and a silicon-tungsten imaging calorimeter.
The 51.2$`\times `$51.2 cm<sup>2</sup> RICH detector , used a 1 cm thick solid sodium flouride (NaF) radiator with a threshold Lorentz factor of 1.5, and a photosensitive MWPC with pad readout to detect the Cherenkov light image and hence measure the velocity of the particles.
The time-of-flight system had two layers one above and one below the tracking stack, each layer made of two 1 cm thick 25$`\times `$50 cm<sup>2</sup> paddles of plastic scintillator. Each paddle had two 5 cm diameter photomultiplier tubes, one at each end. The distance between the two scintillator layers was 1.1 m. The time-of-flight system was used to give a trigger, to measure the time-of-flight and the ionization (dE/dX) losses of the particles. The trigger was a four-fold coincidence between two photomultiplier tubes with signals in the top paddle and two in the bottom scintillator paddle. The threshold of each photomultiplier tube was set high enough to eliminate noise, but low enough to provide an efficiency of nearly 100% to trigger minimum ionizing particles.
The spectrometer consisted of a superconducting magnet and a tracking device equipped with multiwire proportional chambers and drift chambers . The magnet consisted of a single coil of 11 161 turns of copper-clad NbTi wire. The outer diameter of the coil was 61 cm and the operating current was 120 A, producing an inhomogeneous field of about 4 T at the center of the coil. The spectrometer provided 19 position measurements (12 DC and 7 MWPC) in the direction of maximum bending (x) and 12 measurements (8 DC and 4 MWPC) along the perpendicular direction (y). Using the position information together with the map of the magnetic field, the rigidity of the particle was determined.
Finally, the electromagnetic calorimeter consisted of eight 48$`\times `$48 cm<sup>2</sup> planes of silicon strip (3.6 mm wide) detectors with both x and y readouts. These silicon planes were interleaved with layers of tungsten converters, each one radiation length thick. Taking into account all the material, the calorimeter had a total thickness of 7.2 radiation lengths and 0.33 nuclear interaction lengths. The calorimeter provided topological information on both the longitudinal and lateral profiles of the particle’s interaction as well as a measure of the total energy deposited in the calorimeter.
## III Data analysis
The analysis was based on 85800 seconds of data taken at ground, 10000 seconds during the ascent of the payload and 60520 seconds of data taken at float altitude under an average residual atmosphere of 3.9 g/cm<sup>2</sup>.
### A Particle selection
The CAPRICE94 instrument was well suited to measure the muon spectra and charge ratio in the atmosphere against a background of electrons, protons and heavier particles. The background in the muon sample depended strongly on the atmospheric depth. At float the dominant background for positive particles was protons, which outnumbered the positive muons by about a factor of 1000. The upper end of the energy interval for the $`\mu ^+`$ measurement was determined by the ability of the RICH to reject proton events. At increasing atmospheric depths, the abundance of the proton component decreased to a few percent of the positive muon component at ground level. For negatively charged particles at float, the electron component was the dominant background. The electron background rapidly decreased with increasing atmospheric depth becoming smaller than the muon component by 200 $`\mathrm{g}/\mathrm{cm}^2`$ and a small fraction of the muon component at ground level . Because of this varying background different selection criteria were used for $`\mu ^+`$ and $`\mu ^{}`$ and for ground, ascent and float data to maximize the efficiency while keeping the rejecting power for background events at an appropriate level.
The acceptance of the instrument allowed for muons with a range of zenith angles to be measured. The maximum angle was 20 degrees and the mean of the distribution was at 9 degrees. Figure 2 shows the cosine zenith angle distribution of for muons of both signs selected between 0.15 and 2 GV/$`c`$ at ground level, during the ascent and at float. The distributions have been normalized for the total number of events. No significant change is found in the zenithal angle distribution. It is worth pointing out that the distribution narrows as the rigidity increases.
The selections used for identifying muons in the ground, ascent and flight data are summarized in Table I and described below. The figures describing several of these selections show float data since at float the background of other particles was the largest.
#### 1 Tracking
The tracking information was used to determine the rigidity of the particles. In this work the trajectory was determined by fitting only the information from the drift chambers. This made it possible to use the MWPC system for the efficiency estimation of the drift chambers (see section III C 1). To achieve a reliable estimation of the rigidity, a set of conditions were imposed on the fitted tracks:
1. At least 9 (out of 12) position measurements in the x direction and 5 (out of 8) in the y direction were used in the fit.
2. There should be an acceptable chi-square for the fitted track.
3. The estimated error on the deflection should be $`<0.04`$ (GV/$`c`$)<sup>-1</sup>.
These conditions also eliminated events with more than one track in the spectrometer.
#### 2 Scintillators and time-of-flight
The ionization (dE/dX) loss in the top ToF scintillator was used to select minimum ionizing singly charged particles by requiring a measured dE/dX of less than 1.8 times the most probable energy loss by a minimum ionizing particle (mip).
Downward moving particles were selected using the time-of-flight information. The time-of-flight resolution of 280 ps, which was small compared to the flight-time of more than 4 ns, assured that no contamination from albedo particles remained in the selected sample. The particle’s velocity ($`\beta `$) reconstructed from the time-of-flight information was used to select muons against a background of pions, protons and heavier particles. Figure 3 shows $`\beta `$ obtained from the time-of-flight information for positive particles collected at float as a function of rigidity. Muons were selected as particles above the solid line. The ToF $`\beta `$ selection was used for $`\mu ^+`$ identification for the ascent and float portion of the data and for the $`\mu ^{}`$ selection during float.
#### 3 The calorimeter.
With a depth of 7 radiation lengths the calorimeter could identify non-interacting particles in a background of electrons above a momentum of a few hundred MeV/$`c`$. The selection was performed by requiring that the topological information of the signal was consistent with that of a single track. This was accomplished by imposing an upper limit to the number of hits along the track in the calorimeter. This selection reached its highest efficiency above 1 GeV/$`c`$ where electromagnetic showers were well defined in the calorimeter. Below 1 GeV/$`c`$ a non-negligible electron contamination was present. To further reduce this background, another selection criterion, based on the total detected energy in the calorimeter divided by the momentum, was used. An upper limit for this quantity equal to 60 mip/(GeV/$`c`$) was applied. Figure 4 shows this quantity as a function of rigidity for the float data. The two dense band are due to non-interacting particles. Recall that muons above a few hundred MeV/$`c`$ are minimum ionizing particles that release a nearly constant amount of energy in the calorimeter. Below 300 MeV/$`c`$ the calorimeter selection was not used because of the low efficiency. At ground level e<sup>±</sup> amount to less than 0.1% of the muon component above 3 GV/$`c`$ and, consequently, the calorimeter muon selection criteria were not used above this rigidity.
The calorimeter selection criteria also rejected particles that interacted in the calorimeter, namely pions, protons and heavier nuclei, hence contributing to reduce their contamination in the muon sample. However, the rejection factor for these particles was small because the calorimeter depth was only one third of a nuclear interaction length. Therefore, hadrons were removed using selection criteria based on time-of-flight, scintillator and RICH information.
#### 4 The RICH
The RICH was used to measure the Cherenkov angle of the particle and thereby its velocity. The Cherenkov angle was reconstructed from the geometrical distribution of the signals in the pad plane (for a description of the reconstruction method see ). To correctly use the RICH information, a set of conditions was applied on the RICH data . They were:
1. Ionization from charged particles produced significantly higher signals than converting Cherenkov photons. To reject events with multiple charged particles traversing the RICH, we required that an event contained only one cluster of pads with high signals.
2. A good agreement between the particle impact position as determined by the RICH and the tracking system was required. The difference in x and y should be less than three standard deviations, which was rigidity dependent but typically less than 5 mm.
3. More than 3.5 pads with signals from converted Cherenkov photons were required in the fit for the ground data and more that 7.5 for the ascent and float data.
4. The reconstructed Cherenkov angle should not deviate by more than three standard deviations from the expected Cherenkov angle for muons.
Criterion 1 was introduced to eliminate events with more than one charged particle crossing the RICH and this condition was applied over the entire data sets both for $`\mu ^+`$ and $`\mu ^{}`$ . Criterion 2 eliminated events that scattered in the RICH electronics. Criteria 3 and 4 were used to separate muons from the other particles. Figure 5 shows the measured Cherenkov angle for flight data events selected with RICH criteria 1, 2 and 3. The bands corresponding to the different particles are clearly visible. The solid lines indicate the muon selection based on the Cherenkov angle. However, it is important to point out that the solid lines in the figure are only indicative of the selection, since the RICH selection was done on an event by event basis. For each event the Cherenkov angle was obtained and the resolution of this depended on the incident angle of the particle .
These four criteria were used for selecting muons of both signs at float and for selecting $`\mu ^+`$ at ground and during the ascent. Ground and ascent $`\mu ^{}`$ were selected differently: the selection was based upon the dE/dX and calorimeter criteria except below 70 $`\mathrm{g}/\mathrm{cm}^2`$ of residual atmosphere where contamination of interaction products from the payload structure (see section III B 4) was non-negligible. For these small depths the RICH criterion 1 also was used. In addition the full RICH selection with criteria 1 to 4 was used below 500 MeV/$`c`$ since, at low velocities, the RICH was able to separate muons from pions and electrons.
#### 5 The bar
A 17 kg, 1.2 m long aluminum bar with a 7 kg steel hook in the centre was used to connect the payload to the balloon during the flight. This bar was situated 2.3 m above the RICH. Hadrons had a non-negligible probability to interact with the material of the bar and produce secondary particles that could be detected as muons in the apparatus. Hence we chose to reject all particles with extrapolated trajectories that crossed the bar. This procedure resulted in a reduction of the geometrical factor by about 10% as can be seen in Figure 6.
### B Background rejection
The various sources of background in the muon analysis are described below along with the rejections criteria and surviving contamination. This information is summarized in Table II.
#### 1 Electron background
The calorimeter muon selection gave an electron rejection factor that increased from 30 at 0.3 GeV/$`c`$ to more than 1000 above 1 GeV/$`c`$. The RICH separated muons from electrons with an electron rejection factor of more than 100 at 0.1 GeV/$`c`$ decreasing to 10 at 0.3 GeV/$`c`$ and to 1 at 0.7 GeV/$`c`$. Since the electron flux is higher than the muon one only at high altitudes where the electron primary component dominates and is about a factor four larger than the muon flux , the electron and positron contaminations surviving the muon selection were assumed negligible.
#### 2 Proton background.
The ToF $`\beta `$ rejection factor for protons below 1 GeV/$`c`$ was greater than 4000 decreasing to about 300 at 1.2 GeV/$`c`$. The RICH proton rejection factor was greater than 2000 below 1.5 GeV/$`c`$, about 1000 at 2 GeV/$`c`$ and decreased to about 1 at 5 GeV/$`c`$. Below 1.5 GeV/$`c`$, the combined effect of the time-of-flight and RICH selection criteria reduced the proton contamination to a negligible fraction of the selected muon sample. At higher momenta, because of the strong variation of the proton flux, the contamination was dependent on the atmospheric depth. At float altitude the proton contamination became increasingly important above 1.5 GeV/$`c`$ and it dominated the muon sample above 2.5 GeV/$`c`$. Hence, we conservatively limited the positive muon measurements to a momentum range between 0.15 and 2 GeV/$`c`$. In this range the proton contamination was negligible at all atmospheric depths except at float. In the float data we assumed, as a worst case, that all the singly charged particles were protons. Applying the rejecting power of the calorimeter, dE/dX, ToF beta and RICH to this sample we found that a small (less than 20%) proton contamination survived the muon selection in the rigidity range 1.5 to 2 GeV/$`c`$. This contamination was subtracted from the positive muon sample at float. For ground data the range for $`\mu ^+`$ was extended to higher momenta as presented in . The proton contamination was negligible at ground level below 3 GeV/$`c`$, while above 3 GeV/$`c`$ it was calculated and subtracted from the positive muon sample. This calculation was made by rescaling the number of interacting particles in the calorimeter with factors obtained from data at float. The proton candidates were selected if they had a hadronic interaction in the calorimeter. The contamination from muons in the interacting proton sample was studied using negatively charged ground data, while the contamination from electromagnetic showers was negligible at momenta greater than 3 GeV/$`c`$ . The efficiency of this proton selection was estimated using a sample of singly charged particles at float that was assumed to be composed of protons.
In summary, during the ascent positive muons were selected free of proton contamination up to 2 GeV/$`c`$. At float a small proton contamination was subtracted from the positive muon sample between 1.5 and 2 GeV/$`c`$ while at ground the range was extended to much higher momenta subtracting the small estimated proton contamination .
#### 3 Heavier nuclei background.
The case of deuteron background was very similar to the proton one. Helium and heavier nuclei were mainly rejected by the dE/dX selection. The remaining fraction was eliminated using the other selection methods and their contamination was determined to be negligible.
#### 4 Meson background.
Because mesons (pions and kaons) produced in the atmosphere above the payload decay rapidly, they represented a small ($`<2\%`$) fraction of events compared to the muon flux (e.g. see ). However, pions produced in the payload could still present a non-trivial contamination. We have undertaken a careful analysis of the local pion background in order to quantify their abundance.
Evidence for a pion contamination is visible at very low rigidities in Figure 3. These pions, because of their low energies, were presumably produced locally in the RICH or in the dome (part of the gondola above the RICH). However, it is important to stress that no RICH conditions were applied to select the events in Figure 3: the RICH selection rejected multiply charged tracks, particles interacting in the RICH as well as pions below 500 MV/$`c`$. Also the dE/dX selection rejected multiply charged tracks and this leaves the dome as the source of pions produced in the interaction with high energy cosmic rays. The interactions occurred at such an angle that the high energy secondaries missed the active volume of the instrument, but a low energy pion was produced at a large angle and passed through the instrument appearing like a singly charged particle rather than part of a shower. In studying the pions below $`<0.2`$ GV/$`c`$, where they can be identified using the time-of-flight, we found no evidence of a preferred incoming direction. This is consistent with our interpretation of these events since the distribution of material in the gondola was symmetric with respect to the azimuthal angle. The locally produced pion flux entering the apparatus should decrease quickly with energy because of the emission towards the forward direction and the lower probability of having only one charged particle with high momentum traversing the detectors.
To test the correctness of this conclusion the following approach was adopted. Events were selected from float data with:
1. multiple pad ionization clusters in the RICH;
2. multiply charged signal in the top scintillator;
3. rigidity (R) interval: $`0.5<R<2`$ GV/$`c`$;
4. Cherenkov angle of a $`\beta 1`$ particle;
5. time-of-flight of a $`\beta 1`$ particle;
6. dE/dX signal of a minimum ionizing particle in the bottom scintillator.
With these criteria $`\beta 1`$ particles belonging to a shower initiated closely to the top of the apparatus were selected, hence with high probability secondary pions, which entered the calorimeter as singly charged. The criteria 3 to 6 ensured that the selected events were indeed $`\beta 1`$ particles. This sample of locally produced pions was used to estimate their interaction probability in the calorimeter (or, more precisely, the probability of selecting interacting pions in the calorimeter). Then, the same hadronic interaction selection for the calorimeter was applied to the muon sample selected without using the calorimeter muon criteria. Two samples were obtained, one for positive and one for negative pions. These events (about 100 in total) were visually scanned with a graphic program. In this way misidentified muons and, especially, electrons were rejected from the samples. Then, the remaining pion numbers were rescaled by the calorimeter selection efficiency thus obtaining the number of pions in the muon samples. The result was that at float altitude pions could account for a maximum of 20% between 0.5 and 1 GeV/$`c`$ and for less than 10% above 1 GeV/$`c`$ of the muon flux, irrespective of the sign of the charge. We give this result as an upper limit because the procedure is likely to overestimate the number of pions (some of the selected pions could be muons, etc.). For this reason, it was not subtracted from the muon flux and it should be considered as a systematic uncertainty of the measurement. It is important to stress that this uncertainty, due to the similar pion contamination in the $`\mu ^+`$ and $`\mu ^{}`$ samples, affects the $`\mu ^+`$ to $`\mu ^{}`$ ratio less than the corresponding fluxes.
#### 5 Conclusion on background
Clean $`\mu ^+`$ and $`\mu ^{}`$ samples were selected from 0.15 to 0.4 GeV/$`c`$. Above this momentum a non-negligible contamination of locally produced pions could be present. For the float data, this contamination was less than 20% of the muon flux between 0.5 and 1 GeV/$`c`$ decreasing to less than 10% above 1 GeV/$`c`$. For larger atmospheric depths the locally produced pion flux decreased quickly due to the decrease of the interacting proton and helium nuclei, specially at large zenith angle, while the muon flux increased with increasing depth at least up to 100 $`\mathrm{g}/\mathrm{cm}^2`$. Hence the locally produced pion contamination was assumed negligible at all depths except at float.
At float a small proton contamination was subtracted from the positive muon sample between 1.5 and 2 GeV/$`c`$. At ground the momentum range was extended to much higher value with a subtraction of a small proton contamination.
### C Efficiency determination
In order to accurately determine the fluxes of the various types of particles, the efficiency of each detector was carefully studied using both ground and flight data. To determine the efficiency of a given detector, a data set of muons was selected by the remaining detectors. The number of muons correctly identified by the detector under test divided by the number of events in the data set provided a measure of the efficiency. This procedure was repeated for each detector. The efficiency of each detector was determined as a function of rigidity in a number of discrete bins and, then, parameterized to allow an interpolation between bins. This parameterization introduced a systematic error on the efficiency of each detector. Since the parameters were correlated, the error on the efficiency was obtained using the error matrix of the fit for each detector when correcting the measured flux for the detector efficiencies.
#### 1 Tracking efficiency
The drift chamber tracking efficiency was obtained using negative singly charged particles selected by the other detectors similarly as done in . A sample of singly charged particles was selected by requiring a single ionization cluster in the RICH and a dE/dX signal in the top scintillator typical of a minimum ionizing particle. From this sample, negatively charged events were selected by requiring a negative deflection from the fit to the MWPC trajectory measurements. The contamination of spillover protons was eliminated by requiring that the measured impact positions in the RICH and calorimeter agreed with the positions as obtained by extrapolating the particle trajectory derived from the MWPC fit. The resulting sample of negative singly charged particles was used to determine the efficiency of fitting tracks in the drift chamber system. The solid line in Figure 7 shows the tracking efficiency at float altitude. The efficiency varied from ground to float altitude. At ground it was $`95\%`$ above 1 GeV/$`c`$, then just after the launch it was $`87\%`$ increasing to $`93\%`$ at float altitude.
Biases in the efficiency sample were studied using protons (see ). It was found that the criteria used for fitting the tracks using the MWPC slightly reduced the number of scattered tracks in the sample. In order to account for this reduction a systematic uncertainty of 2% was introduced .
Possible charge sign dependence of the efficiency was studied using both the flight data and the data taken on the ground before the flight. No significant dependence was found above 0.3 GV/$`c`$ .
#### 2 Scintillator efficiency
The dE/dX and $`\beta `$ selections were studied using negative events with a minimum ionizing pattern in the calorimeter. A clean RICH signal also was required to reject interaction products from the sample. The dotted lines in Figure 7 show the scintillator efficiency in its two selections. The efficiency was studied using ground, ascent and float data and no variation was found.
#### 3 Calorimeter efficiency
The calorimeter selection efficiency, shown as a dashed line in Figure 7, was obtained using ground data. The result was cross checked with a simulation of the calorimeter and an excellent agreement was found. The calorimeter efficiency also was studied with flight data. The calorimeter only is able to separate muons from electrons above 0.5 GeV/$`c`$, hence the presence of a larger electron contamination had to be taken into account. Inside the errors a good agreement was found as expected since the calorimeter performances were stable over a period of months.
#### 4 RICH efficiency
On the ground and in the first part of the ascent the RICH efficiency was obtained by selecting negative singly charged particles, which did not interact in the calorimeter. At float altitude a large background of interaction products did not permit us to select an unbiased clean sample of muons, hence the efficiency obtained from the ground data was used. This procedure was validated by comparing the RICH efficiency for selecting electrons at ground and at float since an unbiased clean sample of electrons could be selected using the calorimeter. The RICH electron selection criteria is similar to that for muons, namely, same requirements on the impact position, number of effective pads and Cherenkov angle but using the theoretical electron Cherenkov angle. It was found that the RICH electron efficiency for float data reproduced the electron efficiency of ground electrons inside an uncertainty of about 5%. Therefore, it was reasonable to make use of the muon RICH efficiency as obtained from the ground data (dashed-dotted line in Figure 7) for the flight data.
### D Geometrical factor
The geometrical factor, determined with simulation techniques , is shown in Figure 6 for ground and flight data. The difference at low deflections (high rigidities) for the two sets of data is due to the additional geometrical constrain imposed due to the bar.
The geometrical factor was cross checked with two other methods. One adopted the same approach as presented in using, however, a different method to trace the particles: the track fitting algorithm used in the analysis also was used to trace the particle through the spectrometer. This method gave the same results within 1%, at all rigidities. The second used a numerical integration calculation of the geometrical factor that agreed with the previous results within 2% above 0.5 GV/$`c`$ and within 5% below 0.5 GV/$`c`$.
### E Systematic uncertainties
Systematic uncertainties originating from the determination of the detector efficiencies were included in the tables and data points as discussed in session III C. Another possible systematic error was related to the efficiency of the trigger system. The fraction of each trigger combination was compared with the simulated fraction taking into account the position of each paddle and the magnetic field. The excellent agreement between the simulated and experimental fractions permitted us to conclude that a possible systematic error due to a geometrical inefficiency of the trigger was less than 1%.
The residual atmosphere above the gondola was measured by two pressure sensors owned and calibrated by the National Scientific Balloon Facility. The two measurements did not coincide: their difference increased with altitude from less than 1% to about 10% at float. We interpreted this difference as the systematic uncertainty on the atmospheric depth. This uncertainty does not affect the measurement but has to be taken into account when comparing the measured spectra with the simulated ones.
From the discussion in section III D we conclude that the systematic error due to the geometrical factor calculation was less than 5% between 0.3 and 0.5 GV/$`c`$ and less than 2% for rigidities higher than 0.5 GV/$`c`$. For the geometrical factor calculations it was assumed that there was no variation of the muon intensity over the acceptance angle. The effect on the geometrical factor due to the intensity variation was examined using the measured muon spectra and the observed zenithal distribution in our apparatus in the rigidity range 0.2 to 1.5 GV/$`c`$ at ground. We found that the calculated geometrical factor would be reduced by about 3%.
We decided to assign a systematic uncertainty of 5% to the RICH muon efficiency at float to account for possible variations in the RICH performance between ground and float.
The tracking muon selection efficiency varied with time during the ascent. We determined the efficiency for seven time bins and we found that they could be grouped in three in which the efficiency could be assumed constant. Since the efficiency above 1 GV/$`c`$ varied from about 87% at launch to 93% at float we believe that the systematic uncertainty of this procedure is less than 6%.
Assuming that the systematic uncertainties discussed above are uncorrelated, we estimated an overall systematic uncertainty, which is momentum dependent, for ground muons decreasing from $``$ 6% at 0.3 GeV/$`c`$ to about 2% above 1 GeV/$`c`$. It is worth pointing out that the ground muon fluxes measured by the CAPRICE94 apparatus agree at the level of 3% with the measurements from the CAPRICE97 experiment , which was equipped with the same superconducting magnet and calorimeter but with a different tracking system and with a gas RICH. For ascent and float muons this systematic uncertainty decreased from $``$ 9% at 0.3 GeV/$`c`$ to about 7% above 1 GeV/$`c`$. These uncertainties were not included in the data presented in the tables and the figures.
## IV Results
We selected 37864 $`\mu ^{}`$ and 47043 $`\mu ^+`$ between 0.2 and 120 GeV/$`c`$ at ground (1000 $`\mathrm{g}/\mathrm{cm}^2`$); 5081 $`\mu ^{}`$ between 0.3 and 40 GeV/$`c`$ and 2715 $`\mu ^+`$ between 0.3 and 2 GeV/$`c`$ during the ascent (7–850 $`\mathrm{g}/\mathrm{cm}^2`$); 1601 $`\mu ^{}`$ between 0.18 and 20 GeV/$`c`$ and 2063 $`\mu ^+`$ between 0.18 and 2 GeV/$`c`$ at float altitude (3.3–4.6 $`\mathrm{g}/\mathrm{cm}^2`$, mean atmospheric depth of 3.9 $`\mathrm{g}/\mathrm{cm}^2`$). From these we obtained the muon fluxes ($`J_{\mu ^{},\mu ^+}`$) according to:
$`J_{\mu ^{},\mu ^+}(P,X)`$ $`=`$ $`{\displaystyle \frac{1}{T_{live}(X)\times G_{\mu ^{},\mu ^+}\times ϵ(P,X)\times \mathrm{\Delta }P}}`$ (2)
$`\times N_{\mu ^{},\mu ^+}(P,X),`$
where $`X`$ is the atmospheric depth, $`T_{live}`$ is the live time, $`G_{\mu ^{},\mu ^+}`$ are the geometrical factors for $`\mu ^{}`$ and $`\mu ^+`$ , $`ϵ`$ is the combined selection efficiency, $`\mathrm{\Delta }P`$ is the width of the momentum bin corrected for ionization losses to the top of the payload, $`P`$ the momentum and $`N_{\mu ^{},\mu ^+}(P)`$ is the selected number of $`\mu ^{}`$ and $`\mu ^+`$ . The fractional live time decreased from $`0.97240\pm 0.00001`$ at ground to $`0.2690\pm 0.0006`$ at float altitude as indicated in Table III.
Figure 8 and Table IV show the muon spectra at float, corresponding to 3.9 $`\mathrm{g}/\mathrm{cm}^2`$ of residual atmosphere. These muon fluxes are interesting since these muons are the products of the first interaction between the primary cosmic rays and the air nuclei. Hence, along with the simultaneous measurement of the primary spectra of proton and helium nuclei these data provide a useful testbench for studying the pion production in nucleon-air interaction used in the calculation of atmospheric showers. As discussed in section III B 2 we limit the positive muon data to momenta below 2 GeV/$`c`$. The average muon charge ratio on this momentum interval is $`1.59\pm 0.06`$.
In Table III we present the measured muon fluxes at several atmospheric depths and momenta interval. The symbol FAD stands for Flux-weighted Average Depth obtained according to:
$$\mathrm{FAD}(P)=\frac{X(t)ϵ_{live}(t)J(P)𝑑t}{ϵ_{live}(t)J(P)𝑑t},$$
(3)
where $`ϵ_{live}`$ is the fractional live time. The depth and momentum intervals were chosen to match the published data of $`\mu ^{}`$ flux growth curves by the MASS89 and MASS91 experiment .
Figure 9 shows the flux growth curves for (a) negative and (b) positive muons for several momentum bins. For each momentum interval we fitted the data at large atmospheric depths ($`X>190`$ g/cm<sup>2</sup>) with an exponential function :
$$J(P,X)=k(P)e^{X/\mathrm{\Lambda }(P)},$$
(4)
where $`k`$ and $`\mathrm{\Lambda }`$ are obtained from the fits. The resulting best fits are shown in Figure 9 as solid lines. As found in the MASS89 and MASS91 experiments a nearly linear relation exists between the attenuation length ($`\mathrm{\Lambda }`$) and the mean momentum ($`P`$) in unit of GeV/$`c`$ in the 190 to 1000 $`\mathrm{g}/\mathrm{cm}^2`$ range. The relation resulting from CAPRICE94 measurements of $`\mu ^{}`$ is:
$$\mathrm{\Lambda }[\mathrm{g}/\mathrm{cm}^2]=(263\pm 14)+(150\pm 15)\times P,$$
(5)
This relation holds also for the CAPRICE94 $`\mu ^+`$ flux growth curves. It is worth pointing out that in determining relation 5 we used also the $`\mu ^{}`$ fluxes at ground. Equation 5 can be compared with the one determined from MASS89 $`\mu ^{}`$ data as:
$$\mathrm{\Lambda }[\mathrm{g}/\mathrm{cm}^2]=(283\pm 24)+(93\pm 16)\times P.$$
(6)
Both the above expressions reproduce the data within errors over the momentum range of these experiments, but will differ when extended to much larger momenta.
Figure 10 shows the relative difference between the $`\mu ^{}`$ fluxes obtained in this analysis and the MASS89 and MASS91 experiments as a function of atmospheric depth. The comparison is done for muon momenta below 1 GeV/$`c`$ (a) and between 1 and 2 GeV/$`c`$ (b). The dashed lines indicate the average difference between this analysis and MASS89 and the solid lines the average difference with MASS91. Considering the errors in the data points, a good agreement is found in the 1 to 2 GeV/$`c`$ interval among different measurements. However, below 1 GeV/$`c`$ the CAPRICE94 results are significantly higher than the results from the MASS91 and, to a lesser extent, from the MASS89 experiments. These differences could be caused by solar activity or geomagnetic effects. In fact, the MASS89 experiment was launched from Prince Albert, Saschatcewan, Canada, during a period of maximum solar activity while MASS91 was launched from Fort Sumner, New Mexico, and flew at an average geomagnetic cutoff of about 4.5 GV/$`c`$.
Geomagnetic effects also are observed in the muon charge ratio. Figure 11 shows the $`\mu ^+`$ to $`\mu ^{}`$ ratio as a function of atmospheric depth measured in the momentum intervals 0.3–1 GeV/$`c`$ (a) and 1–2 GeV/$`c`$ (b) by this experiment, by the recent CAPRICE98 experiment , which flew from Fort Sumner on 28-29 May 1998, in the range 0.3–0.9 GeV/$`c`$ by the HEAT95 experiment , which flew from Lynn Lake on the 23rd of August 1995, and in the range 0.3–0.9 GeV/$`c`$ (a) and 0.9–1.5 GeV/$`c`$ (b) by the MASS91 experiment . It can be seen that the CAPRICE94 low momenta charge ratios are higher than the MASS91 and CAPRICE98 ones. Moreover, the CAPRICE94 data show a dependence on the atmospheric depth, which is also visible in the HEAT data. Similar latitude effects also can be seen in Figure 12, which shows (a) the CAPRICE94 data at float along with the charge ratios measured by the CAPRICE98 experiment and the MASS91 experiment and (b) the ground muon data reported here and from the CAPRICE97 experiment , which was carried out in Fort Sumner during Spring 1997.
Figure 13 shows the measured spectra of negative muons for nine depth intervals. Above 1.5 GeV/$`c`$ the $`\mu ^{}`$ spectra between 3.9 and 250 $`\mathrm{g}/\mathrm{cm}^2`$(at larger atmospheric depths unacceptable power law fits were found for this momentum range) are power law in momentum with a fairly constant spectral index of $`2.30\pm 0.04`$ that can be compared with $`2.5\pm 0.2`$ between 20 and 400 $`\mathrm{g}/\mathrm{cm}^2`$ above 2 GeV/$`c`$ from MASS89 and with $`2.45\pm 0.05`$ between 25 and 250 $`\mathrm{g}/\mathrm{cm}^2`$ above 1.5 GeV/$`c`$ from MASS91 .
## V Conclusions
In this paper we have presented new results on atmospheric data measured with the CAPRICE94 experiment both for positive and negative muons. The data cover a large atmospheric depth range from close to the top of the atmosphere (3.9 $`\mathrm{g}/\mathrm{cm}^2`$) down to ground level (1000 $`\mathrm{g}/\mathrm{cm}^2`$).
The data were compared with other experimental results that were obtained using the same superconducting magnet but with different identifying detectors. The muon spectra measured by the different experiments at high momenta (above 1 GeV/$`c`$) are in good agreement considering the overall uncertainty of the measurements ($`1015\%`$). At lower energy, the comparison between the results of CAPRICE94 and those of MASS89/91 (the CAPRICE94 $`\mu ^{}`$ fluxes are about 10–20% higher than the ones measured in the other two experiments) indicates solar modulation and geomagnetic effects. It is worth pointing out that the differences between the different measurements cannot account for the discrepancies found at low momenta, while comparing the experimental data with the theoretical calculation, which are in some cases as large as 70% (see ).
###### Acknowledgements.
This work was supported by NASA Grant NAGW-110, The Istituto Nazionale di Fisica Nucleare, Italy, the Agenzia Spaziale Italiana, DARA and DFG in Germany, EU SCIENCE, the Swedish National Space Board and the Swedish Council for Planning and Coordination of Research. The Swedish-French group thanks the EC SCIENCE programme for support. We wish to thank the National Scientific Balloon Facility and the NSBF launch crew that served in Lynn Lake. We would also like to acknowledge the essential support given by the CERN TA-1 group and the technical staff of NMSU and of INFN.
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# Doped Stripes in Models for the Cuprates Emerging from the One-hole Properties of the Insulator
\[
## Abstract
The extended and standard t-J models are computationally studied on ladders and planes, with emphasis on the small J/t region. At couplings compatible with photoemission results for undoped cuprates, half-doped stripes separating $`\pi `$-shifted antiferromagnetic (AF) domains are found, as in Tranquada’s interpretation of neutron experiments. Our main result is that the elementary stripe “building-block” resembles the properties of $`one`$ hole at small J/t, with robust AF correlations across-the-hole induced by the local tendency of the charge to separate from the spin (G. Martins et al., Phys. Rev. B60, R3716 (1999)). This suggests that the seed of half-doped stripes already exists in the unusual properties of the insulating parent compound.
PACS numbers: 74.20.-z, 74.20.Mn, 75.25.Dw
\]
The understanding of high temperature superconductors is among the most important open problems in strongly correlated electrons. A remarkable development in recent years is the accumulation of experimental evidence compatible with stripe formation in the normal state of underdoped cuprates. This includes spin incommensurability (IC) in neutron experiments, results believed to be caused by stripes separating $`\pi `$-shifted AF domains. More recently, it has been shown that the stripes are metallic, result compatible with proposals of the normal state of x=1/8 cuprates as made out of half-doped stripes. Whether stripe formation is beneficial or detrimental to superconductivity is unclear, but it appears that stripes are an important ingredient of the normal state that cannot be ignored.
The theoretical explanation of stripe formation is much debated. Early work reported stripes in the t-J (at large J/t with 1/r repulsions) and Hubbard (Hartree-Fock) models. However, these stripes were insulating with hole density n<sub>h</sub>$``$1.0, different from the experimental n<sub>h</sub>$``$0.5 stripes. Recently, considerable progress was made when doped stripes were reported by White and Scalapino within the standard t-J model (see also Ref.). In Ref. the analysis was performed at couplings where two holes form d-wave pairs, and the stripes are sometimes described as a condensation of these pairs into a stripe domain-wall. However, experiments are usually interpreted as holes moving freely along site-centered stripes. In addition, the “extended” t-J model with hopping beyond neighboring sites, or the standard t-J model with very small J/t, are needed to reproduce the insulator one-hole photoemission (PES) dispersion. Thus, understanding metallic stripe formation requires further work and searching for stripes in the extended t-J model, particularly in regimes without hole binding and where the absence of phase separation (PS) is not controversial, is important to clarify the driving mechanism for these unusual complex structures.
Building upon previous investigations, in this Letter indications of n<sub>h</sub>$``$0.5 stripes compatible with experiments are reported in the extended and standard t-J models on ladders and square clusters. These stripes do not seem composed of hole pairs (although pairs forming domain-walls may be present at larger J/t than studied here). They also exist in the t-J<sub>z</sub> model and using classical spins, implying that the details of the AF spin background are unimportant for its stabilization. Moreover, our most important result is that the basic stripe “building-block” exists already in the insulator where one-hole wave functions have a complex spin structure with strong AF correlations across-the-hole, resembling the stripe spin correlations found here numerically. These results provide a rationalization for stripe formation built upon the $`one`$ hole properties, in regimes where spin and charge are almost separated.
The extended t-J model used here is defined as
$$\mathrm{H}=\mathrm{J}\underset{\mathrm{𝐢𝐣}}{}(𝐒_𝐢𝐒_𝐣\frac{1}{4}\mathrm{n}_𝐢\mathrm{n}_𝐣)\underset{\mathrm{𝐢𝐦}}{}\mathrm{t}_{\mathrm{𝐢𝐦}}(\mathrm{c}_𝐢^{}\mathrm{c}_𝐦+\mathrm{h}.\mathrm{c}.),$$
where $`\mathrm{t}_{\mathrm{𝐢𝐦}}`$ is t(=1) for nearest-neighbors (NN), t for next NN, and t<sup>′′</sup> for next next NN sites, and zero otherwise. The rest of the notation is standard. The t-J<sub>z</sub> model is obtained by J$``$J<sub>z</sub> and $`𝐒_𝐢𝐒_𝐣`$$``$ $`\mathrm{S}_{}^{\mathrm{z}}{}_{𝐢}{}^{}\mathrm{S}_{}^{\mathrm{z}}{}_{𝐣}{}^{}`$, and t$`<`$0 and t<sup>′′</sup>$`>`$0 are relevant to explain PES data. Here the Density Matrix Renormalization Group (DMRG), Lanczos, and an algorithm using a small fraction of the ladder rung-basis (optimized reduced-basis approximation, or ORBA) are used. Results are presented in (i) the small J/t region with t=t<sup>′′</sup>=0.0, and (ii) small and intermediate J/t with nonzero t and t<sup>′′</sup>. These two regions have similar physics, and the extra hoppings are expected to avoid PS. Intuitively, t,t<sup>′′</sup> increase hole mobility, as reducing J/t does, but also avoid ferromagnetism at small J/t. Note also that no coupling fine-tuning is needed: the results below appear in a robust region of parameter space.
In Fig.1, DMRG and ORBA results for 4$`\times `$N clusters are shown. In Fig.1a the rung density for a 4$`\times `$8 (4$`\times `$12) cluster with 4 (2) holes at small J/t is presented. Cylindrical boundary conditions (CBC) are used i.e. open boundary conditions (OBC) along legs and periodic boundary conditions (PBC) along rungs. The four holes separate into two groups of two holes, surprising result since for a square lattice J<sub>c</sub>=0.2 is the critical value for hole pair binding in the t-J<sub>z</sub> model, and in the t-J model J<sub>c</sub> is expected to be larger. Similar results are found in the t-J<sub>z</sub> model (Fig.1a) and at intermediate J/t but with t$``$0, which increases the hole mobility: Fig.1b with six holes show the formation of three groups of two holes as in Fig.1a. This is not spuriously caused by the OBC along legs, as shown in Fig.1c with results using PBC in both directions. As ORBA starting configuration holes clustered (phase separated) or spread apart (free gas) were used, with PBC or CBC, and in both cases the results converged to the same “stripe” answer.
To study the two-hole state internal structure, in Fig.1d the density distribution of one hole around the other is shown, for one of the 2-hole regions of Fig.1a. The largest density is at two lattice spacings along the rung, and the hole distribution does not resemble a tight d-wave bound state. Similar conclusions were reached for the two holes of Fig.1c. The result is actually compatible with the formation of a short site-centered stripe where the two holes form a closed loop with density 0.5 along a rung. These stripes appear to occupy more than one rung in Figs.1a-c, and thus they could also be labeled as bond-centered. However, this effect seems to arise from stripe tunneling between neighboring rungs, as the one-hole projection suggests (Fig.1d). Similar results regarding half-doped stripe formation were also found on 6$`\times `$6 clusters, as exemplified in Fig.2a where sets of 3 holes form individual n<sub>h</sub>$``$0.5 stripes (invariance under reflexions was assumed along the legs). Overall the results are consistent with Tranquada’s description of stripes. They are also consistent with numerical reports for the standard t-J model, although our interpretation of the results (below) is different.
The half-doped stripes reported here also lead to spin IC. For example, in Fig.1e the spin structure factor is shown for the cases of Figs.1a,b. The peak deviation from ($`\pi `$,$`\pi `$) appears in a robust region of parameter space. The spin IC is understood calculating spin-spin correlations when two holes in, e.g., the cluster of Fig.1c are projected into their most probable location (Fig.2b): a $`\pi `$-shift across-the-stripe can be clearly observed. The across-the-stripe AF correlation strength increases reducing J/t and/or increasing t$`<`$0 and t<sup>′′</sup>$`>`$0 in magnitude.
Results compatible with n<sub>h</sub>$``$0.5 stripes and associated $`\pi `$-shifts appear in other clusters as well. On a cylindrical 6$`\times `$4 cluster with PBC along the long direction, the 3-holes ground-state has characteristics compatible with a doped one-dimensional (1D) closed loop along the PBC direction, with $`\pi `$-shifts across-the-stripe (see Fig.2c where one of the two degenerate most dominant ground-state hole configurations is shown). A h-s-h-s-h-s loop (h=hole, s=spin) provides a pictorial representation of our results, but this configuration is not rigid neither along nor perpendicular to the loop. Density correlations along the stripe (Fig.2c) are actually compatible with a 1D n<sub>h</sub>$``$0.5 system at large on-site U interactions, suggesting that the stripes described here are metallic. No indications of a charge-density-wave along the stripe were found. Note also that spin IC induced by antiferromagnetism across holes also exist $`along`$ $`the`$ $`stripes`$, with wavevector $`\pi `$/2 for a half-doped stripe. This spin IC appears also in half-doped 1D models. For an isolated CuO plane, IC should be present in both directions, although with quite different wavevectors and intensities.
Similar results are found in small square clusters: in the 2-holes 4$`\times `$4 lattice with CBC, a 2-hole stripe forms along the PBC direction. With PBC in both directions, the ground-state resembles a mixture of stripes along both axes and since nonzero t-t<sup>′′</sup> avoids PS, our results are not expected to have the boundary effects recently discussed. Tendency to stripe formation is found even in tilted clusters: the PBC $`\sqrt{18}`$$`\times `$$`\sqrt{18}`$ lattice allows for n<sub>h</sub>$``$0.5 closed loops with 3 holes and such structure has a large ground-state weight (Fig.2d). Precursors of the spin structures in Figs.1,2 appear on 2- and 3-leg ladders as well, e.g. in Fig.3a the 2 holes ground-state dominant hole configuration of a 3$`\times `$6 cluster is shown, with its spin correlations. On 2-leg ladders with many holes, $`\pi `$-shifts appear at small J/t (Fig.3b), and each hole is “confined” to a rung, precursor of a rung stripe as the leg number grows. Spin IC is here found both for the 2-leg (Fig.3c) and 3-leg ladders.
The results thus far suggest that doped stripes can form in spin and hole models using realistic couplings. To gain insight into the mechanism driving this complex structure, consider now the $`one`$ hole problem. Fig.3d shows 4-leg ladder spin correlations around a mobile hole for momentum ($`\pi `$,$`\pi `$). The AF correlations across-the-hole are clearly similar to the correlations around the individual holes composing the stripes. The $`\pi `$-shift characteristic of the stripes exists in the one-hole state not only at ($`\pi `$,$`\pi `$) but at several momenta, and, in this sense, the spin IC exists already at the one-hole level, a remarkable result. Similar conclusions are reached for 3- and 2-leg ladders (Fig.3e), and other momenta such as (0,$`\pi `$). Also on small square clusters robust across-the-hole AF correlations exist for one hole. Although spin IC was found in early studies of the t-J model, and the nontrivial structures as in Fig.3d were noticed before, it was only recently tentatively explained as (local) spin-charge separation, similar to the 1D Hubbard model where spins across holes are antiparallel.
The results shown here lead us to believe that the observed doped stripes are made out of one-hole building-blocks (Fig.3d). In this respect the insulator limit already carries the essential information needed to build the stripes, providing an unexpected potential simple link between undoped and doped cuprates. This is compatible with the behavior of the large energy scale PES pseudogap which can be traced back to the one-hole dispersion of the insulator, suggesting a smooth evolution from the undoped to underdoped regimes.
However, further elaboration is needed since for one-hole the lowest energy is found at $`𝐪`$$``$$`(\pi /2,\pi /2)`$ . Naively, hole pockets at $`(\pi /2,\pi /2)`$ should appear at finite hole density. In addition, across-the-hole AF bonds are weaker at $`(\pi /2,\pi /2)`$ than at momenta such as ($`\pi `$,0) or ($`\pi `$,$`\pi `$), although they are still present. To address this issue let us calculate $`\mathrm{n}_𝐪`$=$`\mathrm{c}_𝐪^{}\mathrm{c}_𝐪`$, i.e. the ground-state hole number with a given momentum $`𝐪`$ (note that $`\mathrm{n}_𝐪`$ includes both coherent and incoherent weight). As example, consider the two-hole problem on the 4$`\times `$6 lattice of Fig.2b. The interesting result in Fig.4a is that
the ground state carries dominant weight at momenta around ($`\pi `$,$`\pi `$), and the one-hole states with this momentum have robust AF correlations across-the-hole (Fig.3d), compatible with our conjecture. There are no indications of small hole-pockets in our studies, and the Fermi surface appears open. In this framework the across-the-hole correlations of the, e.g., ($`\pi `$,$`\pi `$) or ($`\pi `$,0) holes can be “linked”, as pictorially shown in Fig.4b, improving the hole mobility since now they share a large region where they do not need to fight against the spin background to move. Creating a stripe loop also avoids the spin frustration intrinsic of the individual hole states when across-the-hole robust correlations are present (Fig.4b). In addition, our results help understanding better the observed stripe density: for n<sub>h</sub>$``$1 the across-the-hole AF bonds in the stripe direction cannot form and holes do not improve their kinetic energy, while for a very hole diluted stripe the finite-size elementary blocks (Fig.3d) do not touch and cannot have a common spin arrangement. For completeness, in Figs.4c,d the one-hole spectral function is exactly calculated on 4-, 3- and 2-leg ladders at small J/t. Note the remarkable small quasiparticle weight, correlated with a robust across-the-hole AF correlation (see also ). The one-hole states contributing to stripes have exotic properties, including a tendency to spin-charge separation.
Summarizing, indications of n<sub>h</sub>$``$0.5 stripes were found in the extended t-J, t-J<sub>z</sub>, and (at small J/t) in the standard t-J models. The gain of kinetic energy against the loss of AF energy appears enough to stabilize stripes, namely the driving force is a one hole process and the seed for stripes is already present in the insulator. Contrary to most approaches to stripe formation, here the small J/t regime was emphasized. The scenario reported here is a generalization of the 1D spin-charge separation involving individual holons, with the twist that stripes of holons are needed in 2D to avoid frustration. This result is compatible with Zaanen’s picture of stripes as “holons in a row”. Charge and spin could be separated in 2D in more subtle ways than anticipated.
The authors thank R. Eder, S. White and J. Zaanen for useful comments and NSF (DMR-9814350), FAPESP-Brazil, and Fundación Antorchas for support.
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# Inclusive production of 𝜔 mesons at large transverse momenta in 𝜋⁻Be interactions at 515 GeV/𝑐
## I INTRODUCTION
The production of mesons at large transverse momenta ($`p_T`$) in hadronic collisions is a process that can be used to study the phenomenology of Quantum Chromodynamics (QCD) and the parton fragmentation process. Measurements of the spectra of particles produced in the fragmentation of partons influence our understanding of high-$`p_T`$ processes and affect the design and implementation of Monte Carlo event generators. Since mesons constitute the majority of the particles produced in hadronic interactions at high-$`p_T`$, such measurements are of interest across a broad spectrum of experiments in high energy physics. Although much is known about pion production in this regime, the hadroproduction of $`\omega `$ mesons has not been as extensively studied . A comparison of $`\omega `$ to $`\pi ^0`$ production, which reflects the overall ratio of vector meson to pseudoscalar meson production ($`V/P`$), can also be used to sharpen the value of the $`V/P`$ phenomenological parameter used in current Monte Carlo event generators.
We report on $`\omega `$ production in $`\pi ^{}`$Be collisions at 515 GeV/$`c`$, as measured in E706, an experiment which was designed to study the production of direct photons, neutral mesons, and associated particles at high-$`p_T`$ using the Meson West Spectrometer at Fermilab . The apparatus included a charged particle spectrometer consisting of silicon microstrip detectors in the target region and multiwire proportional chambers and straw tube drift chambers downstream of a large aperture analysis magnet . The target consisted of two 0.8 mm thick copper foils, immediately upstream of two cylinders of beryllium, one 3.7 cm long and the other 1.1 cm long. Photons were detected in a 3 m diameter, lead and liquid-argon, sampling electromagnetic calorimeter (EMLAC), located $`9`$ m downstream of the target . The EMLAC readout was subdivided azimuthally into octants, each consisting of interleaved, finely segmented, radial and azimuthal views. The radial views were also used to form a fast high-$`p_T`$ event selection trigger. Trigger decisions were based on global (full octant) and local (sixteen 5.5 mm strips) sums of EMLAC energy weighted to measure transverse momentum .
## II DATA ANALYSIS
The data sample corresponds to an integrated luminosity of $`8.6`$ events/pb. The measurement of the $`\omega `$ meson cross section was based on its decay $`\omega \pi ^0\gamma `$. Although this mode has a branching ratio of only 8.5% , its decay chain to three photons yielded a clear signature. In the following section, we briefly describe the analysis procedures used to extract the $`\omega `$ signal in these data. A more detailed description of the procedures used to select and reconstruct the events, and to correct the data for losses due to trigger inefficiencies and selection requirements, can be found elsewhere .
### A Signal extraction
The invariant mass distribution for $`\pi ^0\gamma `$ pairs, subject to only minimal kinematic cuts, is illustrated in Fig. 1. In our $`\omega `$ study, we defined a $`\pi ^0`$ as a combination of two photons with invariant mass, $`M_{\gamma \gamma }`$, in the range 110 MeV/$`c^2<M_{\gamma \gamma }<165`$ MeV/$`c^2`$ and energy asymmetry \[$`A=|E_1E_2|/(E_1+E_2)`$\] less than 0.75. The $`\omega `$ signal in Fig. 1 consists of a shoulder riding on a steeply falling background. To improve the signal-to-background ratio, we investigated the $`\mathrm{cos}\theta ^{}`$ distribution, where $`\theta ^{}`$ is defined as the decay angle of the $`\pi ^0`$ in the $`\pi ^0\gamma `$ rest frame, relative to the $`\pi ^0\gamma `$ line of flight. Monte Carlo studies showed that the $`\mathrm{cos}\theta ^{}`$ distribution for accidental pairings of $`\pi ^0`$’s and $`\gamma `$’s is peaked near $`\pm 1`$ , whereas the distribution for unpolarized $`\omega `$’s is isotropic. (Since we detected no evidence for $`\omega `$ polarization in our data, we assumed for the purposes of the analysis presented in this paper that the signal was dominated by unpolarized $`\omega `$’s.) The requirement $`|\mathrm{cos}\theta ^{}|<0.6`$ eliminates a large fraction of the combinatorial background , resulting in the distribution shown in the insert in Fig. 1, in which the $`\omega `$ signal stands out much more clearly.
Figure 2 displays the weighted mass distribution of $`\pi ^0\gamma `$ pairs in the vicinity of the $`\omega `$ mass. These data have also been subjected to additional kinematic constraints including a $`\pi ^0`$ sideband subtraction which Monte Carlo studies showed to be effective in removing $`\eta \gamma `$ combinations that peaked in the vicinity of the $`\omega `$ signal . The effect of these and other kinematic criteria was estimated using the Monte Carlo simulation described below.
Fits to the $`\pi ^0\gamma `$ mass spectra for different intervals of $`p_T`$ were used to extract the contribution of the $`\omega `$ signal as a function of $`p_T`$. A Gaussian shape for the signal, combined with a third order polynomial for the background, yielded a reasonable description of the data.
### B Monte Carlo
A full event simulation was used to evaluate most corrections to the cross section. This simulation relied on the herwig event generator and a geant Monte Carlo simulation of our apparatus . Events containing leading $`\omega `$ mesons were generated and reconstructed using our standard reconstruction package . The ratio of reconstructed to generated events was parameterized as a function of $`p_T`$, fitted with a Theta function convolved with a Gaussian (Fig. 3), and applied to the $`\omega `$ spectrum. Due to the relatively large opening angle for photons in $`\omega `$ decay, and the resulting spatial spread of these photons in the EMLAC, the decay products infrequently satisfied the E706 high-$`p_T`$ trigger. The overall detection efficiency for the $`\omega `$ (in the $`\pi ^0\gamma `$ decay mode) was $`20`$% at high-$`p_T`$, and differed between the global and local triggers at lower-$`p_T`$.
### C Systematic uncertainties
We estimate the combined systematic uncertainty in the $`\omega `$ cross section due to reconstruction efficiency and normalization to be 15% . The energy response of the electromagnetic calorimeter was calibrated using $`\pi ^0`$, $`\eta `$, $`\omega `$, and $`J/\psi `$ mass peaks. The energy scale uncertainty was determined to be less than 0.5% , and contributed between 5% and 12% (as a function of $`p_T`$) to the systematic uncertainty of the $`\omega `$ cross section.
The uncertainty in the fitted background was estimated by using different parameterizations for the background, and by varying the $`\pi ^0\gamma `$ mass range and the bin sizes used in the fit. We estimate a 10% contribution to the systematic uncertainty in the determination of the signal due to the fitting procedure. The uncertainty in the trigger corrections was estimated by comparing the $`\omega `$ cross section obtained using samples selected with different triggers. This uncertainty ranged from 17% for $`\omega `$ mesons with $`p_T`$ of 3.5 GeV/$`c`$ to 2% for 8 GeV/$`c`$. These quoted systematic uncertainties do not explicitly incorporate contributions due to the possibility of $`\omega `$ polarization, however, as already stated, we detected no evidence for $`\omega `$ polarization over the range $`|\mathrm{cos}\theta ^{}|<0.6`$.
The systematic uncertainties, combined in quadrature, are quoted with the cross sections in Table I. The overall systematic uncertainty on the $`\omega `$ cross section was 30% at $`p_T=`$3.5 GeV/$`c`$, 22% at 5 GeV/$`c`$, and 20% at 8 GeV/$`c`$.
## III Results
Table I lists our measured inclusive invariant cross section for $`\omega `$ meson production in $`\pi ^{}`$Be interactions at 515 GeV/$`c`$ along with statistical and systematic uncertainties. The results are binned in $`p_T`$ from 3.5 to 8 GeV/$`c`$, and averaged over the range of our acceptance in rapidity ($`0.5<y_{\mathrm{CM}}<0.75`$). The $`\omega `$ cross section is also displayed in Figure 4, and compared with expectations from pythia and herwig . The predictions from both Monte Carlos are substantially smaller than the measured $`\omega `$ cross section. Comparison between the relative yields of $`\omega `$’s originating in the Be and Cu target materials resulted in the value $`\alpha =1.12\pm 0.07\pm 0.07`$ using the parameterization $`\sigma _AA^\alpha `$. The Monte Carlo results have been adjusted for this nuclear effect.
Figure 5 and Table I display the relative yields of $`\omega `$ and $`\pi ^0`$ mesons measured in E706, in terms of the ratio of the inclusive differential cross sections as a function of $`p_T`$. The prediction from the pythia generator is consistent with our measured ratio. The ratio from herwig is much smaller than both our measurement and the result from pythia. We include for comparison three previous results on the $`\omega `$ to $`\pi ^0`$ cross section ratio for incident protons on p, Be, and C targets . These earlier measurements were integrated over $`p_T`$ and displayed at their minimum $`p_T`$ value.
The $`\omega `$ to $`\pi ^0`$ ratio can be used to determine the value of $`V/P`$. The quark content of both $`\omega `$ and $`\pi ^0`$ is the same, rendering their production ratio insensitive to beam and target composition. Corrections to $`V/P`$ to account for indirect production were determined using both pythia and herwig . In both cases, the data from E706 were found to require only a relatively small correction for indirect contributions from the decay of higher mass hadrons. We note, however, that in pythia almost all $`\omega `$ mesons are produced directly, which seems to us an extreme assumption, and the herwig $`\omega `$ to $`\pi ^0`$ ratio differs greatly from our data. Nevertheless, if we use these Monte Carlos to correct our measured $`\omega `$ to $`\pi ^0`$ ratio for indirect production, the resulting $`V/P`$ values are $`1.2\pm 0.1`$ using pythia and $`0.9\pm 0.1`$ using herwig.
###### Acknowledgements.
We thank the management and staff of Fermilab, the U. S. Department of Energy, the National Science Foundation, including its Office of International Programs, and the Universities Grants Commission of India, for their support of this research. We are also pleased to acknowledge the contributions of our colleagues on Fermilab experiment E672 to this and other aspects of E706.
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# Anomalous WW𝛾 Vertex in 𝛾𝑝 Collision
## I Introduction
Recently there have been intensive studies to test the deviations from the Standard Model (SM) at present and future colliders. The investigation of three gauge boson couplings plays an important role to manifest the non abelian gauge symmetry in standard electroweak theory. The precision measurement of the triple vector boson vertices will be the crucial test of the structure of the SM.
It is known that present collider measurements at LEP2 and Tevatron can not probe anomalous triple gauge boson self couplings to precision better than O($`10^1`$). Further analyses of $`WW\gamma `$ vertex has been given by several papers for ep collider HERA , projected future colliders LHC and linear electron-positron collider (LC). After LC is constructed $`\gamma e`$, $`\gamma \gamma `$, linac-ring type $`ep`$ and $`\gamma p`$ modes should be discussed and may work as complementary to basic colliders. The Linear Collider design at DESY is the only one that can be converted into an ep collider. $`\gamma p`$ collider mode is an additional advantage of this linac-ring type ep collider where real gamma beam is obtained by the Compton backscattering of laser photons off linear electron beam. Estimations show that the luminosity for $`\gamma `$p collision turns out to be of the same order as the one for ep collision due to the fact that $`\sigma _p>>\sigma _\gamma `$ where $`\sigma _p`$ and $`\sigma _\gamma `$ are transverse sizes of proton and photon bunches at collision point. Since most of the photons are produced at high energy region in the Compton backscattering the cross sections are about one order of magnitude larger than parental ep collider for photoproduction processes. According to present project at DESY 500 GeV electrons from linear electron beam are allowed to collide 820 GeV protons from HERA ring. Parameters of this LC+HERAp and its $`\gamma `$p option are shown in Table I . Therefore such kind of high energy $`\gamma `$p colliders will possibly give additional information to linac-ring type ep colliders for a variety of processes. In this paper we examine the potential of future LC+HERAp based $`\gamma `$p collider to probe anomalous WW$`\gamma `$ coupling and compare the results with those from its basic ep collider and other projected future colliders.
## II Lagrangian and Cross Sections
C and P parity conserving effective lagrangian for two charged W-boson and one photon interaction can be written following the papers
$`L`$ $`=e(W_{\mu \nu }^{}W^\mu A^\nu W^{\mu \nu }W_\mu ^{}A_\nu +\kappa W_\mu ^{}W_\nu A^{\mu \nu }+{\displaystyle \frac{\lambda }{M_W^2}}W_{\rho \mu }^{}W_\nu ^\mu A^{\mu \rho })`$ (1)
where
$`W_{\mu \nu }=_\mu W_\nu _\nu W_\mu `$ (2)
and dimensionless parameters $`\kappa `$ and $`\lambda `$ are related to the magnetic dipol and electric quadrupole moments. For $`\kappa =1`$ and $`\lambda =0`$ Standard Model is restored. In momentum space this has the following form with momenta $`W^+(p_1)`$,$`W^{}(p_2)`$ and $`A(p_3)`$
$`\mathrm{\Gamma }_{\mu \nu \rho }(p_1,p_2,p_3)=`$ $`e[g_{\mu \nu }(p_1p_2{\displaystyle \frac{\lambda }{M_W^2}}[(p_2.p_3)p_1(p_1.p_3)p_2])_\rho `$ (6)
$`+g_{\mu \rho }(\kappa p_3p_1+{\displaystyle \frac{\lambda }{M_W^2}}[(p_2.p_3)p_1(p_1.p_2)p_3])_\nu `$
$`+g_{\nu \rho }(p_2\kappa p_3{\displaystyle \frac{\lambda }{M_W^2}}[(p_1.p_3)p_2(p_1.p_2)p_3])_\mu `$
$`+{\displaystyle \frac{\lambda }{M_W^2}}(p_{2\mu }p_{3\nu }p_{1\rho }p_{3\mu }p_{1\nu }p_{2\rho }])]`$
where $`p_1+p_2+p_3=0`$. There are three Feynman diagrams for the subprocess $`\gamma q_iWq_j`$ and only t-channel W exchange graph contributes $`WW\gamma `$ vertex. One should note that $`\gamma p`$ collision isolates $`WW\gamma `$ coupling but many processes in $`e^+e^{}`$, pp and ep collisions include mixtures of $`WW\gamma `$ and WWZ couplings.
The unpolarized differential cross section for the subprocess $`\gamma q_iWq_j`$ can be obtained using helicity amplitudes from summing over the helicities
$`{\displaystyle \frac{d\widehat{\sigma }}{d\widehat{t}}}={\displaystyle \frac{2}{\widehat{s}M_W^2}}{\displaystyle \frac{\beta }{64\pi \widehat{s}}}{\displaystyle \underset{\lambda _\gamma \lambda _W}{}}{\displaystyle \frac{1}{2}}M_{\lambda _\gamma \lambda _W}^2`$ (7)
where
$`M_{\lambda _\gamma \lambda _W}={\displaystyle \frac{e^2}{\sqrt{2}\mathrm{sin}\theta _W}}{\displaystyle \frac{\widehat{s}}{\widehat{s}+M_W^2}}\sqrt{\beta }A_{\lambda _\gamma \lambda _W},\beta =1{\displaystyle \frac{M_W^2}{\widehat{s}}}`$ (8)
and $`\theta _W`$ is the Weinberg angle.
For the signal we are considering a quark jet and on-shell W with leptonic decay mode
$`\gamma pW^{}+jet\mathrm{}+p_T^{miss}+jet,\mathrm{}=e,\mu `$ (9)
In this mode charged lepton and the quark jet are in general well separated and the signal is in principle free of background of SM.
The total cross section for the subprocess $`\gamma q_iWq_j`$ can be divided into two parts, direct, and resolved-photon production
$`\widehat{\sigma }=\widehat{\sigma }_{dir}+\widehat{\sigma }_{res}`$ (10)
The direct part is given as follows
$`\widehat{\sigma }=`$ $`{\displaystyle \frac{\alpha G_FM_W^2}{\sqrt{2}\widehat{s}}}|V_{q_iq_j}|^2\{(|e_q|1)^2(12\widehat{z}+2\widehat{z}^2)\mathrm{log}({\displaystyle \frac{\widehat{s}M_W^2}{\mathrm{\Lambda }^2}})[(12\widehat{z}+2\widehat{z}^2)`$ (14)
$`2|e_q|(1+\kappa +2\widehat{z}^2)+{\displaystyle \frac{(1\kappa )^2}{4\widehat{z}}}{\displaystyle \frac{(1+\kappa )^2}{4}}]\mathrm{log}\widehat{z}+[(2\kappa +{\displaystyle \frac{(1\kappa )^2}{16}}){\displaystyle \frac{1}{\widehat{z}}}`$
$`+({\displaystyle \frac{1}{2}}+{\displaystyle \frac{3(1+|e_q|^2)}{2}})\widehat{z}+(1+\kappa )|e_q|{\displaystyle \frac{(1\kappa )^2}{16}}+{\displaystyle \frac{|e_q|^2}{2}}](1\widehat{z})`$
$`{\displaystyle \frac{\lambda ^2}{4\widehat{z}^2}}(\widehat{z}^22\widehat{z}\mathrm{log}\widehat{z}1)+{\displaystyle \frac{\lambda }{16\widehat{z}}}(2\kappa +\lambda 2)[(\widehat{z}1)(\widehat{z}9)+4(\widehat{z}+1)\mathrm{log}\widehat{z}]\}`$
where $`\widehat{z}=M_W^2/\widehat{s}`$ and $`\mathrm{\Lambda }^2`$ is cut off scale in order to regularize $`\widehat{u}`$-pole of the colinear singularity for massles quarks. In the case of massive quarks there is no need such a kind of cut. The resolved-photon production cross section can be calculated using Breit-Wigner formula for $`q_\gamma q_pW`$ fusion process
$`\widehat{\sigma }(q_i\overline{q}_jW)={\displaystyle \frac{\pi \sqrt{2}}{3}}G_Fm_W^2|V_{ij}|^2\delta (x_ix_js_{\gamma p}m_W^2)`$ (15)
where $`V_{ij}`$ is the Cabibbo-Kobayashi-Maskawa (CKM) matrix. For the integrated cross sections we need parton distribution functions inside the photon and proton. The photon structure function $`f_{q/\gamma }`$ consists of perturbative pointlike parts and hadronlike parts. The pointlike part can be calculated in the leading logarithmic approximation and is given by the expression
$`f_{q/\gamma }^{LO}(x,Q_\gamma ^2)={\displaystyle \frac{3\alpha e_q^2}{2\pi }}[x^2+(1x)^2]\mathrm{log}{\displaystyle \frac{Q_\gamma ^2}{\mathrm{\Lambda }^2}}`$ (16)
where $`e_q`$ is the quark charge. For the integrated total cross section over the quark distributions inside the proton, photon and the spectrum of backscattered laser photon the following result can be obtained easily
$`\sigma _{res}=\sigma _0{\displaystyle _{m_W^2/s}^{0.83}}𝑑x{\displaystyle _x^{0.83}}{\displaystyle \frac{dy}{xy}}f_{\gamma /e}(y)f_{q_i/p}({\displaystyle \frac{m_W^2}{xs}},Q_p)[f_{q_j/\gamma }({\displaystyle \frac{x}{y}},Q_\gamma ^2)f_{q_j/\gamma }^{LO}({\displaystyle \frac{x}{y}},Q_\gamma ^2)]`$ (17)
with
$`\sigma _0={\displaystyle \frac{\pi \sqrt{2}}{3s}}G_Fm_W^2|V_{ij}|^2`$ (18)
Since the contribution from the pointlike part of the photon structure function was already taken into account in the calculation of the direct part it was subtracted from $`f_{q/\gamma }(x,Q_\gamma ^2)`$ in the above formula to avoid double counting on the leading logarithmic level.
In a similar way direct part of the integrated cross section can be obtained
$`\sigma _{dir}={\displaystyle _{\tau _{min}}^{0.83}}𝑑\tau {\displaystyle _{\tau /0.83}^1}{\displaystyle \frac{dx}{x}}f_{\gamma /e}(\tau /x)f_{q/p}(x,Q^2)\widehat{\sigma }(\tau s)`$ (19)
with $`\tau _{min}=(M_W+M_q)^2/s`$
The sum of resolved and direct contribution, in principle, should not depend on value of the parameter $`\mathrm{\Lambda }`$. The effects of $`\mathrm{\Lambda }`$ on the cross sections and contributions of direct and resolved parts are given in Table II. In Table III integrated total cross sections times branching ratio of $`W\mu \nu `$ and corresponding number of events are shown for various values of $`\kappa `$ and $`\lambda `$. Through the calculations proton structure functions of Martin, Robert and Stirling (MRS A) and photon structure functions of Drees and Grassie (DG) have been used with $`Q^2=M_W^2`$. Here number of events has been calculated using
$`N=\sigma (\gamma pW+Jet)B(W\mu \nu )AL_{int}`$ (20)
where we take the acceptance in the muon channel A and integrated luminosity $`L_{int}`$ as 65% and 200pb<sup>-1</sup>. To give an idea about the comparison with corresponding ep collider the cross sections obtained using Weizsäcker-Williams approximation are also shown on the same table. As the cross section $`\sigma (\gamma pW+Jet)`$ we have considered the sum of $`\sigma (\gamma uW^+d)`$, $`\sigma (\gamma \overline{d}W^+\overline{u})`$, $`\sigma (\gamma \overline{s}W^+\overline{c})`$, $`\sigma (\gamma uW^+s)`$, $`\sigma (\gamma \overline{s}W^+\overline{u})`$, and $`\sigma (\gamma \overline{d}W^+\overline{c})`$. As shown from Table III the cross sections using backscattered laser photons are considerably larger than the case of corresponding ep collision. We assume that the form factor structure of $`\kappa 1`$ and $`\lambda `$ do not depend on the momentum transfers at the energy region considered. Then anomalous terms containing $`\kappa `$ grow with $`\sqrt{\widehat{s}}/M_W`$ and terms with $`\lambda `$ rise with $`\widehat{s}/M_W^2`$. Deviation $`\mathrm{\Delta }\kappa =\kappa 1=1`$ from SM value changes the total cross sections 30-70% whereas the $`\mathrm{\Delta }\lambda =\lambda =1`$ gives 80% changes. Therefore high energy will improve the sensitivity to anomalous couplings. For comparison with HERA energy $`\sqrt{s}=314`$ GeV the similar results would be 20-40% for $`\mathrm{\Delta }\kappa =1`$ and 5% for $`\mathrm{\Delta }\lambda =1`$.
It is important to see how the anomalous couplings change the shape of the transverse momentum distributions of the final quark jet. For this reason we use the following formula:
$`{\displaystyle \frac{d\sigma }{dp_T}}=`$ $`2p_T{\displaystyle _y^{}^{y^+}}𝑑y{\displaystyle _{x_a^{min}}^{0.83}}𝑑x_af_{\gamma /e}(x_a)f_{q/p}(x_b,Q^2)({\displaystyle \frac{x_ax_bs}{x_as2m_TE_pe^y}}){\displaystyle \frac{d\widehat{\sigma }}{d\widehat{t}}}`$ (21)
where
$`y^{}=\mathrm{log}[{\displaystyle \frac{0.83s+m_q^2M_W^2}{4m_TE_p}}\{({\displaystyle \frac{0.83s+m_q^2M_W^2}{4m_TE_p}})^2{\displaystyle \frac{0.83E_e}{E_p}}\}^{1/2}]`$ (22)
$`x_a^{(1)}={\displaystyle \frac{2m_TE_pe^ym_q^2+M_W^2}{s2m_TE_ee^y}},x_a^{(2)}={\displaystyle \frac{(M_W+m_q)^2}{s}}`$ (23)
$`x_a^{min}=MAX(x_a^{(1)},x_a^{(2)}),x_b={\displaystyle \frac{2m_TE_ex_ae^ym_q^2+M_W^2}{x_as2m_TE_pe^y}}`$ (24)
with
$`\widehat{s}=x_ax_b,\widehat{t}=m_q^22E_ex_am_Te^y,\widehat{u}=m_q^2+M_W^2\widehat{s}\widehat{t}`$ (25)
$`m_T^2=m_q^2+p_T^2`$ (26)
The $`p_T`$ spectrum B($`W\mu \nu )\times d\sigma /dp_T`$ of the quark jet is shown in Fig. 1 for various $`\kappa `$ and $`\lambda `$ values at LC+HERAp based $`\gamma `$p collider. Similar distribution is given in Fig. 2 for Weizsäcker-Williams Approximation that covers the major contribution from ep collision. It is clear that the cross section at large $`p_T`$ is quite sensitive to anomalous $`WW\gamma `$ couplings. As $`\lambda `$ increases the cross section grows more rapidly when compared with $`\kappa `$ dependence at high $`p_T`$ region $`p_T>100`$ GeV. The cross sections with real gamma beam are one order of magnitude larger than the case of WWA. Comparison between two figures also shows that the curves become more separable as $`\widehat{s}`$ gets large.
## III Sensitivity to Anomalous Couplings
We can estimate sensitivity of LC+HERAp based $`\gamma `$p collider to anomalous couplings by assuming the 0.02 combined systematic error in the luminosity measurement and detector acceptance for the integrated luminosity value of 200 pb<sup>-1</sup>. We use the simple $`\chi ^2`$-criterion to obtain sensitivity as follows
$`\chi ^2={\displaystyle \underset{i=bins}{}}({\displaystyle \frac{X_iY_i}{\mathrm{\Delta }_i^{exp}}})^2`$ (27)
$`X_i={\displaystyle _{V_i}^{V_{i+1}}}{\displaystyle \frac{d\sigma ^{SM}}{dV}}𝑑V,Y_i={\displaystyle _{V_i}^{V_{i+1}}}{\displaystyle \frac{d\sigma ^{NEW}}{dV}}𝑑V`$ (28)
$`\mathrm{\Delta }_i^{exp}=X_i\sqrt{\delta _{stat}^2+\delta _{sys}^2},V=p_T`$ (29)
We have divided $`p_T`$ region of the final quark into equal pieces for binning procedure and have considered at least 10 events in each bin. For sensitivity, the number of $`W^+`$ and $`W^{}`$ events given with their branching ratios in the $`\mu \nu `$ and $`e\nu `$ channels has been taken into account. Using the above formula the limits on the $`\mathrm{\Delta }\kappa `$ and $`\lambda `$ are given in Table IV for the deviation of the cross section from the Standard Model value at 68% and 95% confidence levels with and without systematic error. On the ground of comparison we give the limits from ep collider in Table V using Weizsäcker-Williams approximation.
From these tables we see that $`\gamma p`$ mode of LC+HERAp probes $`\mathrm{\Delta }\kappa `$ and $`\lambda `$ with better sensitivity than present colliders and comparable with LHC in the case of $`\mathrm{\Delta }\kappa `$ but worse than linear $`e^{}e^+`$ collider. The advantage of the process $`\gamma pWj`$ is that it probes the $`WW\gamma `$ couplings independently of $`WWZ`$ effects. After further improvement of energy and luminosity, linac-ring type $`\gamma p`$ collider possibly will give complementary information to LHC and LC.
## ACKNOWLEDGMENTS
Authors are grateful to the members of the AUHEP group and S. Sultansoy for drawing our attention to anomalous couplings.
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# Second-harmonic interferometric spectroscopy of the buried Si(111)-SiO2 interface
## Abstract
The second-harmonic interferometric spectroscopy (SHIS) which combines both amplitude (intensity) and phase spectra of the second-harmonic (SH) radiation is proposed as a new spectroscopic technique being sensitive to the type of critical points (CP’s) of combined density of states at semiconductor surfaces. The increased sensitivity of SHIS technique is demonstrated for the buried Si(111)-SiO<sub>2</sub> interface for SH photon energies from 3.6 eV to 5 eV and allows to separate the resonant contributions from $`E_0^{}/E_1`$, $`E_2`$ and $`E_1^{}`$ CP’s of silicon.
Second-harmonic generation (SHG) is inherently sensitive to surface and interface properties of centrosymmetric materials. Recently, the spectroscopy of the second-harmonic (SH) intensity has been proved as a promising probe of surfaces and interfacial layers and intensively employed in numerous works for oxidized , reconstructed and H-terminated silicon surfaces. The resonances of the SH intensity are attributed in these cases to direct interband electron transitions. By analogy with the spectrum of the linear dielectric function $`\epsilon (\omega )`$ the spectrum of the quadratic susceptibility $`\chi ^{(2)}(2\omega )`$ of such semiconductor surfaces could be expressed as the superposition of several van Hove singularities (critical points (CP’s) of the combined density of states) $`\chi _m^{(2)}(2\omega )(2\omega \omega _m+i\mathrm{\Gamma }_m)^n`$ with threshold frequencies $`\omega _m`$ and broadenings $`\mathrm{\Gamma }_m`$. The exponent $`n`$ reflects the dimensionality of CP: $`n=1/2,0`$ (logarithmic),$`1/2,1`$ for 3D, 2D, 1D and excitonic CP, respectively. Although the line shapes $`\chi ^{(2)}(2\omega )`$ are quite different for various $`n`$, in most cases a large number of adjustable parameters makes the determination of the type of CP solely from the SH intensity spectrum doubtful, and most of authors interpret the SHG spectroscopy data within excitonic CP line shape.
The single-beam SH interferometry traces back to the mid-1960s and is conventionally used for the determination of the phase of the quadratic susceptibility of adsorbate molecules and their absolute direction and for separation of the SH contributions from thin films and their substrates . Another use of SH phase measurements is a homodyne mixing technique to improve the signal-to-noise ratio for surface SHG probe . The use of external and internal homodynes for dc or low-frequency modulated electric-field-induced SHG allows to measure the in-plane spatial distribution of the electric field vector with micron resolution or to visualize weak nonlinear contributions. Further developments of the SH interferometry are the frequency-domain interferometric SH spectroscopy exploring the broad bandwidth of femtosecond laser pulses, and the hyper-Rayleigh scattering interferometry using the correlation of fluctuations in linear and nonlinear optical properties of thin inhomogeneous films.
In this Letter, a modification of the SH spectroscopy - the SH interferometric spectroscopy (SHIS) which combines both amplitude (intensity) SH spectroscopy and SH interferometry is proposed. The combination of the phase and amplitude SH spectra, extracted from the SHIS data, is shown to be sensitive to type of CP even for systems with interfering SH contributions from close electronic resonances. Additionally SHIS allows to avoid the sign uncertainty of $`\mathrm{Re}(\chi ^{(2)})`$ inherent in the conventional spectroscopy of the SHG intensity. The spectral dependence of the phase and amplitude of the SH waves from the buried Si(111)-SiO<sub>2</sub> interface is measured using the SHIS technique in the spectral range in the vicinity of silicon $`E_2`$ CP. In contrast to $`E_0^{}/E_1`$ CP revealing the excitonic type, the family of $`E_2`$ CP’s of the bulk silicon demonstrates the 2D type in linear response. The resonant contributions to the quadratic susceptibility from $`E_0^{}/E_1`$, $`E_2`$ and $`E_1^{}`$ silicon CP’s are extracted within the simple phenomenological model which accounts the complex Green’s function corrections for the SH wave generation.
The scheme of the SHIS setup is shown in Fig.1(a). The p-polarized output of a tunable nanosecond parametric generator/amplifier laser system (Spectra-Physics MOPO 710) operating in the interval of 490 - 690 nm is focused onto the sample at an angle of incidence of $`45^{}`$. The SH signal is detected by a monochromator, a photomultiplier tube (PMT) and an electronic peak-hold detector. To normalize the SH intensity spectrum over the laser fluence and the spectral sensitivity of the optical detection system a SHG intensity reference channel is used with a slightly wedged z-cut quartz plate and with the detection system identical to the one in the sample channel. The (phase-)reference sample is chosen (i) to be thin enough to avoid Maker fringes in the SH response during tuning the fundamental wavelength $`\lambda _\omega `$, (ii) to be optical inactive for conservation the polarization state of the fundamental radiation while transmitting through it, (iii) to have no resonance features in the tuning region of both the fundamental and SH waves. Therefore the 1 mm-thick plate of fused quartz coated with a 30 nm-thick indium tin oxide (ITO) film is chosen as a reference. The SH interferogram is obtained by translating the reference along the fundamental laser beam varying the distance l between the reference and the sample. The SH signals from the reference, $`I_r^{2\omega }`$, and from the sample, $`I_s^{2\omega }`$, are monitored separately by inserting appropriate filters (yellow or UV, respectively) between the reference and the sample. $`I_r^{2\omega }`$ is adjusted with the angle of incidence of the fundamental beam at the ITO phase reference. The detected SH intensity $`I^{2\omega }`$ is the result of interference of the SH waves from the reference, $`E_r^{2\omega }`$, and from the sample, $`E_s^{2\omega }`$:
$$I^{2\omega }=\frac{c}{8\pi }|E_r^{2\omega }(l)+E_s^{2\omega }|^2=I_r^{2\omega }(l)+I_s^{2\omega }+2\alpha \sqrt{I_r^{2\omega }(l)I_s^{2\omega }}\mathrm{cos}\left(2\pi \frac{l}{L}+\mathrm{\Phi }_{rs}\right),$$
(1)
where $`L=\lambda _\omega (2\mathrm{\Delta }n)^1`$ is the period of SH interferogram with $`\mathrm{\Delta }n=n_{2\omega }n_\omega `$ describing the air dispersion, and $`\alpha <1`$ indicates the laser coherence. The position-dependent phase shift $`2\pi l/L`$ between $`E_r^{2\omega }`$ and $`E_s^{2\omega }`$ comes from the different refractive indices of air for the fundamental and SH waves. The spectral dependences of $`\chi ^{(2)}`$ of the reference and the sample as well as the complex Green’s function corrections for the SH wave produce a position-independent phase shift $`\mathrm{\Phi }_{rs}(\lambda _\omega ,\lambda _{2\omega })`$. The dependence $`I_r^{2\omega }(l)`$ is described by the conventional formula for focused Gaussian beams. The spectrum of the phase of the SH wave from the ITO film, $`\mathrm{\Phi }_r\mathrm{Arg}(E_r^{2\omega })`$, is measured using the 1 mm-thick backside-immersed y-cut quartz as a sample since the phase of the SH wave from the quartz surface is spectrally independent in the whole used spectral region. The $`\mathrm{\Phi }_r`$ spectrum of the ITO film appears to be a constant within the error bars and $`I_r^{2\omega }`$ gradually increases with decreasing $`\lambda _\omega `$ without any resonance features.
The samples are natively oxidized p-doped Si(111) wafers with resistivity of $`10\mathrm{\Omega }cm`$. SHIS has been performed at the maximum of the azimuthal SH rotational anisotropy for the p-in, p-out polarization combination of the fundamental and SH waves. Figure 1(b) shows typical SH interference patterns measured for different $`\lambda _\omega `$. The spectral dependence of the period $`L`$ due to the air dispersion and clear changes in the contrast of the patterns due to the distance dependence of $`I_r^{2\omega }`$ in the focused laser beam are seen. The fit of the set of SH interference patterns by Eq.(1) with $`\mathrm{\Phi }_{rs}`$, $`I_s^{2\omega }`$, $`L`$, and $`\alpha `$ as adjustable parameters leads to the spectra of $`\mathrm{\Phi }_{rs}`$, $`I_s^{2\omega }`$ and $`L`$ shown in Fig. 2(a) and 2(b) and in the inset of Fig. 1, respectively. To emphasize the spectral features of $`I_s^{2\omega }(2\omega )`$, we combine the fitted intensity spectrum with the $`I_s^{2\omega }(2\omega )`$ dependence measured directly with a fine resolution in SH photon energy. $`\mathrm{\Phi }_{rs}`$ increases approximately by 1.2 radians within the interval of 4.2-4.6 eV and decreases outside this energy region. A small, but reliable, non-monotonic feature is seen at 3.8 - 4.0 eV. The $`I_s^{2\omega }`$ spectrum has pronounced peaks centered approximately at 3.9 eV and 4.3 eV. The position of the 4.3 eV resonance is close to $`E_2`$ CP and we associate the observed features of SH phase and intensity spectra for energies between 4.1 and 4.6 eV with direct interband electron transitions at $`E_2`$ CP of Si .
The relative phase $`\mathrm{\Phi }_{rs}`$ measured in SHIS is given by:
$$\mathrm{\Phi }_{rs}=\mathrm{\Phi }_s(\mathrm{\Phi }_r+\mathrm{Arg}(R_{2\omega })),$$
(2)
where $`R_{2\omega }`$ is the Fresnel reflection factor of the p-polarized SH radiation from the Si-SiO<sub>2</sub> interface. The phase $`\mathrm{\Phi }_s\mathrm{Arg}(E_s^{2\omega })`$ originates from both complex surface $`\chi ^{(2),S}`$ and bulk quadrupole $`\chi ^{(2),BQ}`$ quadratic susceptibilities as well as from Green’s function corrections :
$$E_s^{2\omega }=G_{}\chi _{}^{(2)}+G_{}\chi _{}^{(2)},$$
(3)
where $`G_{}`$ and $`G_{}`$ are the Green’s function corrections for the generation and the propagation of in-plane and normal components of $`E_s^{2\omega }`$, and $`\chi _{}^{(2)}`$ and $`\chi _{}^{(2)}`$ are the corresponding effective components of $`\chi ^{(2)}`$. $`\chi _{}^{(2)}`$ and $`\chi _{}^{(2)}`$ are the linear combinations of $`\chi ^{(2),S}`$ and $`\chi ^{(2),BQ}`$ components with nonresonant coefficients depending only on the fundamental wavevector and taking into account the geometry of the nonlinear interaction. This allows to consider the spectral dependences of $`\chi _{}^{(2)}`$ and $`\chi _{}^{(2)}`$ as a superposition of two-photon resonances for different CP’s:
$$\chi _\alpha ^{(2)}(2\omega )=B\underset{m}{}f_m^\alpha \mathrm{exp}(i\varphi _m^\alpha )(2\omega \omega _m+i\mathrm{\Gamma }_m)^n,$$
(4)
where $`\alpha =,`$, and $`m`$ numerates the CP resonances. The oscillator strengths $`f_m^\alpha `$ are supposed to be real numbers. For the sake of simplicity, a slight spectral dependence of the term $`B`$ including the Si resonances with threshold energies below 1.5 eV is neglected, and $`\varphi _m^\alpha `$ are integer multiples of $`\pi /2`$ defining the type of CP. The solid lines in Fig.2 show the fit of $`\mathrm{\Phi }_{rs}`$ and $`I_s^{2\omega }`$ spectra by Eq.(2) and $`|E_s^{2\omega }|^2`$ from Eq.(3), respectively, with the Si dispersion data from Ref.\[\] and expressions for $`G_{}`$ and $`G_{}`$ from Ref.\[\]. Five resonant contributions are included into the fit. The first resonance, centered at $`\omega _1=3.45`$ eV, has excitonic line shape ($`n=1,\varphi _1=0`$) and corresponds to the direct electron transitions at $`E_0^{}/E_1`$ CP. The second resonance at 3.97 eV with 1D maximum line shape ($`n=1/2,\varphi _2=0`$) has no equivalent in the band structure of crystalline bulk silicon. However, a resonance in the close energy interval has been recently observed at the Si(001)-SiO<sub>2</sub> interface and could be associated with transition in Si atoms located at the interface with reduced lattice symmetry. The strong resonant features in the vicinity of 4.3 eV are formed by interference of two resonances centered at $`\omega _3=4.12`$ eV and at $`\omega _4=4.34`$ eV with almost equal amplitudes ($`f_40.9f_3`$). It is mostly reasonable to attribute these peaks to transitions at $`E_2(X)`$ and $`E_2(\mathrm{\Sigma })`$ CP’s. These resonances are fitted with 2D minimum ($`\varphi _3=0`$) and 2D maximum ($`\varphi _4=\pi `$) line shapes ($`\chi _m^{(2)}\mathrm{ln}(2\omega \omega _m+i\mathrm{\Gamma }_m)`$) by analogy with the linear case. Note, that $`\omega _3`$ and $`\omega _4`$ are approximately 0.1 eV red-shifted from resonances of linear $`\chi ^{(1)}`$. This allows to interpret also 4.12 eV-resonance as a contribution from $`E_0`$ CP, which is normally very weak in the linear response. The best representation of the data is obtained with a 2D minimum line shape of the last resonance, centered at 5.15 eV, which can be associated with electron transitions located near $`E_1^{}`$ CP. The error bars for the central frequencies are approximately 0.03 eV and mostly attributed to the relative weight of $`\chi _{}^{(2)}`$ and $`\chi _{}^{(2)}`$ being unresolvable from our data. Excitonic line shapes for all resonances (dotted lines in Fig.2) fit the $`I_s^{2\omega }`$ spectrum with almost the same quality as the CP model, but fit the $`\mathrm{\Phi }_{rs}`$ spectrum obviously worse.
Summarizing, the general scheme of the second-harmonic interferometric spectroscopy is presented. The phase and amplitude of the SH wave from the buried Si(111)-SiO<sub>2</sub> interface are measured simultaneously using the SHIS technique in the interval of SH photon energies from 3.6 eV to 5 eV. The contributions of interband transitions located at $`E_0^{}/E_1`$, $`E_2`$ and $`E_1^{}`$ Si critical points are separated and sensitivity of SHIS to CP line shapes is shown.
This work was supported by the Russian Foundation for Basic Research (RFBR) and Deutsche Forschungsgemeinschaft (DFG): RFBR grant 98-02-04092, DFG grants 436 RUS 113/439/0 and MA 610/20-1, RFBR grant 00-02-16253, special RFBR grant for Leading Russian Science Schools 00-15-96555; NATO Grant PST.CLG975264, Russian Federal Program "Center of Fundamental Optics and Spectroscopy", and Program of Russian Ministry of Science and Technology "Physics of Solid State Nanostructures".
FIGURES
Fig. 1. Panel a: Experimental setup for the SH interferometric spectroscopy. BS, beam splitter; GG and UG, yellow and UV filters, respectively. Panel b: Raw SH interferograms for different SH energies. Solid curves: The dependences given by Eq.(1). Inset: The spectral dependence of the period $`L`$ of the SH interferograms and its fit using a phenomenological expression for air dispersion. Open circles indicate the periods for the curves at the main panel.
Fig. 2. Spectrum of the SH phase $`\mathrm{\Phi }_{rs}`$ (panel a) and SH intensity $`I_s^{2\omega }`$ (panel b). Solid curves are fits to the data within the model of CP line shapes. Dotted lines are fits with excitonic line shapes for all the resonances.
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# Microlensing of multiply–imaged compact radio sources
## 1 Introduction
Gravitationally lensed compact radio sources have many astrophysical and cosmological applications. The foremost being the determination of a time-delay between the individual lensed images in order to constrain the Hubble parameter (e.g. Refsdal 1964). Considerable progress has been made during the last few years in measuring time-delays, both through optical and radio observations (e.g Kundić et al. 1997; Schechter et al. 1997; Lovell et al. 1998; Biggs et al. 1999; Fassnacht et al. 1999; Koopmans et al. 2000). They also allow a detailed study of the mass distribution of the lens galaxy and sometimes the background source, through a large magnification by the lensing potential. Absorption lines in the spectrum of the background source allow the study of the ISM in the lens galaxy and the HI distribution along the lines of sight to the source.
Temporal changes in the brightnesses or spectra of the lensed images also allow the study of uncorrelated external variability. The most important sources of external variability are scintillation at radio wavelengths and microlensing in all wavelength bands. Differencing the image light curves, taking the proper time delay into account, removes intrinsic source variability and leaves only uncorrelated external variability. These difference light curves thus provide valuable information on the compact objects in the lens galaxy (e.g. Schmidt & Wambsganss 1998) and/or on the intervening ionized medium.
The study of the ionized component of the Galactic interstellar medium (ISM) through scattering of radio emission from pulsars has had a long tradition (e.g. Rickett 1977, 1990). Scattering by the ionized ISM can explain long-term variability at meter wavelengths (e.g. Condon et al. 1979), as well as large-amplitude variability in very compact extra-galactic radio sources (e.g. Rickett, Coles & Bourgois 1984). Low-amplitude variability at shorter wavelengths (about 10 cm), called ‘flickering’, has been observed by Heeschen (1982, 1984) and is probably associated with refractive interstellar scattering of an extended source (e.g. Rickett et al. 1984). Strong intra-day variability of very compact radio sources might result from refractive interstellar scattering as well (e.g. Wagner & Witzel 1995). A power-law model of the plasma-density power spectrum (e.g. Rickett 1977, 1990), combined with some distribution of this plasma in our galaxy (e.g. Taylor & Cordes 1993 \[TC93\]) is able to explain most of the observed dispersion measures and variability in pulsars at low frequencies, as well as the variability of extra-galactic radio sources at both low and high frequencies. However, especially for compact flat-spectrum radio sources it remains exceedingly difficult to separate intrinsic variability from scintillation by the Galactic ionized ISM.
Gravitationally lensed (i.e. multiply-imaged) flat-spectrum compact radio sources could offer a solution to this problem. As mentioned previously, these systems provide two or more lines-of-sight through the Galactic ionized ISM. For typical image separations of a few arcseconds, one is looking through parts of the Galactic ionized ISM separated by a few hundred AU. One can expect the scattering of radio waves to be independent for the different lines-of-sight. Differencing the image light curves, after a correction for the appropriate time delay and flux-density ratio, produces a difference light curve that only shows uncorrelated external variability. This difference light curve can be studied to obtain information on the Galactic ionized ISM independent from intrinsic source variability.
However, uncorrelated external variability of the lensed images might also originate from microlensing in the lens galaxy (e.g. Chang & Refsdal 1979). This offers the additional opportunity to study the properties of compact objects in the lens galaxy, if microlensing variability dominates or can be separated from scintillation. Optical microlensing in the lens galaxy of Q2237+0305 has unambiguously been shown (e.g. Irwin et al. 1989; Corrigan et al. 1991; Ostensen et al. 1996; Lewis et al. 1998; Woźniak et al. 2000). In the radio, several suggestions of microlensing variability have been made (e.g. Stickel et al. 1988; Quirrenbach et al. 1989; Schramm et al. 1993; Romero et al. 1995; Chu et al. 1996; Wagner et al. 1996; Lewis & Williams 1997; Takalo et al. 1998; Quirrenbach et al. 1998; Kraus et al. 1999; Watson et al. 1999). In none of these cases, however, has one really been able to convincingly distinguish between intrinsic and external variability. Claims of external variability in singly-imaged radio sources through microlensing should therefore be regarded with some caution.
In this paper, we report the first unambiguous case of external variability of a radio gravitational lens, CLASS B1600+434 (Jackson et al. 1995; Jaunsen & Hjorth 1997; Koopmans, de Bruyn & Jackson 1998 \[KBJ98\]; Koopmans et al. 2000 \[KBXF00\]). The system consists of two compact flat-spectrum radio images, separated by 1.4 arcsec. The background source, at a redshift of $`z`$=1.59, is lensed by an edge-on disk galaxy at a redshift of $`z`$=0.41 (Fassnacht & Cohen 1998). A time delay of $`47_9^{+12}`$ days (95% statistical confidence) was recently found (KBXF00).
What is furthermore of interest is that this system offers two distinct lines-of-sight through the lens galaxy. Image A only passes mainly through the dark-matter halo around the edge–on lens galaxy, whereas image B passes predominantly through its disk and bulge (Koopmans et al. 1998; Maller et al. 2000; CASTLE Survey, Munoz et al. 1999). This makes image A especially sensitive to microlensing by massive compact objects in the halo and image B to microlensing by stars in the disk and bulge. This might even offer an opportunity to study compact objects in the dark–matter halo around the lens galaxy of B1600+434.
The outline of the paper is as follows. In Section 2, we present the VLA 8.5–GHz data from KBXF00 in a different way, unambiguously showing the presence of external variability. We also present additional WSRT 1.4 and 5–GHz monitoring data of B1600+434. In Section 3, we investigate whether Galactic scintillation can explain the fractional rms variabilities (modulation-indices) and time scales of the short-term variability seen in the VLA 8.5–GHz light curves. Similarly, in Sections 4 and 5 the possibility of microlensing by compact objects in the lens galaxy is studied. In Section 6, we present microlensing simulations of a more complex jet structure and compare the results to B1600+434. In Section 7, we discuss a critical test (i.e. the frequency-dependence of the modulation-index) to discriminate between scintillation and microlensing and compare predictions from the VLA 8.5–GHz light curves with the independent multi-frequency WSRT data. In Section 8 our results and conclusions are summarized.
## 2 Short-term variability in B1600+434–A & B
B1600+434 is a compact ($``$1 mas at 8.5 GHz; KBXF00) radio source, which has varied strongly at 8.5 GHz since its discovery in 1994 (Jackson et al. 1995). Its flux density decreased from 58 (48) mJy in March 1994 to only 29 (24) mJy in August 1995 for image A (B) (KBJ98). From February to October 1998, another decrease from 27 (24) to 23 (19) mJy was found (KBXF00). In June 1999, the flux densities appear to have stabilized to 23 (17) mJy. Strong variability was also observed at 5 GHz, where a total flux density was measured of 34–37 mJy in 1987 (GB87; Becker, White and Edwards 1991). Observations in March 1995 gave 45 (37) mJy for image A (B) (KBJ98), whereas in June 1999 this had reduced to only 23 (18) mJy. At 1.4 GHz, the integrated WSRT flux density of B1600+434 has decreased from 60–65 mJy in April–July 1996 (KBJ98) to about 50 mJy in June 1999.
Because of this strong variability, B1600+434 was observed from February to October 1998 with the VLA at 8.5 GHz, in order to measure a time delay between the two lensed images. In KBXF00, the VLA 8.5–GHz light curves of the two lensed images were presented, showing short-term variability, as well as a long-term decrease in the flux density of both lensed images, which we assume to be intrinsic source variability, similar to that seen in for example 0957+561 (e.g. Haarsma et al. 1997).
In Fig. 1, the normalized VLA 8.5–GHz light curves of both lensed images are shown. The curves were created by dividing the light curves by linear fits (KBXF00). The resulting curves show the fractional variability on short time scales (i.e. shorter than the observing period of about 8 months) with respect to the running mean. The 1–$`\sigma `$ error for each point on the light–curve is 0.7 to 0.8% (KBXF00). We omitted the six strongest outliers from both plots, which show clear systematic problems with the data or calibration (KBXF00). The resulting normalized light curves have modulation indices (i.e. fractional rms variabilities) of 2.8% and 1.6% for images A and B, respectively. We use these values throughout this paper. Lensed images A and B have modulation-indices significantly larger than expected on the basis of the measurement errors only, indicating the clear presence of short-term variability, either intrinsic or external. The structure functions (e.g. Simonetti, Cordes, & Heeschen 1985) of the normalized light curves are shown in Fig. 6 (Sect. 3.2).
### 2.1 Intrinsic or external variability?
To show that most of the short-term variability is of external origin, the linearly-interpolated light curve of image B was subtracted from the light curve of image A, taking a flux-density ratio of 1.212 and a time delay of 47 days into account (KBXF00). The time delay was determined with the minimum-dispersion method from Pelt et al. (1996). Hence, the modulation-index of the normalized difference light curve, shown in Fig. 2, is by definition a lower limit. This is illustrated by the dotted and dashed lines, which show the normalized difference light curves for a time delay of 41 d and 52 d (i.e. the 68% confidence region), respectively. Both curves have larger modulation-indeces, as should be expected.
The normalized difference light curve has a modulation-index of 2.8%. This is significantly larger than the modulation-index of 1.0–1.1% (i.e. the shaded region in Fig. 2), which one would expect on the basis of the measurement errors only. Most of the short-term variability present in the light–curves must therefore be of external origin. A $`\chi ^2`$–value of 377 was determined from the 58 points composing the difference light curve. This is inconsistent with a flat difference light curve at the 14.6–$`\sigma `$ confidence level.
If the short–term variability of the two lensed images is uncorrelated, the expected modulation-index in the normalized difference curve should be $`\sqrt{(2.8^2+1.6^2)}`$$``$3.2%, which is slightly larger than the observed value of 2.8%. It remains hard to assess whether individual features in the light curves might be of intrinsic origin. In this paper we will therefore assume that all short–term variability in both lensed images is of external origin.
### 2.2 WSRT 1.4 & 5–GHz monitoring
Before proceeding with the analysis of the VLA 8.5–GHz light curves, we first present multi-frequency WSRT total-flux-density data of B1600+434 obtained in 1998/9 at 1.4 and 5 GHz. This data will play an essential role in distinguishing between the different physical mechanisms causing the external variability observed in the VLA 8.5–GHz lensed-image light curves, as we will see in Sect. 7.
Starting in August 1998 the WSRT was outfitted with a series of new multi-frequency front-ends (MFFE’s) that can operate at frequencies from 0.3 to 8.5 GHz. When the WSRT observations of B1600+434 were begun in August 1998 only 6 front-ends were available at 1.4 GHz. The available number of telescopes with MFFE’s increased at a rate of about 1 per month until the full array was outfitted in February 1999. Towards the end of 1998 the monitoring was extended to include 5 GHz. In the spring of 1999 we also included 8.5–GHz observations. However, the analysis of the 8.5–GHz data is still encountering some problems and we therefore do not report on the results of the 8.5–GHz observations here.
Each run consisted of two sets of observations (at up to 3 frequencies) on B1600+434 and the nearby reference source B1558+439. The latter is a strong steep spectrum double radio source (0.8–arcsec size) about 40 arcmin north-east of B1600+434. In each run we also observed two primary calibration sources (3C286, 3C343 and CTD93). Although we changed the details of the observing sequence during the year the basic structure did not change.
The resolution of the WSRT is about 12$`\times `$18 arcsec (1.4 GHz) and 3.5$`\times `$5 arcsec (5 GHz). B1600+434 is therefore always unresolved at 1.4 GHz. At 5 GHz the source, however, shows slight hour-angle dependent resolution effects \[the WSRT is an east-west synthesis array, hence the instantaneous synthesized response is a fan beam rotating clockwise on the sky\]. Because the observations were scheduled at random hour angles the resolution effect is therefore variable from session to session. We minimized the magnitude of this effect by determining the flux density using only baselines up to 1300 meter; the residual effect on the flux density is below the thermal noise error.
The amplitude and flux calibration was performed in NEWSTAR using standard procedures starting with a selfcalibration on the primary reference source. The complex telescope gains were then transferred to the target and reference source. The flux densities of B1600+434 and B1558+439 were then determined using a uv-plane fitting algorithm in the program NMODEL. Selfcal phase solutions on the reference source B1558+439 usually showed only very slight decorrelation effects due to slow instrumental/atmospheric phase-drifts. Because these would be very similar for B1600+434 and the reference source we decided not to apply a phase selfcal solution.
The flux density of the reference source B1558+439 was found to show a scatter of about 1–2% around 700 mJy at 1.4 GHz and about 1.0% around 204 mJy at 5 GHz. This scatter is still larger than the noise error on the flux density. We believe these can be attributed to small changes in atmospheric opacity and small instrumental gain drifts (e.g. due to pointing). Normalizing the amplitudes of B1600+434 by those of the reference source should eliminate them. We therefore expect that the final errors on the flux density of B1600+434 are determined by the thermal noise only. Still, to be on the safe side we adopted as a final error on the flux density of B1600+434 the quadrature sum of the thermal noise level and a 1% scale error. This amounts to a typical error of 1.0% and 1.3% at 1.4 and 5 GHz, respectively.
In Fig. 3, we have shown the calibrated WSRT 1.4–GHz (L–band) and 5–GHz (C–band) light curves, together with the total-flux-density VLA 8.5–GHz (X–band) light curve (KBXF00). We note the following properties: (i) All light curves are dominated by systematic trends with decreasing intensity at 8.5 and 1.4 GHz, but increasing at 5 GHz. We believe these changes to be due to intrinsic variations. (ii) If the VLA 8.5–GHz and WSRT 5–GHz long-term intrinsic flux-density variations are correlated, there has been a clear trend-break in the gradient of the light curves around day 300 (Fig. 3). The WSRT 1.4–GHz light curve still shows a similar gradient as the VLA 8.5–GHz light curves in 1998, hence there appears to be a time-lag of at least 200 days between long-term intrinsic source variability at 8.5 and 1.4 GHz. This trend-break is supported by the first results from the 1999 VLA campaign (Koopmans et al. in prep.). (iii) At those epochs where the light curves overlap, the higher-frequency light curves show a larger short-term modulation-index around the long-term linear gradient.
An important statistical property of the light curves is their modulation-index on short time scales (i.e. time scales shorter than the length of the light curves), as function of frequency. To calculate the modulation-indices, we divide the light curves through a linear fit (see Sect. 2) in order to remove most of the presumably intrinsic source variability. The results are normalized light curves, similar to the normalized VLA 8.5–GHz light curves shown in Fig. 1. We only calculate the 1.4 and 5–GHz modulation indices ($`m_{\mathrm{part}}`$) for those epochs, where we have both a 1.4 and 5–GHz WSRT flux-density measurements. The normalized light curves of these epochs are shown in Fig. 4. The resulting modulation-indices are listed in Table 1. Because the modulation-index at 1.4 GHz is very close to the estimated flux-density error, we regard it as an upper limit.
The total flux-density modulation-indices in Table 1 are a combination of the individual modulation-indices of images A and B. If we assume that the ratio of modulation-indices are equal to those at 8.5 GHz (i.e. $`(m_\mathrm{A}/m_\mathrm{B})_\lambda `$$``$2.8/1.6=1.75) for each wavelength, we only make a slight error ($``$15%) compared to the assumption that image B does not vary at all on short time scales. Because we will only use modulation-index ratios (see Sect. 7), which are independent from these assumptions in first order, this is of negligible importance.
### 2.3 Possible origins of the external variability
What can be the origin of the external variability seen in the VLA 8.5–GHZ light curves, why does the modulation-index differ between the two image light curves and what causes the short-term variability in the WSRT 1.4 and 5–GHz light curves? Below we have listed different physical mechanisms which can introduce external variability in the flux density of compact radio sources:
1. Scintillation
* Weak scattering (Sect. 3.1)
* Refractive scattering (Sect. 3.2)
* Diffractive scattering (Sect. 3.3)
2. Microlensing (Sect. 4–5)
In the next three sections, we will investigate in detail whether one or more of these can explain the external variability seen in the VLA 8.5–GHz light curves of B1600+434–A & B. In Sect. 7, we will combine the conclusions from Sections 3–5 with the results from the WSRT 1.4 and 5–GHz observations to further constrain the scintillation and microlensing hypotheses.
## 3 Scintillation
Even though the difference in modulation-index between images A and B seems to require a considerable change in the properties of the Galactic ionized ISM over an angular scale of 1.4 arcsec (Sect. 3.2), we still proceed to investigate whether the short-term variability, superposed on the gradual and presumably intrinsic long-term decrease of the flux density of the lensed images, can be the result of scintillation. We will follow the prescription of Narayan (1992) and its numerical implementation by Walker (1998; \[W98\]) for the Galactic ionized ISM model from TC93. This assumes that the inhomogeneities of the ionized ISM can be described by a Kolmogorov power-law spectrum (e.g. Rickett 1977; Rickett 1990) and that the ionized ISM model from TC93 is approximately valid. Support for the approximate validity of the TC93 model in the direction to B1600+434 is given by the dispersion and scattering measures of nearby pulsars, showing no apparent deviations from this model. Also, no evidence is found in the low-frequency (327–MHz) WENSS catalogue (e.g. Rengelink et al. 1997) for diffuse HII emission or SN remnants that could introduce small-scale perturbations in the Galactic ionized ISM.
Depending on the line-of-sight through the galaxy and the observing frequency, the scattering strength of radio waves, expressed in the scattering strength $`\xi `$=$`r_\mathrm{F}/r_{\mathrm{diff}}`$, can be strong ($`\xi `$$`>`$1) or weak ($`\xi `$$`<`$1), where $`r_\mathrm{F}`$ is the Fresnel scale and $`r_{\mathrm{diff}}`$ is the diffractive scale (e.g. Narayan 1992). The transition between these two regimes occurs near a transition frequency ($`\nu _0`$). For B1600+434 at a Galactic latitude of $`b`$=+48.6, we find $`\nu _0`$=4.2 GHz (W98). At the observing frequency of 8.5 GHz, scattering should therefore be in the weak regime ($`\xi `$$``$0.3). The transverse velocity of the ISM with respect to the line-of-sight to B1600+434 is determined by projecting the velocity vector of the earth’s motion on the sky as function of time. We find a transverse velocity ($`v`$) between 20 and 40 km s<sup>-1</sup>. We will therefore adopt an average value of $`v`$=30 km s<sup>-1</sup> throughout this paper. For lack of better knowledge, we assume any intrinsic transverse motion of the scattering medium to be zero.
### 3.1 Weak scattering
The modulation-index of a point source in the weak scattering regime is (Narayan 1992)
$$m_\mathrm{p}=\left(\frac{\nu _0}{\nu }\right)^{17/12}.$$
(1)
In the simplest case that B1600+434 is a point source smaller than the Fresnel scale of 3.9 $`\mu `$as (W98), we would expect a modulation-index around 35%. This is significantly larger than the observed modulation-indices of 2.8% for image A and 1.6% for image B. Hence, the total angular size of the images must be larger than the Fresnel scale. This does, however, not exclude that part of the source might still be compact.
The variability time scale for a point source is given by (Narayan 1992)
$$t_\mathrm{p}=\frac{3.3}{v_{30}}\times \sqrt{\frac{\nu _0}{\nu }}\text{ h,}$$
(2)
where $`v_{30}`$ is the transverse velocity of the scintillation pattern with respect to the line-of-sight to the source in units of 30 km s<sup>-1</sup>. Inserting the values for the observing frequency ($`\nu `$), the transition frequency ($`\nu _0`$) and $`v_{30}`$=1.0, a time scale around 2 h is found, which is much smaller than the apparent variability time scales of days to weeks seen in a significant fraction of the light curves (Figs 12).
If 5 to 10% of the flux density of the source is contained in a compact region ($``$$`\theta _\mathrm{F}`$), the rms fluctuations decreases to the observed modulation-index of 2–3% for B1600+434–A and B. The variability time scale would still remain $``$2 h. We have observed B1600+434 with the WSRT at 5 GHz during several 12 h periods and find no evidence for short-term variability $``$2% over a $``$12 h time scale (Koopmans et al. in prep.). This excludes the posibility that the longer time-scale variations are purely the result of undersampled light curves. Hence, a simple compact source structure, embedded in a more extended non-scintillating region of emission, can not explain the observed variability. A more extended source ($``$$`\theta _\mathrm{F}`$) is therefore required, if we want to explain the observed modulation-index in terms of scintillation.
In case the source is extended, with a size $`\theta _\mathrm{s}`$$`>`$$`\theta _\mathrm{F}`$, both the modulation-index and variability time-scale change. The modulation-index decreases as follows (Narayan 1992)
$$m=m_\mathrm{p}\times \left(\frac{\theta _\mathrm{F}}{\theta _\mathrm{s}}\right)^{7/6},$$
(3)
whereas the time scale of variability increases as
$$t=t_\mathrm{p}\times \left(\frac{\theta _\mathrm{s}}{\theta _\mathrm{F}}\right).$$
(4)
Combining these two equations, using the transition and observing frequencies for B1600+434, gives the relation
$$m=\left(\frac{0.90^\mathrm{h}}{v_{30}t}\right)^{7/6},$$
(5)
for $`tt_\mathrm{p}2.25^\mathrm{h}/v_{30}`$.
If the lensed source has an angular radius of about 40 $`\mu `$as, the modulation-index reduces to 2–3%, as observed. The variability time scale should then be around 1 day, still significantly smaller than the observed modulation timescale in a major part of the light curves. Although part of the very short-term ($``$ few days) variability could be due to scintillation, many of the long-term variations seen in Figs 12 certainly can not be explained this way.
The relatively nearby extragalactic radio source J1819+387 has a modulation-index $`m`$$``$0.5 and a variability time scale $`t`$$``$0.5–1 hour (Dennett–Thorpe & de Bruyn 2000). Using the relation between $`m`$ and $`t`$ (e.g. Eq. (5)), this translates to a variability time scale less than a day for a modulation-index of a few percent, in agreement with expectations from the Galactic ionized ISM model from TC93. Similarly, variations with a time scale of more than a week (Figs 12) require a source size of 0.3 mas, reducing the modulation-index to 0.2%, which is well below the noise level in the VLA 8.5–GHz light curves.
In Fig. 5, we have summarized the weak and strong scattering regimes, as function of the modulation-index, the variability time scale and the source size. The variability seen in image A (but also image B) is especially hard to explain by weak scattering without either invoking unlikely high values for the equivalent distance of the phase screen ($``$10 kpc) or a persistently low transverse velocity ($``$few km s<sup>-1</sup>). From this, we conclude that weak scattering has great difficulties in accounting for the observed modulation-indices and longer variability time scales ($``$1 day) seen in the VLA 8.5–GHz light curves of the lensed images. However, it remains difficult to determine a reliable variability time scale from the image light curves, partly because of the relatively poor sampling (i.e. every 3.3 days). In any case, sigificant variability on short time scales ($``$12 h) is excluded (see above). For further discussion we refer to Sect. 3.2, where we present the structure functions of the observed variations.
### 3.2 Strong scattering
In the strong scattering regime, we can not use the numerical results derived from the TC93 model, from which one expects B1600+434 to be in the weak-scattering regime at 8.5 GHz. We therefore make direct use of the relation between the scattering measure (SM), the distance to the equivalent phase screen ($`D_{\mathrm{kpc}}`$), the observing frequency ($`\nu _{\mathrm{GHz}}`$) and the scattering strength ($`\xi `$) (e.g. W98)
$$\xi =2.610^3\times \mathrm{SM}^{6/10}D_{\mathrm{kpc}}^{1/2}\nu _{\mathrm{GHz}}^{17/10}.$$
(6)
The scattering strength and the Fresnel scale determine both the modulation-index and variability time scale of a source, given the source size. The Fresnel scale, given by
$$\theta _\mathrm{F}=8/\sqrt{D_{\mathrm{kpc}}\nu _{\mathrm{GHz}}}\mu \text{as},$$
(7)
specifies the angular distance from the source over which there is about one radian phase difference between rays, due to the difference in path length. The scattering measure (e.g. TC93) for an extra-galactic source is defined as
$$\mathrm{SM}=_0^{\mathrm{}}C_\mathrm{n}^2dl,$$
(8)
where $`C_\mathrm{n}^2`$ is the structure constant normalizing the Kolmogorov power-law spectrum of the ionized ISM inhomogeneities (e.g. Cordes, Weisberg & Boriakoff 1985). From now on, we assume that SM has units of kpc m<sup>-20/3</sup> and $`C_\mathrm{n}^2`$ units of m<sup>-20/3</sup>. The distance to the equivalent phase screen (e.g. W98) is defined as
$$D\frac{1}{\mathrm{SM}}_0^{\mathrm{}}C_\mathrm{n}^2sds.$$
(9)
Despite the fact that the difference in modulation-index of the lens images seems to require very different properties of the Galactic ionized ISM on a scale of 1.4 arcsec, we will investigate the two distinct strong scattering regimes, i.e. refractive and diffractive (e.g. Rickett 1990; Narayan 1992), in more detail in the next two sections.
#### 3.2.1 Refractive scintillation
Using Eq. (6) and the scaling laws from Narayan (1992), one finds for a point source in the strong scattering regime that the modulation-index is
$`m_\mathrm{p}`$ $`=`$ $`\xi ^{1/3}`$
$``$ $`7.310^2\mathrm{SM}^{1/5}D_{\mathrm{kpc}}^{1/6}\nu _{\mathrm{GHz}}^{17/30},`$
whereas the variability time scale is
$$t_\mathrm{p}=\frac{\theta _\mathrm{F}D}{v}\xi \frac{3.3}{v_{30}}D_{\mathrm{kpc}}\mathrm{SM}^{6/10}\nu _{\mathrm{GHz}}^{11/5}\text{ yrs}.$$
(11)
We furthermore use $`D_{\mathrm{kpc}}`$=0.5 (TC93), $`v_{30}`$=1.0 and $`\nu _{\mathrm{GHz}}`$=8.5 throughout this section. From Eqs. (10–11) it is immediately obvious that for a point source in the refractive regime, an extremely high value for SM is needed ($`10^5`$ kpc m<sup>-20/3</sup>) to obtain the modulation-index of images A and B. The time scale of variability would be around 15 years. Clearly the point-source approximation is not valid.
For extended sources, the modulation-index and time scale of variability scale as $`(\theta _\mathrm{r}/\theta _\mathrm{s})^{7/6}`$ and $`(\theta _\mathrm{s}/\theta _\mathrm{r})`$, respectively, where $`\theta _\mathrm{s}`$ is the source size and $`\theta _\mathrm{r}`$ the size of the scattering disk (Narayan 1992). At 8.5 GHz, we find
$$\theta _\mathrm{r}=\theta _\mathrm{F}\xi 180\times \mathrm{SM}^{6/10}\mu \mathrm{as},$$
(12)
which is independent of the distance to the equivalent phase screen. If we subsequently use the scaling laws, combined with Eqs. (10–11), we find a relation between the time scale of variability and the modulation-index:
$$m=\xi ^{5/6}\left(\frac{\theta _\mathrm{F}D}{vt}\right)^{7/6},$$
(13)
which is valid only if $`\theta _\mathrm{s}>\theta _\mathrm{r}`$. Inserting the usual numerical values for $`v`$, $`D`$ and $`\nu `$, we find
$$m2.0\mathrm{SM}^{1/2}t_\mathrm{d}^{7/6},$$
(14)
with $`t_\mathrm{d}`$ in units of days. We find that a scattering measure SM$``$$`210^2`$ is needed to explain modulations with a time scale of $``$1 week in image A. From TC93 we find that SM=$`210^4`$ in the direction of B1600+434, corresponding to a time scale of one day. For deep modulations of about 1 month a scattering measure SM$``$0.5 is needed for image A. Both values are larger than can be expected on the basis of the ionized ISM model from TC93.
##### Differences in the scattering measure:
A strong argument against refractive scattering is the large difference between the modulation-indices of images A and B. If this is due to a difference in the scattering measure, it requires $`\mathrm{SM}_\mathrm{A}/\mathrm{SM}_\mathrm{B}`$$``$$`(0.028/0.016)^2`$$``$3.1 (Eqs. (13–14)), which is substantial over an angular scale of only 1.4 arcsec. The same factor is found for the weak-scattering regime.
##### The structure function:
We have also calculated the structure function (SF; Simonetti et al. 1985):
$$\mathrm{D}^{(1)}(\tau )=\frac{1}{2}<(\delta \mathrm{F}(t)\delta \mathrm{F}(t+\tau ))^2>,$$
(15)
following Blandford, Narayan & Romani (1986 \[BNR86\]), who investigated intensity fluctuations (i.e. “flickering”) of extended radio sources, caused by refractive scattering. $`\mathrm{F}(t)`$ is the normalized light curve as shown in Fig. 1. BNR86 take a slightly steeper spectrum of the phase fluctuations with a power-spectrum slope $`\beta `$=4, instead of a Kolmogorov slope of $`\beta `$=11/3.
Fitting the theoretical SF from BNR86 to the observed SFs<sup>1</sup><sup>1</sup>1The structure function D$`{}_{}{}^{(1)}(\tau )`$, calculated from the normalized light–curves, corresponds to the structure function D$`{}_{}{}^{(2)}(\tau )`$, calculated from the observed light curves (KBXF00), because SFs of order $`M`$ remove polynomials of order $`M`$–1 from the time series (e.g. Simonetti et al. 1985). (Fig. 6) gives a scale length of $`L`$$``$0.9 kpc, which corresponds to an equivalent phase screen distance of 0.9$`/\sqrt{\pi }`$0.5 kpc (Eq. (9)). The ‘best-fit’ image sizes are 62 and 108 $`\mu `$as, respectively, for images A and B. The saturation time scales (e.g. BNR86) found from these fits are $`\tau _\mathrm{s}^\mathrm{A}`$=2.5 days and $`\tau _\mathrm{s}^\mathrm{B}`$=4.4 days, even though there are clear variations with longer time scales in the light curves. The presence of variabilty with longer time scales ($``$week) has been confirmed by new multi-frequency VLA observations in 1999/2000 (Koopmans et al. in prep.).
To test the reliability of these saturation time scales, we replaced the normalized flux densities at each epoch in Fig. 1 by Gaussian–distributed values with a 1–$`\sigma `$ scatter equal to the observed modulation-index of the light curve. In Fig. 6 the result is shown, from which it is immediately clear that the light curves are undersampled such that the SFs and the saturation time scales for time lags $`\tau `$$``$4 days become highly unreliable.
##### Scatter-broadening:
The difference between the modulation-indices of images A and B can be explained by a difference in the scattering measure of the Galactic ionized ISM towards both images, as well as by a difference in their respective image sizes (see above). However, image B has a smaller magnification due to the lensing potential and should therefore be smaller than image A. Consequently, image B should show stronger variability than image A, whereas it does not. The only viable alternative to obtain a larger size for image B is through scatter-broadening by the ionized ISM in the lens galaxy.
The expected scattering disk at 8.5 GHz due to the Galactic ionized ISM is $``$1 $`\mu `$as and cannot account for the apparent difference in the image size, derived from the SFs. This requires a scattering disk of $``$90 $`\mu `$as for image B at 8.5 GHz, if image A is not scatter-broadened, implying that SM<sub>B</sub>$``$1 in the lens galaxy. If we take into account that scattering occurs at a frequency of 8.5$`\times (1+z_\mathrm{l})`$$``$12.0 GHz, this implies a scattering disk of 3 mas at 1.7 GHz, which is nearly equal to the very conservative upper limit of $``$4 mas on the image sizes found from 1.7–GHz VLBA observations (Neal Jackson, private communication).
Recent polarization observations by Patnaik et al. (1999) gave rotation measures RM=40 rad m<sup>-2</sup> and has RM=44 rad m<sup>-2</sup> for images A and B, respectively. The difference of 4$`\pm `$4 rad m<sup>-2</sup> is rather low and certainly does not support a high scattering measure in the disk/bulge of the lens galaxy.
Hence, although scatter-broadening cannot be excluded, to fully explain the observed difference between the modulation-indices of images A and B in terms of galactic scintillation, one would require an extremely high scattering measure in the lens galaxy.
#### 3.2.2 Diffractive scintillation
For diffractive scintillation at 8.5 GHz to be at work, one requires both a very high scattering measure and a very small source, neither of which seems plausible. However, to be complete we briefly discuss this possibility.
The modulation-index is unity for a point source,
$$m_\mathrm{p}=1,$$
(16)
much larger than seen in both lensed images. However, for a source larger than the scale on which there are phase changes of about 1 radian ($`\theta _\mathrm{d}`$), the modulation-index becomes $`m=(\theta _\mathrm{d}/\theta _\mathrm{s})`$, where $`\theta _\mathrm{s}`$ is the source size. We find (e.g. W98)
$`\theta _\mathrm{d}`$ $`=`$ $`\theta _\mathrm{F}\xi ^1`$
$``$ $`3.110^3\mathrm{SM}^{6/10}\nu _{\mathrm{GHz}}^{6/5}D_{\mathrm{kpc}}^1\text{ }\mu \text{as},`$
and for the point–source variability time–scale
$$t_\mathrm{d}=t_\mathrm{F}\xi ^1\frac{0.25}{v_{10}}\mathrm{SM}^{6/10}\nu _{\mathrm{GHz}}^{6/5}\text{ min},$$
(18)
where $`t_\mathrm{F}`$=$`\theta _\mathrm{F}Dv^1`$=$`4.010^4v_{30}^1\sqrt{D_{\mathrm{kpc}}/\nu _{\mathrm{GHz}}}`$ sec. The time scale increases by $`(\theta _\mathrm{s}/\theta _\mathrm{d})`$, if the source size is larger than $`\theta _\mathrm{d}`$. Combining the equations above, we find the relation
$$\xi =\frac{t_\mathrm{F}}{mt},$$
(19)
where $`m`$ is the observed modulation-index. This relation is independent of source size, as long as the source is larger than the diffractive scale $`\theta _\mathrm{d}`$. Using this equation, we immediately find that $`\theta _\mathrm{d}>\theta _\mathrm{F}`$ for the deep modulations of $``$1 week, which is only true in the weak scattering regime. Thus, diffractive scattering offers no solution either, which comes as no surprise.
## 4 Radio Microlensing: Theory
Microlensing is unlikely for bright ($``$1 Jy) flat-spectrum radio sources at low frequencies ($``$few GHz), which have typical angular sizes of $``$1 mas, determined by the inverse Compton limit on their surface brightness (e.g. Kellerman & Pauliny–Toth 1969). This angular size is much larger than the typical separation of a few $`\mu `$as between caustics in the magnification pattern, thereby reducing any significant microlensing variability (see Sect. 5).
High-frequency ($``$10 GHz) sources with flux densities less than a few tens of mJy, however, can be as small as several tens of $`\mu `$as. If part of the source is moving with relativistic velocities ($`\beta _{\mathrm{app}}c`$), Doppler boosting allows an even smaller angular size. In those cases, microlensing can start to contribute significantly to the short-term variability seen in these sources (e.g. Gopal-Krishna & Subramanian 1991).
This is in stark contrast to optical microlensing, where the variability time scales are dominated by the transverse velocity ($`v_{\mathrm{trans}}`$) of the galaxy with respect to the line-of-sight to the static source, as in the case of Q2237+0305 (e.g. Wyithe, Webster & Turner 1999). Microlensing time scales between strong caustic crossings in the optical waveband are therefore several orders of magnitude (i.e. $`\beta _{\mathrm{app}}c/v_{\mathrm{trans}}`$$``$$`10^3`$) larger than in the radio waveband. This makes superluminal radio sources the perfect probes to study microlensing by compact objects in strong gravitational lens galaxies, using relatively short ($``$1 year) monitoring campaigns.
### 4.1 Relativistic jet-components
If the lensed source consists of a static core and a synchrotron self-absorbed jet-component, which moves away from the core with a velocity $`\beta _{\mathrm{bulk}}`$=$`v/c`$, the Doppler factor of this jet-component is given by
$$𝒟=[\gamma (1\beta _{\mathrm{bulk}}n)]^1,$$
(20)
where $`\gamma `$=$`(1\beta _{\mathrm{bulk}}^2)^{1/2}`$ and $`n`$ is the direction of the observer (e.g. Blandford & Königl 1979).
The observed flux density of a circular-symmetric radio source with an observed angular radius $`\mathrm{\Delta }\theta `$ is
$$S(\nu )=\frac{2kT_\mathrm{b}\nu ^2}{c^2}\pi \mathrm{\Delta }\theta ^2,$$
(21)
where $`\nu `$ is the observing frequency and $`T_\mathrm{b}`$ the observed brightness temperature of the source. We assume the source has a constant surface brightness. However, due to the Doppler boosting, the apparent brightness temperature of a relativistic jet–component moving towards the observer can be significantly brighter than the inverse Compton limit of about $`T_{\mathrm{b},\mathrm{ic}}`$$``$$`510^{11}`$K (e.g. Kellerman & Pauliny-Toth 1969). The true comoving surface brightness temperature of a flat-spectrum radio source is related to the observed surface brightness temperature ($`T_\mathrm{b}`$) through
$$T_\mathrm{b}=10^{12}\times T_{12}^{\mathrm{b},\mathrm{ic}}\left(\frac{𝒟}{1+z}\right)\mathrm{K},$$
(22)
where $`z`$ is the redshift of the source (e.g. Blandford & Königl 1979). If we substitute Eq. (22) into Eq. (21), we find that the flux density of the relativistically moving jet-component is
$$S_{\mathrm{knot}}=0.23\left[\frac{𝒟T_{12}^{\mathrm{b},\mathrm{ic}}}{1+z}\right]\nu _{10}^2\left(\frac{\mathrm{\Delta }\theta _{\mathrm{knot}}}{\mu \mathrm{as}}\right)^2\text{ mJy},$$
(23)
with $`\nu `$=10$``$$`\nu _{10}\mathrm{GHz}`$. Inverting this equation, we find an approximate angular radius of the jet-component
$$\mathrm{\Delta }\theta _{\mathrm{knot}}=2.1\left[\left(\frac{S_{\mathrm{knot}}}{\text{mJy}}\right)\frac{(1+z)}{𝒟T_{12}^{\mathrm{b},\mathrm{ic}}\nu _{10}^2}\right]^{1/2}\text{ }\mu \text{as}.$$
(24)
Given the observed frequency, the redshift of the jet-component and the Doppler boosting $`𝒟`$, we can subsequently set a limit on the angular radius of the jet-component.
In the case of B1600+434, the redshift of the lensed quasar is $`z`$=1.59 (Fassnacht & Cohen 1998). The observing frequency is $`\nu _{10}`$=0.85 and $`S_{\mathrm{knot}}`$$``$$`S_{8.5}^{\mathrm{tot}}f`$, where $`f`$ is the fraction of the total average source flux density in the relativistic jet-component, $`S_{8.5}^{\mathrm{tot}}`$$``$$`25/\mu `$ mJy (KBXF00) and $`\mu `$ is the average magnification at the image position. From the singular isothermal ellipsoidal (SIE) mass model (Kormann, Schneider & Bartelmann 1994), we find $`\mu _\mathrm{A}`$$``$1.7 ($`\kappa `$=0.2) and $`\mu _\mathrm{B}`$$``$1.3 ($`\kappa `$=0.9). We use $`\mu `$$``$1.5 as a typical value.
After inserting all the known observables into Eq. (24) and adopting $`T_{12}^{\mathrm{b},\mathrm{ic}}`$=0.5, we find an approximate relation between the fraction of the total observed flux–density of B1600+434 contained in the jet–component and its angular size in the source plane
$$\mathrm{\Delta }\theta _{\mathrm{knot}}20\sqrt{\frac{f}{𝒟}}\text{ }\mu \text{as}.$$
(25)
This equation can be used to put constraints on the light-curve fluctuations seen in B1600+434, and decide whether they are the result of microlensing of a single relativistic jet-component.
### 4.2 Microlensing time scales
The typical time scale which one would expect between relatively strong microlensing events is the angular separation between strong caustic crossings divided by the angular velocity of the jet-component in the source plane.
In case the source is not lensed, the apparent velocity (in units of $`c`$) of the jet-component is
$$\beta _{\mathrm{app}}=\frac{n\times (\beta _{\mathrm{bulk}}\times n)}{1\beta _{\mathrm{bulk}}n}=\frac{\beta _{\mathrm{bulk}}\mathrm{sin}(\psi )}{1|\beta _{\mathrm{bulk}}|\mathrm{cos}(\psi )},$$
(26)
where $`\psi `$ is the angle between the jet and the line–of–sight and $`\beta _{\mathrm{bulk}}`$ the bulk velocity of the jet–component (e.g. Blandford & Königl 1979). The apparent angular velocity (in vector notation) of the jet–component becomes
$`{\displaystyle \frac{\mathrm{d}\theta _\mathrm{s}}{\mathrm{d}t}}`$ $`=`$ $`{\displaystyle \frac{\beta _{\mathrm{app}}c}{(1+z_\mathrm{s})D_\mathrm{s}}}`$
$`=`$ $`{\displaystyle \frac{1.2\beta _{\mathrm{app}}}{(1+z_\mathrm{s})}}\left({\displaystyle \frac{D_\mathrm{s}}{\mathrm{Gpc}}}\right)^1{\displaystyle \frac{\mu \mathrm{as}}{\mathrm{week}}},`$
where $`z_\mathrm{s}`$ and $`D_\mathrm{s}`$ are the redshift and the angular diameter distance to the stationary core, respectively. If the jet-component moves with superluminal velocities ($`\beta _{\mathrm{app}}1`$), one expects angular velocities in the order of several tenths of $`\mu \mathrm{as}`$ per week. Because the source is lensed by the foreground galaxy, its observed angular velocity (in the lens plane) becomes
$$\frac{\mathrm{d}\theta _\mathrm{d}}{\mathrm{d}t}=\left[𝒜^1\frac{\mathrm{d}\theta _\mathrm{s}}{\mathrm{d}t}\right]\left(\frac{D_\mathrm{d}}{D_\mathrm{s}}\right),$$
(28)
where
$$𝒜=\left[\begin{array}{cc}1\kappa \gamma _1& \gamma _2\\ \gamma _2& 1\kappa +\gamma _1\end{array}\right]$$
(29)
is the local transformation matrix of the source plane to the lens plane, with $`\kappa `$ and $`\gamma _{1,2}`$ being the local convergence and shear (e.g. Gopal-Krishna & Subramanian 1991; Schneider et al. 1992).
We calculate the source and caustic structure in the source plane, however. Thus, the angular velocity and source structure undergo the inverse transformation of the angular velocity in the source plane (Eq. (28)). The angular velocity that we need to use, is therefore given by Eq. (27). Using the observed redshift $`z`$=1.59 (Fassnacht & Cohen 1998) of B1600+434 the angular velocity in the source plane reduces to
$$\frac{\mathrm{d}\theta _\mathrm{s}}{\mathrm{d}t}=0.34\beta _{\mathrm{app}}\frac{\mu \mathrm{as}}{\mathrm{week}},$$
(30)
where we assume a flat Friedmann–Robertson–Walker universe with $`\mathrm{\Omega }_\mathrm{m}`$=1 and H<sub>0</sub>=65 km s<sup>-1</sup> Mpc<sup>-1</sup>.
### 4.3 Microlensing modulation-indices
The normalized modulation-index ($`m_\mu `$) of a superluminal jet component, due to microlensing, is
$$m_\mu \sqrt{\frac{\mu ^2}{\mu ^2}1},$$
(31)
where $`\mu `$=$`\mu (t)`$ is the microlensing light curve of the jet component. We determine $`m_\mu `$ as function of the angular size of the source by averaging over randomly-oriented simulated microlensing light curves and find that $`m_\mu `$ can be well approximated by the analytical function
$$m_\mu (\mathrm{\Delta }\theta _{\mathrm{knot}})=\frac{C_\mu }{1+\left(\frac{\mathrm{\Delta }\theta _{\mathrm{knot}}}{\mathrm{\Delta }\theta _\mathrm{b}}\right)}$$
(32)
for $`\mathrm{\Delta }\theta _{\mathrm{knot}}`$$``$20$`\mu `$as, where $`\mathrm{\Delta }\theta _\mathrm{b}`$ is the turnover size of the source after which the modulation-index $`m_\mu `$ decreases linearly with source size and $`C_\mu `$ is the asymptotic modulation-index for a source with $`\mathrm{\Delta }\theta _{\mathrm{knot}}`$$``$0. We fit this function to the numerical results to obtain both $`C_\mu `$ and $`\mathrm{\Delta }\theta _\mathrm{b}`$.
We subsequently combine Eqs. (25) and (32) with the fact that the observed modulation-index in the lensed images is $`m_\mu ^{\mathrm{obs}}=fm_\mu `$ and find
$$m_\mu ^{\mathrm{obs}}(f,𝒟)\frac{C_\mu f}{1+20\mathrm{\Delta }\theta _\mathrm{b}^1\sqrt{f/𝒟}},$$
(33)
where $`\mathrm{\Delta }\theta _\mathrm{b}`$ is in units of $`\mu `$as.
Many jets consist of more than a single jet-component. If we assume that (i) the jet consists of $`N`$ similar jet-components, each containing a fraction $`f_1`$=$`f/N`$ of the total flux–density and (ii) the magnification curves of the jet-components to be uncorrelated, we expect the modulation-index of the combined jet-components to decrease roughly as $``$$`1/\sqrt{N}`$. Hence, we find that $`m_\mu ^{\mathrm{obs}}(N)`$$``$$`fm_\mu (f_1)/\sqrt{N}`$, or
$$m_\mu ^{\mathrm{obs}}(f,𝒟)\frac{C_\mu f/\sqrt{N}}{1+20\mathrm{\Delta }\theta _\mathrm{b}^1\sqrt{f/(𝒟N)}},$$
(34)
where $`f`$=$`Nf_1`$$``$1. Hence, multiple jet-components will in general decrease the modulation-index. However, if the individual jet-components are very compact – i.e. are much smaller than the typical separation between strong caustics –, the microlensing variability will be dominated by single caustic crossings, creating strong isolated peaks in the light curve. We then expect Eq. (34) to break down, such that the factor $`1/\sqrt{N}`$ in the numerator can be removed.
In the case of scintillation, $`N`$ compact jet-components ($`\mathrm{\Delta }\theta _{\mathrm{knot}}`$$``$$`\theta _\mathrm{F}`$) always vary independently, such that the observed modulation-index roughly decreases as $`1/\sqrt{N}`$. However, in the case of microlensing, the modulation-index, caused by the same compact jet-components moving over a magnification pattern, can be independent from $`N`$ or even increase as $`\sqrt{N}`$.
### 4.4 Source constraints
At this point, we give a qualitative recipe to obtain constraints on the jet-component parameters from the observed light curves.
1. First, the observed modulation-indices of both lensed images can be used to solve for the fraction of flux density in the jet-component ($`f`$), as well as its Doppler factor ($`𝒟`$), by comparing them to those found from numerical simulations for different mass functions (MFs), which fix both $`C_\mu `$ and $`\mathrm{\Delta }\theta _\mathrm{b}`$). One obtains a set of two equations (i.e. Eq. (33)) with two constraints ($`m_\mu ^\mathrm{A}`$ and $`m_\mu ^\mathrm{B}`$) and two unknown variables ($`𝒟`$ and $`f`$), which in some cases can be uniquely solved for. The combinations of MFs, which do not give a consistent solution, can be excluded, thereby putting constraints on the allowed MFs in the line-of-sight towards the lensed images.
2. Second, combining the typical observed variability time scale between strong microlensing events and the angular separation of these from the numerical simulations, one can obtain a value for $`\beta _{\mathrm{app}}`$ (Eq. (27)). Combining Eqs. (20) and (26), one then also readily solves for both $`\psi `$ and $`\beta _{\mathrm{bulk}}`$.
Thus, given (i) a mass model of the lens galaxy found from (macroscopic) lens modeling, (ii) some plausible range of MFs near the lensed images and (iii) the observed modulation-indices and variability time scales in the light curves of the lensed images, one can in principle solve for several parameters of the simple jet/jet–component structure: $`𝒟`$, $`f`$, $`\psi `$, $`\beta _{\mathrm{app}}`$, $`\beta _{\mathrm{bulk}}`$ and $`\mathrm{\Delta }\theta `$. However, one should keep in mind that some of the parameters might be degenerate and that we also do not know intrinsic brightness temperature of the components.
## 5 Radio Microlensing: Results
In this section, we examine the observed variability in the light curves of B1600+434 A and B in terms of microlensing and, following the procedure delineated in Sect. 4.4, derive constraints on the properties of the jet-component, as well as on the MFs near both lensed images. In Sect. 7, we will use these results to determine the microlensing modulation-index as function of frequency and compare this to the results from independent WSRT 1.4 and 5–GHz monitoring data of B1600+434 (Sect. 2). An overview of the combined microlensing/scintillation situation in B1600+434 is given in Fig. 8, which might act as guide to the overall situation.
From now on, we will (i) use a source structure consisting of a static core plus a single relativistically moving spherically-symmetric jet-component, (ii) assume that the core and jet parameters do not change appreciably over the time-span of the observations, (iii) assume that all short-term variability in the image light curves (Fig. 1) is dominated by microlensing and that scatter-broadening is negligible. The static core does not contribute to the short-term variability, because its velocity with respect to the magnification pattern of the lens is much smaller than that of the jet-component (i.e. $`v_{\mathrm{core}}`$$``$$`v_{\mathrm{jet}}`$). At this moment, we feel that a more detailed model is not warranted. One should, however, keep in mind that the conclusions drawn below depend on these assumptions !
### 5.1 Numerical simulations
From KBJ98, we find for B1600+434 that the local convergence and shear are close to $`\kappa `$=$`\gamma `$=0.2 for image A, using a singular isothermal ellipsoidal (SIE) mass distribution for the lens galaxy. Similarly for image B, we find $`\kappa `$=$`\gamma `$=0.9. Using these input parameters, we simulate the magnification pattern of a $`50\eta _\mathrm{e}\times 50\eta _\mathrm{e}`$ field for different MFs, where $`\eta _\mathrm{e}`$ is the Einstein radius of a 1–M star projected on the source plane. We use the microlensing code developed by Wambsganss (1999). Using the lens redshift $`z_\mathrm{l}=0.41`$, the source redshift $`z_\mathrm{s}`$=1.59 and H<sub>0</sub>=65 km s<sup>-1</sup> Mpc<sup>-1</sup> in a flat ($`\mathrm{\Omega }_\mathrm{m}`$=1) FRW universe, we find that $`\eta _\mathrm{e}`$=$`2.1\mu `$as. The magnification pattern is calculated on a grid of 1000$`\times `$1000 pixels. Each pixel has a size of 0.107 $`\mu `$as by 0.107 $`\mu `$as.
We simulate 100 randomly-oriented light curves on this grid, for a range of jet-component sizes ($`\mathrm{\Delta }\theta _{\mathrm{knot}}`$=0.125, 0.25,…, 16.0 $`\mu `$as) and MFs for the compact objects (see Sect. 5.2). Each simulated light curve is 54 $`\mu `$as long. For each step (5 pixels) on the light curve, we calculate the magnification of a circular-symmetric jet-component with a constant surface-brightness and a radius $`\mathrm{\Delta }\theta `$, by convolving the magnification pattern with its surface-brightness distribution (e.g. Wambsganss 1999).
We also assume that the surface density near image B is dominated by stellar objects in disk and bulge of the lens galaxy, even though the halo does contribute to the line-of-sight surface density. If we assume that the halo surface density near image B is equal to that near image A ($`\kappa `$=0.2), the density of the caustics in the resulting magnification pattern remains completely dominated by the significantly higher number density of compact objects in the disk and bulge ($`\kappa `$=0.7). We have tested this by modifying the microlensing code to allow multiple mass functions. The rms variabilities calculated from these modified magnification patterns are the same within a few percent from those without the halo contribution (for the range of mass functions and surface densities that we used in this paper), which justifies this simplification.
### 5.2 The mass function of compact objects in the lens galaxy of B1600+434
We use a range of different MFs, subdivided in two classes: (i) power-law MFs and (ii) single-mass MFs.
#### 5.2.1 The power-law MF
In the solar neighborhood it appears that the stellar MF can be represented by a single or segmented power-law of the form $`\mathrm{d}N(m)/\mathrm{d}MM^\alpha `$ (e.g. Salpeter 1955; Miller & Scalo 1979). Recent observations towards the Galactic bulge (e.g. Holtzman et al. 1998) suggest a break in the MF around 0.5–0.7 M, with a shallower slope at lower masses. A lower-mass cutoff is not well constrained, although the break in the MF suggests it might lie around a few tenths of a solar mass. Evolution of the MF from $`z`$=0.4 to $`z`$=0.0 affects the upper mass cutoff only and has no significant influence on microlensing, which is dominated by the mass concentrated around the lower mass cutoff.
1. First, we simulate the magnification pattern, using a power-law MF with $`\alpha =2.35`$ (Salpeter 1955) and a mass range between 0.01–1.0 M.
2. Second, we also investigate the slopes $``$2.85 and $``$1.85 for image B – going through the bulge/disk – assuming again a mass range of 0.01–1 M.
3. Third, we use an MF mass range of 0.1–10 M near image B, with a slope of $``$2.35. This gives an average stellar mass of about 0.3 M, more in line with observations of the bulge of our Galaxy (e.g. Holtzman et al. 1998).
We use the power-law MFs for image B only, even though it is clearly a rough approximation of the true MF. A similar power-law MF for image A, passing through the halo, seems unlikely, especially if the halo consists of stellar remnants (e.g. Timmes, Woosley & Weaver 1996).
The results from these simulations – $`m_\mu `$ as function of $`\mathrm{\Delta }\theta _{\mathrm{knot}}`$ – are fitted by Eq. (32). The values for $`C_\mu `$ and $`\mathrm{\Delta }\theta _\mathrm{b}`$ are listed in Table 2 (models BP1–4).
#### 5.2.2 The single-mass MF
In steep ($`\alpha `$$`<`$$`1`$) MFs most of the mass is concentrated close to the lower-mass cutoff in the MF. It therefore seems appropriate to approximate the MF by a delta-function. We simulated single-mass MFs for both image A and B, for 0.1, 0.2, 0.4, 0.6, 0.8, 1.0, 1.5 M and additionally for 2.5 and 5.0–M for image A. The results are listed in Table 2.
From these simulation, we notice several things: First, $`C_\mu `$ is almost independent from the average mass of the compact objects, for a given surface density. Second, there appears to be a strong correlation between the average mass of the compact objects and the turnover angle ($`\mathrm{\Delta }\theta _\mathrm{b}`$), which means that for a given surface density, shear and source size, the modulation-index is larger for a higher average mass of the compact objects. The results from our simulations are consistent with the results in Deguchi & Watson (1987) and Refsdal & Stabell (1991, 1997).
### 5.3 Microlensing in B1600+434?
As we can see from Table 2, the modulation-index in image A can be explained by a relatively small jet-component with a moderate boosting factor. However, one might expect image B to show similar, if not stronger variability for a similar MF. This is not the case, however. Before we examine this in terms of a different mass function in the disk/bulge and halo of the lens galaxy, we first explore two alternative explanations:
* What if the dimensionless surface density near image B is near unity (i.e. $`\kappa `$$``$1)? In that case, the magnification pattern can become very dense and suppress the modulation-index (e.g. Deguchi & Watson 1987; Refsdal & Stabell 1991, 1997). For B1600+434 we used $`\kappa _\mathrm{B}`$$``$$`\gamma _\mathrm{B}`$ (Sect. 5.1), in which case the density of the caustics remains relatively small, in contrast with the case where $`\kappa `$$``$1 and $`\gamma `$$``$0 (e.g. Refsdal & Stabell 1997). This is supported by simulations with $`\kappa =\gamma =0.999`$ near image B, which give nearly the same results as for $`\kappa =\gamma =0.9`$.
* In Sect. 3.2 we showed that scatter-broadening of image B can suppress scintillation caused by the Galactic ionized ISM. Similarly, is can suppress variability due to microlensing. An overview of the situation is given in Fig. 8. Image A, however, is seen through the galaxy halo at about 4 $`h^1`$ kpc above the disk (KBJ98). It is therefore extremely unlikely to pass through a region with a high scattering measure (Sect. 3.2). In the remainder of the paper, however, we assume that both images A and B are not affected by scatter-broadening (situation 2 in Fig. 8).
We will now explore the only other plausible solution; a very different MF of compact objects near images A and B.
#### 5.3.1 Limits on the MF and source structure
In Table 3, we have listed all combinations between the simulated MFs for images A and B (Table 2) that reproduce the observed short-term modulation-indices of 2.8% (A) and 1.6% (B) for a consistent set of parameters $`f`$ and $`𝒟`$ (Sect. 4.4). Several examples of simulated light curves are shown in Fig. 7.
##### Constraints on the MFs in B1600+434:
From Table 3 one finds that a significantly higher average mass of compact objects in the halo is needed than in the bulge/disk to explain the modulation indeces of both images. No consistent solutions are found for an average mass of compact objects in the halo $`<`$1 M for the range of MFs that we investigated. If we furthermore put a conservative upper limit of $``$0.5 M on the average mass of compact objects in the bulge/disk of B1600+434 – which lies around the break in the Galactic bulge MF (e.g. Holtzman et al. 1998) – only MFs BS1–3 and BP1–4 (Table 2) remain viable MFs for the bulge/disk of B1600+434.
##### Constraints on $`\mathrm{\Delta }\theta `$, $`f`$ and $`𝒟`$:
Using the MFs assumed viable above, we find from Table 3: 0.05$``$$`f`$$``$0.11 and 1.1$``$$`𝒟`$$``$4.0. Using Eq. (25) and the values of $`f`$ and $`𝒟`$ listed in Table 3, the jet-component size then lies between 2$``$$`\mathrm{\Delta }\theta _{\mathrm{knot}}`$$``$5 $`\mu `$as.
##### Constraints on $`\beta _{\mathrm{app}}`$:
To estimate a time scale for strong microlensing variability, we calculate the average power spectrum of the 100 light curves for each MF, for the source size of 2 and 4 $`\mu `$as. The power-spectrum is typically relatively flat at low frequencies, smearing out the long-period modes in the light–curves. At higher frequencies the power drops rapidly. We therefore expect the strongest Fourier modes in the light curves to lie around the turn-over frequency, where the power drops to about 50%. Consequently, we define the typical angular scale of variability ($`\theta _\mathrm{t}`$) to correspond with the half-power frequency in the power-spectrum. In Table 4, we listed the results for those MFs that give a consistent solution (Table 3).
If we now take a separation of $``$2 weeks as indicative for the separation of strong modulations in light curve of image A (see days 80–140 in Figs 12), we find an angular velocity of the jet-component in the source plane between 3 and 9 $`\mu `$as/week. Using Eq. (30), we then derive 9$``$$`\beta _{\mathrm{app}}`$$``$26. This range strongly dependents on the local structure of the magnification pattern, which can differ strongly from place to place (e.g. Wambsganss 1990).
##### Constraints on $`\beta _{\mathrm{bulk}}`$ and $`\psi `$:
Using the constraints of the Doppler factor ($`𝒟`$) and angular velocity of the jet-component ($`\beta _{\mathrm{app}}`$), one also obtains constraints on the angle between the jet–component direction with respect to the line-of-sight to the observer and the bulk velocity of the jet-component ($`\beta _{\mathrm{bulk}}`$). From Fig. 9, we subsequently find: $`4^{}`$$``$$`\psi `$$``$$`13^{}`$ and $`\beta _{\mathrm{bulk}}`$$``$$`0.995`$, for the allowed ranges of these parameters.
Thus, several combinations of MFs for images A and B (Table 3) give solutions that reproduce the observed modulation-indices of both images for a consistent, although not unique, set of jet-component parameters. The derived constraints on the jet-component, however, do agree with observations of confirmed superluminal sources (e.g. Vermeulen & Cohen 1994).
#### 5.3.2 Microlensing by compact halo objects
It appears we have found a lower limit of $``$1 M on the mass of compact objects in the halo around the lens galaxy, under the assumptions that all variability we see is due to microlensing, the jet is dominated by a single component and there is no scatter-broadening. Let us now explore the implications of this in more detail, first mentioning several potential problems.
* Could microlensing be due to a globular cluster (GC) in the halo of the lens galaxy? It is easy to show that the probability of seeing a lensed image through a GC surface density $`>\kappa _{\mathrm{GC}}`$ is
$$P(>\kappa _{\mathrm{GC}})\frac{N}{4\kappa _{\mathrm{GC}}^2}\left(\frac{\sigma _{\mathrm{GC}}}{\sigma _{\mathrm{LG}}}\right)^4.$$
(35)
We take a population of $`N`$$``$150 GCs inside the Einstein radius of the lens galaxy (LG), with a velocity dispersion of $`\sigma _{\mathrm{GC}}`$=7 km s<sup>-1</sup>. These values are typical of those found for our galaxy (e.g. Binney & Merrifield 1999). Using $`\sigma _{\mathrm{LG}}`$$``$200 km s<sup>-1</sup> (KBJ98), $`\kappa _{\mathrm{GC}}`$=$`\kappa _\mathrm{A}`$=0.2, we then find a probability $`P(>\kappa _{\mathrm{GC}})`$$``$$`1.710^3`$ that the microlensing optical depth ($`\tau _{\mathrm{ml}}`$$``$$`\kappa _{\mathrm{ml}}`$) of the GC exceeds that of the dark matter halo. The probability that $`\kappa _{\mathrm{GC}}`$ exceeds $`\kappa _\mathrm{B}`$=0.9, thereby causing similar or larger microlensing variability, is $`710^5`$. Hence, it is very unlikely that a GC in the line-of-sight to lensed image A could enhance the microlensing optical depth significantly.
* What is the influence of binary systems on the mass limit of a compact object in the dark matter halo? We know that a large fraction of stars in the bulge is locked up in binary systems (e.g Holtzman et al. 1998). In the case of high microlensing optical depths, one can consider a binary as a single microlensing object, with a mass equal to the sum of the individual masses.
From Table 2 we see that a higher mass of compact objects gives a higher microlensing modulation-index for a given source size. Hence, if we know the typical stellar mass of objects in the bulge/disk and assume they are all single objects, not in a binary system, we underestimate its modulation-index. One has to take this effect into account.
Having dealt with these possible complications, let us explore the lower mass limit of the compact objects in more detail. If we assume that (i) all compact object in the bulge/disk are not in binary systems and (ii) all compact objects are in binary systems and (iii) use the lowest average mass of objects in the bulge/disk (Tables 2–3), we find a very conservative lower limit of 0.5 M on the mass of individual compact objects in the halo. In the more realistic case where most of the stars in the bulge/disk are in binaries and most compact objects are probably not, a lower limit of 1.5 M is found, assuming that the average bulge/disk stellar mass in the halo is $``$0.1 M. If the bulge/disk stars have average masses somewhere between 0.1 and 0.3 M and are foremost in binaries, the lower limit increases to $``$2.5 M.
As in the case of scintillation, scatter-broadening of image B also suppresses microlensing (Sect. 5.3). If this happens, one would underestimate the true microlensing modulation-index of image B. This would give one more freedom to decrease $`f`$ and/or increase $`𝒟`$ for the microlensed jet-component, thereby changing the required average mass of compact objects in the halo. It would, however, never eliminate the need for them.
## 6 Microlensing of a realistic jet structure
So far, we have only investigated a very simple model of the source structure, consisting of a core plus a single jet-component. To see how a real source behaves, when microlensed by similar MFs as in the lens galaxy of B1600+434, we simulated light curves of a more complex jet structure.
### 6.1 The jet structure of 3C120
We used the inner jet-structure seen in 3C120 (Gómez et al. 1998; Gómez, Marscher & Alberdi 1999) as a ‘template’, because it is one of the best-studied nearby jet structures. The inner jet consists of a core plus at least 15 distinguishable jet-components (Gómez et al. 1999), which move superluminally, with velocities typically around 4.5$`c`$ (H<sub>0</sub>=65 km s<sup>-1</sup> Mpc<sup>-1</sup>; Gómez et al. 1998). The jet was observed at 22 GHz. The total flux density of 3C120 at this frequency is 5.7$`\pm `$0.9 Jy (O’Dell et al. 1978), whereas the fitted inner-jet-components contain about 2.5 Jy (Gómez et al. 1999). We therefore assume that the inner jet contains about 40% of the total flux density of 3C120.
Note that 22 GHz corresponds to approximately X–band (8.5 GHz) observations, when placed at the source redshift of B1600+434 ($`z`$=1.59). We thus scale the size of the jet structure down by the ratio of angular diameters distances for 3C120 and the lensed source of B1600+434 and its flux density by the ratio of luminosity distances squared.
At a redshift of $`z`$=1.59, 3C120 would be a 1.8 mJy source. To obtain a source of about 25 mJy total , as observed for B1600+434, we scale each jet component’s size by a factor $`\alpha _\mathrm{r}`$=3.7 in radius and assume that about 60% of its flux density is contained in an extended structure, which is not sensitive to microlensing. We assume their surface brightness temperature to be conserved and associate the radius of the source with half the FWHM of the component size determined by Gómez et al. (1999). The resulting jet structure has a flux density of about 25 mJy at 8.5 GHz and an angular size of the inner jet less than 1 mas. If we only use the inner jet and not assume the additional 60% of extended emission, one should scale the inner-jet-components by a factor of $`\alpha _\mathrm{r}`$=5.7 to obtain a total flux density of about 25 mJy.
### 6.2 Microlensed light curves of 3C120
The jet structure is randomly placed on the magnification pattern. We recalculate the jet structure and the resulting normalized light curves at epochs separated by 3.3 days, which is the average sampling of the light curves of B1600+434. For image A, we use the MF AS7, which corresponds to a halo filled with stellar remnants, such as black holes and neutron stars. For image B, we use the MF BS2.
We calculate light curves with a total time span of 35 weeks, corresponding to the length of the observed VLA light curves of B1600+434. We subsequently scale the light curves by a factor 0.4 ($`\alpha _\mathrm{r}`$=3.7 for 40% of the flux density from the inner jet) or 1.0 ($`\alpha _\mathrm{r}`$=5.7 for 100% of the flux density from the inner-jet). We repeat these simulations using an apparent velocity three times larger ($`\alpha _\mathrm{v}`$=3).
In Fig. 10, one simulated light curve is shown for the images A and B, for each size scale factor ($`\alpha _\mathrm{r}`$) and velocity scale factor ($`\alpha _\mathrm{v}`$). The modulation-index and estimated variability time scale between strong microlensing events are listed in Table 6.
### 6.3 A comparison with B1600+434
Not only does the modulation-index correspond well with that seen for B1600+434–A and B, also the time scale of variability is in the order of several weeks to months (depending on the choice of $`\alpha _\mathrm{v}`$; Table 6). The similarity between some of the simulated light curves and those observed for B1600+434 is remarkable, knowing that we have not resorted to extreme assumptions.
We therefore regard these simulations as ‘proof of principle’, showing that microlensing of multiply-imaged compact flat-spectrum radio sources, of which more and more are being discovered – for example in the CLASS/JVAS survey (e.g Browne et al. 1997; Myers et al. 2000) – can be a very common occurrence, enabling us to study both the structure of these high–$`z`$ radio sources, as well as the MF of compact objects in the intermediate–$`z`$ lens galaxies.
## 7 Microlensing versus Scintillation
We have seen many arguments in Sections 3–5, for and/or against scintillation and microlensing as the cause of variability in compact flat-spectrum radio sources. Both can in principle explain the observed modulation-index and variability time scales in the VLA 8.5–GHz light curves in the individual lensed images of B1600+434, although it remains very difficult to explain the longer ($``$1 day) variations in the light curves or the difference in modulation-index between the two lensed images in terms of scattering. In the case of microlensing, one would expect to see scintillation at some level as well, possibly complicating a straightforward analysis. How can we then separate these mechanisms as the dominant cause of variability?
Scintillation and microlensing have different dependencies on frequency. Although microlensing is achromatic, the frequency dependence of the source structure predicts a clear dependence of the microlensing variability as function of frequency. For flat-spectrum synchrotron self-absorbed source sources, the source size is inverse proportional to frequency (Eq. (24)). Thus the modulation-index decreases with decreasing frequency (Eq. (32)). In the case of weak and strong refractive scattering (i.e. flickering), however, the modulation-index usually increases with decreasing frequency (e.g. BNR86; Narayan 1992; Rickett et al. 1995). The key to testing whether the observed short-term variability in gravitationally lensed flat-spectrum radio sources is partly or fully dominated by microlensing or scintillation is therefore their strong opposite dependence on frequency.
In Fig. 11, we have plotted the dependence of the modulation index in the case of weak and strong refractive scintillation versus that of microlensing. We assume that the source or jet-component size scales as $`\nu ^1`$ and that scatter-broadening is negligible. All curves are normalized at $`m_{\mathrm{part}}`$=3.7% at 6 cm, as measured with the WSRT in 1999 (Table 1). We determine the modulation-index from scintillation following Rickett et al. (1995) and that from microlensing using Eq.(32). In the case of microlensing, we use the maximum range of the turn-over scale $`\mathrm{\Delta }\theta _\mathrm{b}`$=0.9–4.5 $`\mu `$as and the jet-component size at 8.5 GHz $`\mathrm{\Delta }\theta _{\mathrm{knot}}`$=2–5 $`\mu `$as, found from Tables 2 and 5. We furthermore assume that the short-term variability is dominated by image A, as observed in the VLA 8.5–GHz light curves. The resulting curves show a clear opposite trend as function of wavelength and therefore act as a strong discriminator between microlensing and scintillation. The constraints on the microlensing curves were determined from the VLA 8.5–GHz light curves only and therefore independent from the WSRT 1.4 and 5–GHz modulation-indices.
From the WSRT modulation-indices ($`m_{\mathrm{part}}`$) at 1.4 and 5 GHz (Table 1), one finds $`m_{21}/m_6`$$``$0.31, as indicated by the open circle in Fig. 11. Although this result is based on two frequencies only, it clearly agrees much better with the predictions from microlensing and not that from scintillation! The latter would require a $``$8 times larger value for $`m_{21}`$ (i.e about 9%). We do not plot the VLA 8.5–GHz modulation-index, because it was determined from a different epoch. The modulation-index from microlensing and scintillation might change as function of time, whereas the ratio of modulation-indices, measured simultaneously, is less likely to change.
Thus, if all short-term variability seen in the light curves of B1600+434 is external, it follows the predictions from microlensing. Moreover, if the short-term variability seen in the WSRT 1.4 and 5–GHz light curves is intrinsic, this would be very hard to reconcile with the fact that in 1998 almost all of the VLA 8.5–GHz short-term variability was shown (see Sect. 2) to be external. The most logical conclusion from all this is that the short-term variability at 1.4, 5 and 8.5 GHz is dominated by microlensing. This explains both the modulation-indices as function of frequency at 1.4 and 5 GHz and the longer variability time scale in the VLA 8.5–GHz light curves. In the case of scintillation, one would require either different sizes of the lensed images or a very different ionized ISM towards the images and also a different time scale and frequency-dependence of scattering from that expected from a Kolmogorov spectrum of inhomogeneities of the ionized ISM. All evidence thus far is therefore only consistent with microlensing as the dominant cause of the observed short-term variability.
## 8 Summary & Conclusions
We have shown unambiguous evidence of external variability in the CLASS gravitational lens B1600+434. The difference between the 8.5–GHz VLA light curves of the two lensed images shows external variability at the 14.6–$`\sigma `$ confidence level. The modulation indices of the short-term variability are 2.8% for image A and 1.6% for image B. The difference light curve has an rms scatter of 2.8%, indicating that the short-term variability in both light curves is mostly of external origin (Sect. 2).
We have investigated two plausible sources of this external variability: (i) scattering by the ionized component of the Galactic interstellar medium (ISM) and (ii) microlensing by massive compact objects in the bulge/disk and halo of the lens galaxy.
Based on the ‘standard’ theory of scintillation (e.g. Narayan 1992; Rickett et al. 1995) there should be a considerable increase in the modulation-index with wavelength (Sections 3 and 7). From simultaneous WSRT 1.4 and 5–GHz observations we find, however, that $`m_{21}`$=1.2% and $`m_6`$=3.7% (Table 1), which is a considerable decrease. Scintillation theory predicts $`m_{21}`$=9.0% for $`m_6`$=3.7% (Sect. 7). If the 1.4 and 5–GHz short-term variability is intrinsic, it is hard to reconcile with the fact that in 1998 the VLA 8.5–GHz light curves were dominated by external variability during the full eight months of monitoring (Sect. 2), although it can not be fully excluded yet. Moreover, from microlensing simulations, we expect that $`m_{21}`$=1.2–2.4% if $`m_6=3.7\%`$ (Fig. 11), based on constraints on the source structure and mass function of compact objects found from the VLA 8.5–GHz light curves (Sections 4, 5 and 7). This range agrees remarkably well with the observed modulation index $`m_{\mathrm{part}}`$=1.2% at 21 cm.
Supplementary to this argument, the difference in modulation-index between the lensed images would, in the case of scintillation, argue for either a very different Galactic ionized ISM (SM$`{}_{\mathrm{A}}{}^{}/`$SM$`{}_{\mathrm{B}}{}^{}`$3.1; Section 3.1–2) towards the lensed images or a different image size ($`\mathrm{\Delta }\theta _\mathrm{B}/\mathrm{\Delta }\theta _\mathrm{A}`$$``$1.75; Sect. 3.2), although the latter might result from scatter-broadening. Furthermore, the longer variability time scales at 8.5 GHz ($``$1 day; Figs 12) are also difficult to explain in terms of scintillation, as well as the absence of variability with short time scales in several 12 h WSRT observations at 5 GHz (Koopmans et al. in prep.).
However, the strongest argument against scintillation remains the dominant presence of short-term external variability at 8.5 GHz in 1998, combined with the fact that in 1999 significant short-term variability is seen at 5 GHz, but almost none at 1.4 GHz.
Under the microlensing hypothesis, we find a consistent, although not unique set of jet-component parameters. A core plus a single-jet-component with a size of 2–5 $`\mu `$as, containing 5–11% of the flux density and moving superluminally with 9$``$$`\beta _{\mathrm{app}}`$$``$26, can explain the modulation-index and variability time scale in both lensed images (Sections 4–5). For image A we find a significantly higher average mass of compact objects ($``$0.5 M), compared with those near image B. A much lower mass of compact object would result in a finer magnification pattern and thus in less variability. If image B is scatter-broadened, its microlensing modulation-index is reduced, which might change the lower-limit on the compact object mass.
If one, based on the evidence gathered thus far, accepts that the 1.4, 5 and 8.5–GHz short-term variability in B1600+434–A and B is dominated by microlensing, the profound consequence is that the dark-matter halo at $``$6 kpc above the plane of the disk-galaxy lens in B1600+434 is partly filled with massive compact objects. New WSRT, VLA and VLBI multi-frequency data is being obtained at the moment, which combined with a more comprehensive statistical analysis should provide us with refined constraints on the mass function of compact objects and the source structure (Koopmans et al. in prep.).
###### Acknowledgements.
The authors would especially like to thank Joachim Wambsganss for making available and giving support in using his microlensing code. They thank Frank Briggs for providing useful supermongo code, and Joachim Wambsganss, Jane Dennett-Thorpe, Penny Sackett, Jean-Pierre Macquart and Roger Blandford for useful discussions. They also thank the referee, Andreas Quirrenbach, for pointing out several important issues. LVEK and AGdeB acknowledge the support from an NWO program subsidy (grant number 781-76-101). This research was supported in part by the European Commission, TMR Programme, Research Network Contract ERBFMRXCT96-0034 ‘CERES’. The National Radio Astronomy Observatory is a facility of the National Science Foundation operated under cooperative agreement by Associated Universities, Inc. The National Radio Astronomy Observatory is a facility of the National Science Foundation operated under cooperative agreement by Associated Universities, Inc. The Westerbork Synthesis Radio Telescope (WSRT) is operated by the Netherlands Foundation for Research in Astronomy (ASTRON) with the financial support from the Netherlands Organization for Scientific Research (NWO).
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# Cosmological Constraints from the Cosmic Microwave Background
## 1 Introduction
The newest and perhaps most powerful tool in the cosmologist’s toolbox are the temperature fluctuations in the cosmic microwave background (CMB). Their very existence (e.g., Smoot et al. 1992) lends much credence to the general picture of gravitational instability forming galaxies and the observed large–scale structure. The last two decades have witnessed the elaboration of this idea, with numerical studies of gravitational growth permitting quantitative comparison to actual survey data, and with the introduction of a physical mechanism – namely, Inflation – for the creation of the initial perturbations. Daring in scope, the resulting scenario would encompass the evolution of the Universe from possibly the Planck era to the present, explaining not only the origin of the required density perturbations, but also dispelling a handful of misgivings about the initial conditions of Big Bang model, such as the impressive homogeneity on large scales (Kolb & Turner 1990; Peebles 1993; Peacock 1999) Like any theory, it can never be proved; but it may be tested. And like any good theory, it provides a physics, dependent on the validity of the theory, specific enough to be used as a tool to other ends: it is now well appreciated that detailed study of the CMB fluctuations may be applied to determine the fundamental parameters of the Big Bang model itself (Bond et al. 1994; Knox 1995; Jungmann et al. 1996).
Our object with the present study is to examine what may be learned from present CMB data within the Inflationary context. It should be noted at the outset that there are other models contending to explain the origin of density perturbations, for example, defect models (Durrer 1999). Neither these nor any other alternative shall be our concern in the following, and it is important to emphasize that our results are therefore limited to the Inflationary context. Very specifically, we shall be concerned only with temperature fluctuations caused by adiabatic, coherent and passively evolving density (scalar) perturbations, dominated by cold dark matter (CDM) (Bond & Efstathiou 1984; Vittorio & Silk 1984). Inflation can generate isocurvature modes, but we ignore these in the following, along with gravitational waves, the only non–scalar modes expected in Inflationary scenarios, and reionization. All model predictions have been calculated using the CAMB Boltzmann code developed by Lewis, Challinor & Lasenby (1999) and built upon CMBFAST (Seljak & Zaldarriaga 1996; Zaldarriaga, Seljak & Bertschinger 1998) The reason for these restrictions is of course that it is impossible to explore the otherwise vast parameter space.
Even in this seemingly restrained setting, the physics of CMB temperature anisotropies is quite rich, but it may nevertheless be boiled down to two regimes: large scales where purely gravitational effects operate \[Sachs–Wolfe (SW) effect; Sachs & Wolfe 1967\], and small scales, within the horizon at recombination where causal physics plays its role. In the latter regime, pressure of the coupled baryon–photon fluid resists the inward pull of gravity and therein establishes oscillating sound waves that will be observed in the CMB fluctuation power spectrum as a sequence of power peaks (so–called Doppler peaks; e.g., Hu & Sugiyama 1996). These are damped towards the smallest scales by smoothing due to the finite thickness of the last scattering surface from which emanates the CMB. The well defined aspect of these power peaks owes to the coherent nature and passive evolution of Inflation generated perturbations; the loss of both of these in defect models has the effect of broadening or completely smearing out the peaks (see, e.g,. Durrer 1999 and references therein). The very existence of such a series of power peaks is thus a strong discriminator between models.
Besides the exact details of the physics of Inflation, the development of such perturbations depends, unsurprisingly, on the constituents of the primordial plasma in which they reside, as well as on the metric background in which they evolve. This provides the link between the observable sky temperature fluctuations and the fundamental cosmological parameters. Actual data from the first generation of CMB experiments<sup>1</sup><sup>1</sup>1we use this term to refer to those experiments, with the exception of COBE, which were not specifically designed for map making; our detailed list is given in Table 2 and discussed below. already indicates the existence of the first Doppler peak, a fact that leads to interesting and non–trivial conclusions, as noted by many authors and as detailed below.
In this paper we present our constraints over a six–dimensional parameter space (see Table 1) resulting from a large compilation of first generation experiments (see Table 2). Our conclusions are based not on a traditional $`\chi ^2`$ approach, but on methods developed and presented elsewhere which attempt to account for the non–Gaussian nature of power spectrum estimates as well as information at times lost in simple band–power estimates \[Bartlett et al. 1999 (BDBL); Douspis et al. 2000 (DBBL); Douspis, Bartlett & Blanchard (DBB)\]; we have tested these methods against complete likelihood analyses over subsets of the present CMB data. Our parameter space here is spanned by the quantities listed in Table 1, namely, the Hubble constant $`H_\mathrm{o}`$; the total energy density $`\mathrm{\Omega }_{\mathrm{tot}}\mathrm{\Omega }_\mathrm{m}+\lambda _\mathrm{o}`$, where $`\mathrm{\Omega }_\mathrm{m}`$ is the matter density and $`\lambda _\mathrm{o}\mathrm{\Lambda }/3`$; the vacuum energy density parameter $`\lambda _\mathrm{o}`$; the baryon density, $`\mathrm{\Omega }_\mathrm{b}`$, in terms of $`\eta _{10}`$ (baryon–to–photon number ratio)$`\times 10^{10}`$ ($`\mathrm{\Omega }_\mathrm{b}h^2=0.00366\eta _{10}`$); and the primordial scalar spectral parameters $`n`$ and $`Q`$, i.e., the slope and normalization. Our data set is listed in detail in Table 2.
While some recent analyses have covered a larger parameter space (e.g., Tegmark & Zaldarriaga 2000), they are based on traditional $`\chi `$–squared methods. In a key work, Dodelson & Knox (1999) applied approximate likelihood methods similar to those employed here, although over a more restricted range of parameters. The present work thus extends the approximate likelihood approach to a larger parameter space and includes an appropriate GoF statistic. On the other hand, we have not considered the effects of calibration uncertainties; the agreement between our results and those of previous authors supports the idea that these do not make a drastic difference to the final conclusions.
The paper develops as follows: we begin in the next section with a presentation of our analysis methods; for the most part, this is a brief review of work presented in BDBL, DBBL and DBB. Our parameter constraints are then given in Section 3, followed by our conclusions in Section 4. The basic and robust result at this stage must be considered as the conclusion that the spatial curvature is close to zero and the primordial spectral index close to its scale–invariant value of unity. The results of our adapted statistical analysis are thus in agreement with many others (Lineweaver et al. 1997; Bartlett et al. 1998ab; Bond & Jaffe 1998; Efstathiou et al. 1998; Hancock et al. 1998; Lahav & Bridle 1998; Lineweaver & Barbosa 1998ab; Lineweaver 1998; Webster et al. 1998; Lasenby et al. 1998; Dodelson & Knox 1999; Melchiorri et al. 1999; Tegmark & Zaldarriaga 2000; Knox & Page 2000) and confirm the two fundamental predictions of the Inflationary model.
## 2 Analysis Method
The experimental results of Table 2 are shown in Figure 1 as standard band–power estimates, together with several theoretical curves. As for most current analyses, this will be our starting point. Simple as it may at first appear, a correct statistical approach to these data is in fact a non–trivial issue. The temptation is of course to apply a traditional $`\chi ^2`$ minimization to the ensemble of points and errors. This, however, is not strictly allowed, for these points, being power estimates, are not Gaussian distributed. In addition, the given errors are usually computed from the the band–power likelihood function and do not therefore necessarily represent the errors on the power estimate (which a frequentist would argue must be found by simulation). In general, the best statistical method to constrain parameters with CMB data would be a likelihood analysis based on the original pixel values (or temperature differences; in this paper, we use the term ‘pixel’ to refer to differences as well), because this is guaranteed to use all relevant experimental information. Straightforward to construct in the (present) context of Gaussian theories as a multi–variant Gaussian in the pixel values, whose covariance matrix depends on the underlying model parameters (and noise characteristics), the likelihood is in practice computationally cumbersome due to the required matrix operations (Bond, Jaffe & Knox 1998; Borrill 1999ab; Kogut 1999).
This difficulty has motivated us (BDBL; DBBL; DBB) and others (Bond, Jaffe & Knox 1998; Wandelt, Hivon & Górski 1998) to propose methods based on power estimates and aimed at reproducing as far as possible a complete likelihood analysis. Computation time is greatly reduced by working in the power plane (Bond, Jaffe & Knox 1998; Tegmark 1997), due to the much smaller number of data elements, but all proposed methods must be benched against the ultimate goal of recovering the likelihood results. In DBBL we evaluated at length the viability of several approximate methods by comparison to a complete likelihood analysis. We found that the traditional $`\chi ^2`$ method over power estimates (such as shown in Figure 1) is subject to bias. Other approaches based on approximate band–power likelihood functions fare better, but still do not always fully reproduce the complete likelihood results. As discussed in DBBL, this may be traced to relevant experimental information lost by the simple band–power representation of an experiment; we demonstrated this explicitly by showing that some MAX and Saskatoon data are sensitive to the local slope of the power spectrum. In other words, the observations constrain not only an in–band power, but also a local effective spectral slope. This information is correctly incorporated by a complete likelihood analysis, but obviously missed by any method based solely on the band–power estimates of Figure 1.
These comments motivate an approach in which all relevant experimental information is first identified by a set of parameters that are constrained by the data. A general likelihood function may then be constructed over this parameter space using the original pixel set; this need be done only once. Expressing a general power spectrum by these same parameters then permits one to assign a likelihood value to any model by extrapolation of the pre–calculated likelihood function. This value incorporates all pertinent information and should be the best possible approximation to the exact likelihood. The gain is that one manipulates the likelihood function with the full pixel set only once, and then simply interpolates over a much reduced set of parameters. Once the likelihood function has been calculated, the technique is hardly more complicated than the traditional $`\chi ^2`$ employed for its facility. Unfortunately, the information needed for a general application of this technique to all of the current experimental results is not readily available. Because the number of points we are able to treat in this fashion is thus small, we elect to analyze the entire data set with the band–power approximation developed in BDBL.
We examine a set of Inflationary models over the parameter space spanned by $`(H_\mathrm{o},\mathrm{\Omega }_{\mathrm{tot}},\lambda _\mathrm{o},\eta _{10},n,Q)`$, with respective ranges and step sizes given in Table 1. The likelihood of each model is calculated as just described, and the best model is found by maximizing the likelihood function over the explored space. We present our results as a series of two–dimensional contour plots of the likelihood projected onto various parameter planes. The contours are defined in the full six–dimensional space with values of $`\mathrm{\Delta }\mathrm{log}()=1,4`$ (dashed, in green), motivated as the 1 and 2 $`\sigma `$ contours of a Gaussian distribution when projected onto one of the axes, and of $`\mathrm{\Delta }\mathrm{log}()=2.3,6.17`$ (solid, in red), motivated by a Gaussian distribution with two degrees of freedom. Since the likelihood is not actually Gaussian, the confidence percentages associated with our contours are not exactly 1 and 2 “$`\sigma `$”; the technique is however standard practice. Our final results shown in the following figures have been obtained by excluding the Python V points. As discussed below, this is because the inclusion of this data set leads to a poor GoFfor the entire class of models considered. On formal grounds we would then be lead to reject the models, while on a technical note we worry that likelihood contours for a poor best–fit may give misleading constraints. Perhaps somewhat arbitrarly, we thus proceed in our analysis without these data. This difficulty would in part be alleviated by a complete treatment of calibration errors.
Another aspect of our analysis not incorporated, as far as we are aware, in previous work is the application of an adequate goodness–of–fit statistic (GoF). Once the best model, i.e., the most likely model, has been found, one is obliged to evaluate the quality of its description of the data. As for the likelihood function itself, our situation is complicated by the fact that the power estimates shown in Figure 1 are not Gaussian distributed variables; in particular, a traditional $`\chi ^2`$ GoF statistic is inadequate for the task. In DBB we proposed a GoF statistic readily applicable, if necessarily approximate, to band–power estimates. One requires distribution of these power estimates, $`\widehat{\delta T_{\mathrm{fb}}}`$, for a given, underlying model; this distribution is not the same as the band–power likelihood function (frequentist point–of–view). Remarkably, we found in DBB that the same parameters introduced in BDBL to approximate the band–power likelihood function could be re–employed in a slightly different fashion to yield the distribution of the power estimator. The technique was tested with Monte Carlo simulations of experiments for which we performed complete likelihood analyses; details may be found in DBB. The important point is that with just the best power estimate and a confidence interval, we may construct an approximation to the complete distribution of the power estimate from an experiment and, hence, a GoF statistic for the data ensemble shown in Figure 1.
## 3 Results
Although perhaps at first glance the observational situation shown in Figure 1 seems confused, in fact the first Doppler peak would appear clearly detected. Several different experiments viewing different regions of sky on this scale, such as BOOMERanG, Python V, Saskatoon, Toco and QMAP, all indicate the presence of a rise in power and, in some cases, a hint of the subsequent fall–off, which is also supported by other experiments at higher $`l`$. This is not to say that all the data follow exactly the same party line, one example being the rather low MSAM points around the peak, but one could argue that the general trend favors the presence of a rise in power over the scales expected for first Doppler peak of Inflationary scenarios.
The way to quantify these statements is by fitting a model to the data and examining its GoF statistic. Over the parameter space explored and excluding the Python V points for the moment, the data identify the model with $`(H_\mathrm{o},\mathrm{\Omega }_{\mathrm{tot}},\lambda _\mathrm{o},\mathrm{\Omega }_\mathrm{b}h^2,n,Q)=(60\mathrm{km}/\mathrm{s}/\mathrm{Mpc},1.0,0.3,0.032,1.06,16.0\mu \mathrm{K})`$ as the best fit; the corresponding model spectrum is shown in Figure 1 as the solid curve. The quality of the fit may be judged from the distribution of our GoF statistic (referred to as a generalized $`\chi _{\mathrm{gen}}^2`$) as shown in Figure 2. Assuming the data come from the adopted model, a value of $`\chi _{\mathrm{gen}}^2`$ as big or larger than the observed value (indicated by the heavy arrow) occurs with a probability of 0.014 (i.e., 1.4% of the time; indicated by the shaded area under the curve). Admittedly, this is a little low, but it does give a rather satisfying numerical representation of the impression given by the data (it seems reasonable). Although perhaps the fit is marginal, we certainly do not find the result sufficiently conclusive to eliminate the entire class of models from consideration. We are further comforted in this direction by the fact that the GoF is dominated by only a few outliers (see DBB). We thus do not hesitate to accept the model and move on to see what constraints the data provide over the parameter space considered. We remark in passing that a (incorrect) standard $`\chi ^2`$ statistic is even less kind to the model with a value only 0.0016 (i.e., 0.16%) probable. Keeping the Python V points in the analysis yields a very poor best–fit model, quantified by a value of $`\chi _{\mathrm{gen}}^2`$ only $`2\times 10^5`$ probable (marked on the figure by the dashed (green) curve and arrow). This is the reason for which we choose to exclude Python V in the following; thus, all our final results quoted hereafter and shown in the figures exclude this data set. Fully aware that such a procedure should always be undertaken only with caution (perhaps the fluctuations are not Gaussian, for example), we nevertheless feel that this is the most constructive approach at present. In some sense this is a purely formal argument, but we do worry about the interpretation of the contours from a likelihood with a poor GoF. In any case, the best–fit model is only slightly changed by inclusion of Python V: $`(H_\mathrm{o},\mathrm{\Omega }_{\mathrm{tot}},\lambda _\mathrm{o},\mathrm{\Omega }_\mathrm{b}h^2,n,Q)=(60\mathrm{km}/\mathrm{s}/\mathrm{Mpc},1.0,0.5,0.032,1.12,13.0\mu \mathrm{K})`$. For comparison, we also calculated the GoF for the so–called “concordance model” $`(H_\mathrm{o},\mathrm{\Omega }_{\mathrm{tot}},\lambda _\mathrm{o},\mathrm{\Omega }_\mathrm{b}h^2,n,Q)=(65\mathrm{km}/\mathrm{s}/\mathrm{Mpc},1.0,0.3,0.018,1,20.0\mu \mathrm{K})`$: we find a GoF 0.1% probable.
The presence of the first Doppler peak in the data permits the elaboration of non–trivial constraints over our parameter space. Within our present context of adiabatic Inflation–generated perturbations, this peak appears on the physical scale of the horizon at the moment of recombination, $`H_\mathrm{o}^1\sqrt{\mathrm{\Omega }_\mathrm{m}}`$. As the distance to the last scattering surface is also proportional to $`H_\mathrm{o}^1`$ ($`D_\mathrm{a}=H_\mathrm{o}^1d_\mathrm{a}(\mathrm{\Omega }_\mathrm{m},\lambda _\mathrm{o})`$) the Hubble constant has relatively little influence on the projected angular scale of the peak; rather, $`\mathrm{\Omega }_\mathrm{m}`$ and, most notably, the curvature of space (light ray focusing) control this observable scale (Blanchard 1984). For this reason, one should expect that the most robust result coming out of the present data set would be constraints in the $`(\mathrm{\Omega }_{\mathrm{tot}},\lambda _\mathrm{o})`$–plane, as shown in Figure 3. The remarkable conclusion is that a flat universe is prefered and that, in particular, models with negative curvature (low $`\mathrm{\Omega }_{\mathrm{tot}}`$) are eliminated. On the other hand, the data place only weak constraints on the value of the cosmological constant. This degeneracy between $`\mathrm{\Omega }_\mathrm{m}`$ and $`\lambda _\mathrm{o}`$ is consistent with the expectation that one constrains instead their combination defining the quantity $`\sqrt{\mathrm{\Omega }_\mathrm{m}}/d_\mathrm{a}(\mathrm{\Omega }_\mathrm{m},\lambda _\mathrm{o})`$.
If spatial flatness may be considered as one of the motivating principals and a key “prediction” of the overall Inflation paradigm, then another is certainly the form of the primordial density perturbation spectrum, $`n`$. Figure 4 shows our constraints in the $`(\mathrm{\Omega }_{\mathrm{tot}},n)`$–plane. The two most simple “predictions” of Inflation are rapidly evaluated on this diagram, and we see that the model fares quite well: zero spatial curvature ($`\mathrm{\Omega }_{\mathrm{tot}}=1\mathrm{\Omega }_\kappa =1`$) and the spectral index of a scale–invariant spectrum ($`n=1`$) both fall within the inner contour.
As already mentioned, the best–fit model appears acceptable (marginally) according to our GoF statistic. Since our statistic is based on the assumption of Gaussianity, this in particular implies that the data are consistent with Gaussian anisotropies, as also expected in the simplest Inflationary scenarios. The general, overall conclusion from the first generation CMB anisotropy experiments must then be their coherence with the rudimentary concepts of Inflation.
One surprise of our analysis concerns the baryon density, $`\mathrm{\Omega }_\mathrm{b}`$. As seen in Figure 5, the data indicate extremely high values of $`\mathrm{\Omega }_\mathrm{b}h^2`$, and this almost independently of the value of the other parameters. So–called low D/H values observed in some QSO absorption systems yield $`\eta _{10}5`$ (Tytler et al. 2000), while here the CMB data prefer even higher baryon densities; although within “$`2\sigma `$” the two remain consistent. We believe the origin of this result to be related to the relative height of the first Doppler peak. Some data, such as Toco 97 & 98 and Saskatoon, suggest a high peak followed by a deep trough; low power at higher $`l`$ is also supported by other experiments, like Viper. For illustration, we show as the dotted line in Figure 1 the best–fit model with a fixed $`\mathrm{\Omega }_\mathrm{b}h^2=0.011`$. The essential differences between this model and the overall best–fit model are (in the fomer case) a lower first peak and the absence of a deep trough before the appearance of the second peak. One important aspect (besides the anticipated large improvement in overall data quality) of the next generation mapping experiments (e.g., Archeops<sup>2</sup><sup>2</sup>2http://www-crtbt.polycnrs-gre.fr/archeops/
general.html, BOOMERanG<sup>3</sup><sup>3</sup>3http://astro.caltech.edu/l̃gg/boom/boom.html, and MAXIMA<sup>4</sup><sup>4</sup>4http://cfpa.berkeley.edu/group/cmb/gen.html) is their ability to calibrate on the CMB dipole. This new calibration method, plus the fact that the entire first peak will be covered by a single instrument, should help to reduce the uncertainty surrounding the peak heights. As we have just argued, this is particularly important for determination of such quantities as the baryon density, and it will be an important test of the present result favoring a high baryon density.
## 4 Conclusions
Our purpose in this paper has been to analyze the ensemble of first generation CMB anisotropy experiments to see what conclusions may be drawn concerning certain fundamental cosmological parameters from the CMB data alone. Our approach employs approximate likelihood methods that are adapted to power estimates, and which have been detailed elsewhere (BDBL, DBBL and DBB). Our primary conclusions are that a flat geometry ($`\mathrm{\Omega }_\kappa 0`$ or $`\mathrm{\Omega }_{\mathrm{tot}}1`$) and a scale–invariant primeval spectrum ($`n1`$) are favored, while strongly hyperbolic models are ruled out with high significance – in short, Inflation remains a good theory. Specifically, the best–fit model parameters are $`(H_\mathrm{o},\mathrm{\Omega }_{\mathrm{tot}},\lambda _\mathrm{o},\mathrm{\Omega }_\mathrm{b}h^2,n,Q)=(60\mathrm{km}/\mathrm{s}/\mathrm{Mpc},1.0,0.3,0.032,1.06,16.0\mu \mathrm{K})`$. Our analysis includes a GoF statistic that indicates that this model, and therefore the entire class (Inflation with adiabatic, scalar perturbations and without re–ionization) provides an acceptable description of the data. Many authors have recently explored these issues (Lineweaver et al. 1997; Bartlett et al. 1998ab; Bond & Jaffe 1998; Efstathiou et al. 1998; Hancock et al. 1998; Lahav & Bridle 1998; Lineweaver & Barbosa 1998ab; Lineweaver 1998; Webster et al. 1998; Lasenby et al. 1998; Dodelson & Knox 1999; Melchiorri et al. 1999; Tegmark & Zaldarriaga 2000; Knox & Page 2000), with various combinations of the present data set, and most would agree with these conclusions. The extensive spectral coverage (in multipole order $`l`$) over the first, and perhaps second Doppler peaks, and the low noise expected of the next generation instruments should qualitatively change the confidence in and precision of this kind of study.
Very little may be said about $`H_\mathrm{o}`$ or about the relative contributions of matter, $`\mathrm{\Omega }_\mathrm{m}`$, and vacuum, $`\lambda _\mathrm{o}`$, to the total energy density, the latter due to a well–known degeneracy when considering CMB data alone. One must turn to other observational constraints, coming from, for example, cluster evolution (Oukbir & Blanchard 1992), cluster baryon fractions (White et al. 1993), SNIa Hubble diagrams (Reiss et al. 1998; Perlmutter et al. 1999), weak cosmic shear (Blandford et al. 1991; Mellier 1999), etc…, to eliminate such degeneracies (so–called “cosmic complementarity”, Eisenstein, Hu & Tegmark 1999). Many of the authors listed above have included such constraints in their analysis; our own work along these lines is left to a forthcoming paper.
A surprising note is the very high baryon densities that we have found are prefered by the data set: $`\eta _{10}8.9`$ or $`\mathrm{\Omega }_\mathrm{b}h^20.032`$. This is even higher than the values indicated by the “low” D/H values found in several QSO absorption systems (Tytler et al. 2000), although within “$`2\sigma `$” (a little more for high values if $`H_\mathrm{o}`$; see Figure 5) the results are consistent. It remains to be seen if this intriguing result bears the scrutiny of the next generation experiments. In particular, their ability to significantly reduce overall calibration uncertainty by using the CMB dipole will be crucial to such issues.
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# Group Theoretical Quantum Tomography
## 1. introduction
In Quantum Mechanics a physical system is associated with a Hilbert space $``$: the states are described by positive trace-class trace-one operators $`T`$ on $``$, the physical quantities by self-adjoint operators $`A`$ on $``$ and the physical content of the theory is given by the expectation values $`\mathrm{Tr}(AT)`$. The state $`T`$ is completely determined by $`\mathrm{Tr}(Q_nT)`$ for $`Q_n`$ running on a suitable set $`\{Q_n\}`$ of observables and, for arbitrary operator $`A`$, $`\mathrm{Tr}(AT)`$ can be computed in terms of $`\mathrm{Tr}(Q_nT)`$. In order to implement this scheme one has to estimate $`\mathrm{Tr}(Q_nT)`$ experimentally, facing the problems arising from statistical errors and instrumental noise. Moreover, the number of experimental data is clearly finite, while $`A`$ and $`T`$ are operators on an infinite dimensional Hilbert space and the set $`\{Q_n\}`$ is infinite.
The problem of determining the state of a quantum system entered the realm of experiments in the last decade, in the domain of quantum optics. Many authors, see for example , proposed and used various techniques to reconstruct the density operator of a single mode of the e.m. field from the probability distributions of its quadratures. These methods were originally based on the use of the Radon transform, as in medical tomographic imaging. Due to this analogy the name quantum tomography is currently used to refer to these techniques. Their common feature, for a review see , is the use of a set of observables $`\{Q_n:nX\}`$, called quorum, parametrised by a space $`X`$ endowed with a probability measure $`\mu `$. The fundamental property of the quorum is that any observable $`A`$ can be expressed as integral transform on the space $`X`$
$$A=_X[A](n)𝑑\mu (n)$$
in such a way that, for all $`nX`$, the operator $`[A](n)`$ is a function of $`Q_n`$ in the sense of the functional calculus. Then, if $`T`$ is the state, one has that
(1)
$$\mathrm{Tr}(AT)=_{X\times }\sigma (A)(n,\lambda )\omega (n,\lambda )𝑑\mu (n)𝑑\lambda ,$$
where $`\lambda \omega (n,\lambda )`$ is the probability density of $`Q_n`$ in the state $`T`$, i.e.
$$\mathrm{Tr}(TQ_n)=_{}\lambda \omega (n,\lambda )𝑑\lambda ,$$
and $`\lambda \sigma (A)(n,\lambda )`$ is the function defined by $`[A](n)`$ using the functional calculus, i.e.
$$\mathrm{Tr}(T[A](n))=_{}\sigma (A)(n,\lambda )\omega (n,\lambda )𝑑\lambda $$
(in the above formulas we assumed for simplicity that each $`Q_n`$ has a continuous spectrum). Selecting randomly $`Q_n`$ in the quorum according to the probability measure $`\mu `$ and measuring it, the outcome probability of obtaining the value $`\lambda `$ is given by $`\omega (n,\lambda )d\mu (n)d\lambda `$. Then, by means of Eq. (1), the expectation value $`\mathrm{Tr}(AT)`$ can be reconstructed, by averaging the function $`\sigma (A)`$ over $`X\times `$ endowed with the probability measure $`\omega d\mu d\lambda `$. We notice that the function $`\sigma (A)`$, called the estimator of $`A`$, does not depend on $`T`$, and that the same set of data can be used to estimate all the expectation values $`\mathrm{Tr}(AT)`$.
In and a general method has been proposed to realize a quorum and define estimators in terms of suitable unitary representations of Lie groups (for a self contained synthetic exposition see and ). The present paper is devoted to the mathematical foundation of this method using the theory of square-integrable representations of unimodular Lie groups. In Section 2 we present the mathematical theory and in Section 3 we apply it to two examples: the homodyne tomography related to the Weyl-Heisenberg group and the angular momentum tomography associated with the rotation group.
## 2. Group-dynamical quorum
In this section we define a quorum associated with a square-integrable representation of a Lie group.
Let $`G`$ be a unimodular connected Lie group $`G`$ and $`K`$ a central closed subgroup. The quotient space $`H=G/K`$ is a unimodular connected Lie group. We denote by $``$ its Lie algebra, by $`m+1`$ the (real) dimension of $``$ as a vector space, by $`dv`$ the Lebesgue measure on $``$ and by $`dh`$ the Haar measure of $`H`$, uniquely defined up to a positive constant, which will be fixed in the following.
Denoted by $`\mathrm{exp}`$ the exponential map from $``$ to $`H`$, we assume that there is an open subset $`V`$ of $``$ such that $`\mathrm{exp}(V)`$ is open in $`H`$, its complement has zero measure with respect to $`dh`$ and $`\mathrm{exp}`$ is a diffeomorphism from $`V`$ onto $`\mathrm{exp}(V)`$. This hypothesis implies that, given $`fL^1(H,dh)`$,
(2)
$$_Hf(h)𝑑h=D_{}f(\mathrm{exp}(v))|\mathrm{det}(d(\mathrm{exp})_v)|\chi _V(v)𝑑v,$$
where $`d(\mathrm{exp})_v`$ is the differential of the exponential map at $`v`$, i.e.
$$d(\mathrm{exp})_v(w)=\left(\frac{d}{dt}\mathrm{exp}(v)\mathrm{exp}(v+tw)\right)_{t=0}w,$$
$`\mathrm{det}()`$ is the determinant and $`D`$ is a positive constant, see for example Th. 1.14, Ch. I of . We normalize the Haar measure $`dh`$ of $`H`$ in such a way that $`D=1`$.
###### Remark 1.
The density $`\mathrm{det}(d(\mathrm{exp})_v)`$ can be easily computed observing that, if $`\lambda _1,\mathrm{},\lambda _{m+1}`$ are the (possibly repeated) eigenvalues of $`d(\mathrm{exp})_v`$, viewed as linear operator on $``$, then
$$\mathrm{det}(d(\mathrm{exp})_v)=\frac{1e^{\lambda _1}}{\lambda _1}\mathrm{}\frac{1e^{\lambda _{m+1}}}{\lambda _{m+1}},$$
with $`\frac{1e^0}{0}=1,`$ see for example Th. 1.7, Ch. I of .
Let $`U`$ be an irreducible continuous unitary representation of $`G`$. We denote by $``$ the (complex separable) Hilbert space where the representation acts and by $`,`$ the scalar product, linear in the second argument.
We assume that the representation $`U`$ is square-integrable modulo $`K`$, i.e. there is a non-zero vector $`v`$ such that
(3)
$$_H|U_{c(h)}v,v|^2𝑑h<\mathrm{},$$
where $`c`$ is a section from $`H`$ to $`G`$, i.e. a measurable map $`c:HG`$ such that
$`c(e_H)`$ $`=`$ $`e_G`$
$`\pi (c(h))`$ $`=`$ $`hhH,`$
with $`\pi `$ being the canonical projection from $`G`$ to $`H`$. Notice that the value of the integral in Eq. (3) is independent of the choice of the section and that Eq. (3) implies that the function $`hU_{c(h)}u,w`$ is square integrable for all $`u,w`$, .
###### Remark 2.
In many examples $`K`$ is trivial, i.e. $`K=e_G`$, so that $`H=G`$ and Eq. (3) reduces to the usual notion of square-integrability. Nevertheless, there are cases, as the Weyl-Heisenberg group, that require the full theory. Moreover, in this framework one can easily consider projective representations. Indeed, let $`\widehat{U}`$ be a projective representation of a Lie group $`\widehat{H}`$ with multiplier $`m`$. Define $`G`$ as the central extension of the torus $`K`$ by $`\widehat{H}`$ associated with $`m`$. Then $`K`$ is a central closed subgroup of $`G`$, $`H`$ is canonically isomorphic with $`\widehat{H}`$ and there is a unitary representation $`U`$ of $`G`$ such that
$$\widehat{U}_{\pi (g)}=U_ggG.$$
Clearly, the fact that $`U`$ is square-integrable modulo $`K`$ is equivalent to the fact that $`\widehat{U}`$ is a square-integrable projective representation of $`\widehat{H}`$.
Being $`U`$ square-integrable modulo $`K`$, one can prove that there is a constant $`d>0`$, called the formal degree of $`U`$, such that, for all $`u_1,u_2,v_1,v_2`$,
(4)
$$_H\overline{U_{c(h)}v_1,u_1}U_{c(h)}v_2,u_2𝑑h=\frac{1}{d}u_1,u_2v_2,v_1.$$
Using the above relation we can represent the Hilbert-Schmidt operators as square integrable functions on $`H`$. Indeed, let $`^2()`$ be the Hilbert space of the Hilbert-Schmidt operators with the scalar product
$$(A,B)\mathrm{Tr}(A^{}B),$$
where $`\mathrm{Tr}()`$ denotes the trace and $`A^{}`$ is the adjoint operator of $`A`$. If $`u,v`$, let $`uv^{}`$ be the operator in $`^2()`$
$$(uv^{})(w)=v,wuw.$$
Given a section $`c`$, we define $`\mathrm{\Sigma }(uv^{})`$ as the function from $`H`$ to $``$ given by
$$\mathrm{\Sigma }(uv^{})(h)=U_{c(h)}v,uhH.$$
From Eq. (4), it follows that $`\mathrm{\Sigma }(uv^{})`$ is square-integrable with respect to $`dh`$ and
$$\mathrm{\Sigma }(uv^{})_{L^2(H,dh)}^2=\frac{1}{d}u^2v^2=\frac{1}{d}uv^{}_{^2()}^2.$$
Taking into account that the set $`\{uv^{}:u,v\}`$ is total in $`^2()`$, it follows that $`\mathrm{\Sigma }`$ is defined uniquely by continuity on $`^2()`$ and, if $`A,B^2()`$,
(5)
$$\mathrm{Tr}(A^{}B)=d\mathrm{\Sigma }(A),\mathrm{\Sigma }(B).$$
Moreover, if $`A`$ is of trace-class, then for almost all $`hH`$
(6)
$$\mathrm{\Sigma }(A)(h)=\mathrm{Tr}(U_{c(h)^1}A).$$
Indeed, let
$$A=\underset{i}{}\lambda _ie_if_i^{}$$
be the canonical decomposition of $`A`$, where $`(e_i)`$ and $`(f_i)`$ are orthonormal sequences in $``$, $`(\lambda _i)`$ is an $`\mathrm{}_1`$-sequence and the series converges in trace-norm and, hence, in the Hilbert-Schmidt norm. Since $`\mathrm{\Sigma }`$ is continuous, then
$$\mathrm{\Sigma }(A)=\underset{i}{}\lambda _i\mathrm{\Sigma }(e_if_i^{}),$$
where the series converges in $`L^2(H,dh)`$. On the other hand, fixed $`hH`$, since $`A`$ is of trace class, so is $`U_{c(h)^1}A`$, hence
$`\mathrm{Tr}(U_{c(h)^1}A)`$ $`=`$ $`{\displaystyle \underset{i}{}}f_i,U_{c(h)^1}Ae_i`$
$`=`$ $`{\displaystyle \underset{i}{}}\lambda _iU_{c(h)}f_i,e_i`$
$`=`$ $`{\displaystyle \underset{i}{}}\lambda _i\mathrm{\Sigma }(e_if_i^{})(h),`$
where the series converges pointwise. The claim is now clear.
We are now ready to define a quorum associated with the square-integrable (modulo $`K`$) representation $`U`$ of $`G`$.
Let $`T`$ be a state of $``$, i.e. a positive trace-class operator of trace one, and $`A`$ a Hilbert-Schmidt operator on $``$. Taking into account Eq. (5) and Eq. (6),
$`\mathrm{Tr}(TA)`$ $`=`$ $`d\mathrm{\Sigma }(T),\mathrm{\Sigma }(A)_{L^2(H,dh)}`$
$`=`$ $`d{\displaystyle _H}\overline{\mathrm{Tr}(U_{c(h)^1}T)}\mathrm{\Sigma }(A)(h)𝑑h,`$
so that
$$\mathrm{Tr}(AT)=d_H\mathrm{\Sigma }(A)(h)\mathrm{Tr}(TU_{c(h)})𝑑h.$$
By means of Eq. (2), the above equation becomes
$$\mathrm{Tr}(AT)=d_{}\mathrm{\Sigma }(A)(\mathrm{exp}v)\mathrm{Tr}(TU_{c(\mathrm{exp}v)})\chi _V(v)|\mathrm{det}(d(\mathrm{exp})_v)|𝑑v.$$
Let $`S^m`$ be the sphere in $``$. Then, for all $`nS^m`$, the map
$$tU_{c(\mathrm{exp}(tn))}$$
is a projective representation of $``$. Since all the multipliers of $``$ are equivalent to an exact one, there is a selfadjoint unbounded operator $`Q_n`$ and a measurable complex function $`\alpha _n`$ with modulo 1 such that, for all $`t`$,
(7)
$$U_{c(\mathrm{exp}(tn))}=\alpha _n(t)e^{itQ_n}.$$
Using polar coordinates in the above equation, one has that
$`\mathrm{Tr}(AT)`$ $`=`$ $`dC_m{\displaystyle _{S^m}}𝑑\mathrm{\Omega }(n){\displaystyle _0^{\mathrm{}}}𝑑tt^m\mathrm{\Sigma }(A)(\mathrm{exp}(tn))\alpha _n(t)`$
$`\mathrm{Tr}(Te^{itQ_n})\chi _V(tn)|\mathrm{det}(d(\mathrm{exp})_{tn})|,`$
where $`d\mathrm{\Omega }`$ is the normalized measure on the sphere $`S^m`$, $`C_m`$ is the volume of $`S^m`$ and $`dt`$ is the Lebesgue measure on the real line. The set of self-adjoint operators $`\{Q_n:nS^m\}`$, labelled by the probability space $`(S^m,d\mathrm{\Omega })`$, is called the quorum defined by the representation $`U`$. We notice that Eq. (7) defines $`Q_n`$ uniquely up to an additive constant, see, also, Remark 3 below.
Since $`Q_n`$ is selfadjoint, by means of the spectral theorem, there is a projection valued measure $`EP_n(E)`$ defined on $``$ such that
$$\mathrm{Tr}(TQ_n)=_{}\lambda 𝑑\mathrm{Tr}(TP_n(\lambda )),$$
where $`d\mathrm{Tr}(TP_n(\lambda ))`$ denotes the positive bounded measure on $``$
$$E\mathrm{Tr}(TP_n(E)).$$
Using this equation, one obtains that
$`\mathrm{Tr}(AT)`$ $`=`$ $`dC_m{\displaystyle _{S^m}}𝑑\mathrm{\Omega }(n){\displaystyle _0^{\mathrm{}}}𝑑t{\displaystyle _{}}𝑑\mathrm{Tr}(TP_n(\lambda ))`$
$`e^{i\lambda t}\mathrm{\Sigma }(A)(\mathrm{exp}(tn))\alpha _n(t)\chi _V(tn)|\mathrm{det}(d(\mathrm{exp})_{tn})|t^m.`$
In order to obtain a reconstruction formula for $`\mathrm{Tr}(AT)`$, we would like to interchange the integrals in $`dt`$ and in $`d\mathrm{Tr}(TP_n(\lambda ))`$.
We consider first the case when $`\mathrm{\Sigma }(A)`$, which is only square-integrable, is in fact integrable with respect to $`dh`$, i.e.
(10)
$$_H|\mathrm{\Sigma }(A)(h)|𝑑h<\mathrm{}.$$
By means of Fubini theorem, this condition implies that, for almost all $`nS^m`$, the map $`t\mathrm{\Sigma }(A)(\mathrm{exp}(tn))`$ is integrable with respect to the measure
(11)
$$dt_n=\chi _V(tn)|\mathrm{det}(d(\mathrm{exp})_{tn})|t^mdt.$$
Then the map from $`S^m\times `$ to $``$
(12)
$$\sigma (A)(n,\lambda )=dC_m_0^{\mathrm{}}e^{i\lambda t}\mathrm{\Sigma }(A)(\mathrm{exp}(tn))\alpha _n(t)\chi _V(tn)|\mathrm{det}(d(\mathrm{exp})_{tn})|t^m𝑑t,$$
is well-defined and it is called the estimator of the observable $`A`$. We notice that the estimator does not depend on $`T`$ and, given the representation $`U`$, can be computed analytically.
Since the measure $`d\mathrm{Tr}(TP_n(\lambda ))`$ is bounded, by means of Fubini theorem, one can interchange the integrals in Eq. (2) obtaining
(13)
$$\mathrm{Tr}(AT)=_{S^m}𝑑\mathrm{\Omega }(n)_0^{\mathrm{}}𝑑\mathrm{Tr}(TP_n(\lambda ))\sigma (A)(n,\lambda ).$$
The above integral transform is the core of the quantum tomography and is a concrete realization of the scheme proposed in the introduction, compare with Eq. (1). Indeed, $`d\mathrm{\Omega }(n)d\mathrm{Tr}(TP_n(\lambda ))`$ is the probability to obtain the value $`\lambda `$ when the observable $`Q_n`$, chosen randomly in the quorum according to $`d\mathrm{\Omega }`$, is measured. Moreover, by means of Eq. (13), the expectation value $`\mathrm{Tr}(AT)`$ can be reconstructed as average of the estimator $`\sigma (A)`$ over many random measures of the observables $`Q_n`$ in the quorum.
###### Remark 3.
There is a choice for the section that simplifies the expression of the estimator. Indeed, denoted by $`𝔊`$ the Lie algebra of $`G`$, since the differential $`d\pi `$ of $`\pi `$ is a surjective linear map from $`𝔊`$ onto $``$, there is an injective linear map $`j`$ from $``$ to $`𝔊`$ such that $`d\pi (j(v))=v`$ for all $`v`$. Since $`\mathrm{exp}`$ is a diffeomorphism from $`V`$ onto $`\mathrm{exp}(V)`$, it is well defined a smooth map $`\widehat{c}`$ from $`\mathrm{exp}(V)`$ to $`G`$ such that
$$\widehat{c}(\mathrm{exp}(v))=\mathrm{exp}(j(v))v.$$
Clearly $`\widehat{c}`$ is a section and the relation $`U_{\widehat{c}(\mathrm{exp}(tn))}=U_{\mathrm{exp}(tj(n))}`$ shows that one can always choose $`\alpha _n(t)=1`$ in Eq. (7). Hence $`U_{\widehat{c}(\mathrm{exp}(tn))}=e^{itQ_n}`$.
One can easily prove that, if one changes $`jj+l`$ in such a way that $`d\pi (j(v)+l(v))=v`$, then the quorum transforms according to $`Q_nQ_n+q_nI`$. However, in most of the cases, there is a natural choice for the map $`j`$, so that the quorum $`Q_n`$ is, in fact, defined uniquely by the representation $`U`$.
###### Remark 4.
Once the quorum $`\{Q_n\}`$ is fixed, Eq. (12) is independent of the choice of the section $`c`$. Indeed if $`c^{}`$ is another section, then, for all $`hH`$, $`c^{}(h)=k(h)c(h)`$ and $`k(h)K`$. Since $`K`$ is central in $`G`$ and $`U`$ is irreducible, then $`U_{k(h)}=\beta (h)I`$, where $`\beta (h)`$ is a complex number of modulo one. Hence, with obvious notations, for almost all $`hH`$ and for all $`t`$,
$`\mathrm{\Sigma }^{}(A)(h)`$ $`=`$ $`\overline{\beta (h)}\mathrm{\Sigma }(A)(h)`$
$`\alpha _{}^{}{}_{n}{}^{}(t)`$ $`=`$ $`\beta (h)\alpha _n(t),`$
so that $`\sigma (A)`$ is invariant with respect to the change $`cc^{}`$.
###### Remark 5.
If $`A`$ is of trace class and satisfies Eq. (10), using Eq. (6) one obtains a more explicit formula for the estimator of $`A`$
$$\sigma (A)(n,\lambda )=dC_m_0^{\mathrm{}}e^{i\lambda t}\mathrm{Tr}(Ae^{itQ_n})\chi _V(tn)|\mathrm{det}(d(\mathrm{exp})_{tn})|t^m𝑑t.$$
Moreover, in most examples the set $`V`$ is sufficiently nice so that the map $`n\chi _V(tn)`$ is continuous for almost all $`t`$. In this case, if one chooses the section $`\widehat{c}`$ as in Remark 3, taking into account that the function $`g\mathrm{Tr}(TU_g)`$ is continuous (since the ultra-weak operator topology is equivalent to the weak operator topology on the unit ball of $`()`$), it follows that the estimator $`\sigma (A)`$ is continuous on $`S^m\times `$. This property is important in order to approximate the integral of Eq. (13) by a finite sum.
###### Remark 6.
We notice that this procedure is unbiased since the observables $`Q_n`$ are chosen randomly and the integral given by Eq. (13) can be approximated by a finite sum since $`d\mathrm{\Omega }(n)d\mathrm{Tr}(TP_n(\lambda ))`$ is a probability measure. This means that this approach is not affected by the systematic errors that were present in the first tomographic scheme , due to the cutoff needed in the inversion of the Radon transform, see .
###### Remark 7.
If $`H`$ is compact then $`dh`$ is finite and any irreducible representation is square-integrable. Since the Hilbert space $``$ where the representation acts is finite dimensional, $`^2()`$ coincides with the space of the all operators. Moreover, since $`L^2(H,dh)L^1(H,dh)`$, Eq. (10) holds for every operator.
###### Remark 8.
If $`U`$ is an integrable representation (modulo $`K`$), there exists a dense set $`S`$ in $``$ such that, if $`u,vS`$, then $`\mathrm{\Sigma }(uv^{})`$ satisfies Eq. (10).
If condition (10) does not hold, it may happen that, for a non-negligible set of $`nS^m`$, the map $`t\mathrm{\Sigma }(A)(\mathrm{exp}(tn))`$ is not integrable with respect to the measure $`dt_n`$ defined by Eq. (11) (it is only square-integrable), so that the estimator $`\sigma (A)`$ given by Eq. (12) is not well defined.
In these cases, in order to define the estimator one has to use a suitable regularization procedure. For example, fixed $`L>O`$ let, for all $`nS^m`$ and $`\lambda `$,
(14)
$$\sigma _L(A)(n,\lambda )=dC_m_0^Le^{i\lambda t}\mathrm{\Sigma }(A)(\mathrm{exp}(tn))\alpha _n(t)\chi _V(tn)|\mathrm{det}(d(\mathrm{exp})_{tn})|t^m𝑑t.$$
It may be the case that there exists a function $`\sigma (A)`$ such that
$`\underset{L\mathrm{}}{lim}{\displaystyle _{S^m}}𝑑\mathrm{\Omega }(n){\displaystyle _{}}𝑑\mathrm{Tr}(TP_n(\lambda ))\sigma _L(A)(n,\lambda )=`$
$`{\displaystyle _{S^m}}𝑑\mathrm{\Omega }(n){\displaystyle _{}}𝑑\mathrm{Tr}(TP_n(\lambda ))\sigma (A)(n,\lambda ).`$
Then, as an easy consequence of dominated convergence theorem, one has that
$$\mathrm{Tr}(AT)=_{S^m}𝑑\mathrm{\Omega }(n)_{}𝑑\mathrm{Tr}(TP_n(\lambda ))\sigma (A)(n,\lambda ).$$
Analogous regularization procedures could be used to extend $`\mathrm{\Sigma }(A)`$ to non-Hilbert-Schmidt operators. Although this problem is physically relevant (many observables of interest are unbounded) it is out of the scope of the present paper.
## 3. examples
### 3.1. The Weyl-Heisenberg group
Let $`G`$ be the Weyl-Heisenberg group, i.e. $`G=^3`$ with the composition law
$$(\eta _1,a_1,b_1)(\eta _2,a_2,b_2)=(\eta _1+\eta _2+\frac{b_1a_2a_1b_2}{2},a_1+a_2,b_1+b_2).$$
It is known that $`G`$ is a connected simply-connected nilpotent unimodular Lie group.
The set $`K=\{(\eta ,0,0):\eta \}`$ is clearly a central closed subgroup of $`G`$ and the quotient group $`H=G/K`$ can be identified with the vector group $`^2`$. One has the following facts.
1. The canonical projection $`\pi `$ is given by $`\pi (\eta ,a,b)=(a,b)`$.
2. A smooth section $`c`$ is given by $`c(a,b)=(0,a,b)`$.
3. A Haar measure on $`H`$ is the Lebesgue measure $`dadb`$ of $`^2`$.
4. The Lie algebra $``$ of $`H`$ can be identified with $`^2`$ so that the exponential map is the identity and, for all $`v`$, $`\mathrm{det}(d(\mathrm{exp})_v)=1`$.
5. The constant $`D`$ in Eq. (2) is equal to $`1`$.
It follows that the choice $`V=`$ satisfies the assumptions of the previous section.
Let $`U`$ be the representation of $`G`$ acting in $`=L^2(,dx)`$ as
$$(U_{(\eta ,a,b)}u)(x)=e^{i(\eta +\frac{ab}{2})}e^{ixa}u(x+b)$$
where $`x`$, $`uL^2(,dx)`$ and $`(\eta ,a,b)G`$. It is known that $`U`$ is a unitary continuous irreducible representation of $`G`$, called the Schrödinger representation.
We prove that $`U`$ is square-integrable modulo $`K`$. Given $`u`$, the map
$$(x,b)\overline{u(x+b)}u(x)$$
is measurable and
$$_{}𝑑x_{}𝑑b|\overline{u(x+b)}u(x)|^2=u^4<\mathrm{}.$$
By Fubini theorem, for almost all $`b`$, the map
$$x\overline{u(x+b)}u(x)$$
is square-integrable and, since both $`xu(x+b)`$ and $`xu(x)`$ are square-integrable, it is integrable. Then, the map
$$a_{}𝑑xe^{iax}\overline{u(x+b)}u(x)$$
is well defined and square-integrable with respect to $`da`$ for almost all $`b`$. Moreover, the map
$$(a,b)_{}𝑑xe^{iax}\overline{u(x+b)}u(x)$$
is measurable and, by means of the isometry of the Fourier transform,
$`{\displaystyle _{}}𝑑b{\displaystyle _{}}𝑑a\left|{\displaystyle _{}}𝑑xe^{iax}\overline{u(x+b)}u(x)\right|^2`$ $`=`$ $`2\pi {\displaystyle _{}}𝑑b{\displaystyle _{}}𝑑x|\overline{u(x+b)}u(x)|^2`$
$`=`$ $`2\pi u^4.`$
Since
$$U_{c(a,b)}u,u=_{}𝑑xe^{iax}\overline{u(x+b)}u(x),$$
by Fubini theorem one has that $`U`$ is square-integrable and the formal degree is $`d=\frac{1}{2\pi }`$.
Since $`U`$ is square-integrable modulo $`K`$, according to Section 2, it defines a quorum. In order to explicit it, we observe that, with the notation of the previous section,
$$S^m=\{n_\mathrm{\Phi }:=(\mathrm{cos}(\mathrm{\Phi }),\mathrm{sin}(\mathrm{\Phi })):\mathrm{\Phi }[0,2\pi ]\},$$
$`m=1`$, $`C_1=2\pi `$ and $`d\mathrm{\Omega }=\frac{d\mathrm{\Phi }}{2\pi }`$. Moreover, since $`tU_{c(tn_\mathrm{\Phi })}`$ is a one parameter subgroup, then
$$U_{c(tn_\mathrm{\Phi })}=e^{itY_\mathrm{\Phi }}$$
where $`Y_\mathrm{\Phi }`$ is a selfadjoint operator (in this example $`\alpha _{n_\mathrm{\Phi }}(t)=1`$). If $`u`$ is a Schwartz function, one has that
$$Y_\mathrm{\Phi }u=\mathrm{cos}(\mathrm{\Phi })Qu+\mathrm{sin}(\mathrm{\Phi })Pu,$$
where $`Q`$ is the multiplicative operator by $`x`$, i.e. the position operator, and $`P`$ is $`i`$ times the weak derivative operator, i.e. the momentum operator. Hence the quorum defined by $`U`$ is given by the set of self-adjoint operators
$$\{Y_\mathrm{\Phi }:\mathrm{\Phi }[0,2\pi ]\}$$
labelled by the space $`[0,2\pi ]`$ with the uniform measure $`\frac{d\mathrm{\Phi }}{2\pi }`$.
The above quorum has the following property. For all $`\mathrm{\Phi }[0,2\pi ]`$, there is a unitary operator $`W_\mathrm{\Phi }`$ such that
(15)
$$Y_\mathrm{\Phi }=W_\mathrm{\Phi }QW_\mathrm{\Phi }^1.$$
To prove it, given $`\mathrm{\Phi }[0,2\pi ]`$, let $`f_\mathrm{\Phi }`$ from $`G`$ to $`G`$
$$f_\mathrm{\Phi }(\eta ,a,b)=(\eta ,\mathrm{cos}(\mathrm{\Phi })a\mathrm{sin}(\mathrm{\Phi })b,\mathrm{sin}(\mathrm{\Phi })a+\mathrm{cos}(\mathrm{\Phi })b).$$
One can easily check that $`f_\mathrm{\Phi }`$ is a continuous group automorphism of $`G`$, so that $`gU^{f_\mathrm{\Phi }}`$ is a unitary irreducible continuous representation of $`G`$ and the restriction to $`K`$ is the character $`\eta e^{i\eta }`$. From the unicity of the Schrödinger representation, it follows that there exists a unitary operator $`W_\mathrm{\Phi }`$ such that
$$U^{f_\mathrm{\Phi }}=W_\mathrm{\Phi }UW_\mathrm{\Phi }^1.$$
Then
$$U_{c(tn_\mathrm{\Phi })}=U_{(0,t,0)}^{f_\mathrm{\Phi }}=W_\mathrm{\Phi }U_{(0,t,0)}W_\mathrm{\Phi }^1,$$
and Eq. (15) follows by Stone theorem.
Let now $`T`$ be a state of $``$. Recalling that the spectral measure $`P_Q`$ of $`Q`$ is the one given by the multiplicative operators by characteristic functions, then, by means of Eq. (15), for each $`\mathrm{\Phi }[0,2\pi ]`$ there is a $`L^1(,d\lambda )`$-function $`\lambda \omega (\mathrm{\Phi },\lambda )`$ such that
$$\mathrm{Tr}(TP_\mathrm{\Phi }(E))=\mathrm{Tr}(W_\mathrm{\Phi }^1TW_\mathrm{\Phi }P_Q(E))=_E\omega (\mathrm{\Phi },\lambda )𝑑\lambda ,$$
where $`EP_\mathrm{\Phi }(E)`$ is the spectral measure associated with $`Y_\mathrm{\Phi }`$. The map $`\omega `$ can always be chosen measurable as a function on $`[0,2\pi ]\times `$ and, in doing so, it is a probability density on $`[0,2\pi ]\times `$ with respect to the measure $`\frac{d\mathrm{\Phi }}{2\pi }d\lambda `$.
Finally, fix a Hilbert-Schmidt operator $`A`$ in $``$ such that $`\mathrm{\Sigma }(A)`$ is integrable with respect to $`dadb`$. According to Eq. (12), the estimator of $`A`$ is, if $`\mathrm{\Phi }[0,2\pi ]`$ and $`\lambda `$,
$$\sigma (A)(\mathrm{\Phi },\lambda )=_0^{\mathrm{}}t\mathrm{\Sigma }(A)(t\mathrm{cos}(\mathrm{\Phi }),t\mathrm{sin}(\mathrm{\Phi }))e^{i\lambda t}𝑑t,$$
and the reconstruction formula Eq. (13) is explicitly given by
$$\mathrm{Tr}(AT)=_0^{2\pi }_{}\sigma (A)(\mathrm{\Phi },\lambda )\omega (\mathrm{\Phi },\lambda )\frac{d\mathrm{\Phi }}{2\pi }𝑑\lambda .$$
The representation $`U`$ is actually integrable and, if $`(u_n)`$ is the basis of eigenvectors of the number operator, then $`\mathrm{\Sigma }(u_nu_{n+l}^{})L^1(H,dh)`$ and one has the explicit formula
(16)
$$\sigma (u_nu_{n+l}^{})(\mathrm{\Phi },\lambda )=\frac{(i)^l}{2^{\frac{l}{2}}}\sqrt{\frac{n!}{(n+l)!}}e^{il\mathrm{\Phi }}_0^{\mathrm{}}t^{l+1}L_n^l(\frac{t^2}{2})e^{\frac{t^2}{4}+i\lambda t}𝑑t,$$
where $`L_m^k`$ are the associated Laguerre polynomials. The statistical reliability of Eq. (16) has been verified in .
This example is physically realized by the homodyne tomography, . The quantum system is the harmonic oscillator representing a single mode of the e.m. field with annihilation and creation operators $`\widehat{a}`$ and $`\widehat{a}^{}`$. In terms of such operators, one has the following translation table
$`Q`$ $`=`$ $`{\displaystyle \frac{\widehat{a}+\widehat{a}^{}}{\sqrt{2}}}`$
$`P`$ $`=`$ $`{\displaystyle \frac{\widehat{a}\widehat{a}^{}}{\sqrt{2}i}}`$
$`U_{(\eta ,a,b)}`$ $`=`$ $`e^{i\eta }e^{\left(\alpha \widehat{a}^{}\overline{\alpha }\widehat{a}\right)}`$
$`Y_\mathrm{\Phi }`$ $`=`$ $`\sqrt{2}{\displaystyle \frac{\widehat{a}^{}e^{i\mathrm{\Phi }}+\widehat{a}e^{i\mathrm{\Phi }}}{2}}=:\sqrt{2}X_\mathrm{\Phi }`$
where $`\alpha =\frac{b+ia}{\sqrt{2}}`$, $`e^{\left(\alpha \widehat{a}^{}\overline{\alpha }\widehat{a}\right)}`$ is displacement group and $`X_\mathrm{\Phi }`$ is the quadrature with phase $`\mathrm{\Phi }[0,2\pi ]`$.
The measuring apparatus is a homodyne detector with tunable phase with respect to the local oscillator. The function $`\sqrt{2}\omega (\mathrm{\Phi },\sqrt{2}\lambda )`$ is the probability density (with respect to $`d\lambda `$ ) to obtain the value $`\lambda `$ measuring the quadrature $`X_\mathrm{\Phi }`$, chosen randomly according to the measure $`\frac{d\mathrm{\Phi }}{2\pi }`$. Moreover, the explicit form of the estimator of $`A`$, being $`A`$ of trace class, is
$$\sigma (A)(\mathrm{\Phi },\sqrt{2}\lambda )=\frac{1}{2}_0^{\mathrm{}}\mathrm{Tr}(Ae^{it(X_\mathrm{\Phi }\lambda )})t𝑑t.$$
One could consult for an example of an experimental realization of the above tomographic method.
###### Remark 9.
In this example one is able to obtain an estimator also for for monomials in $`\widehat{a}`$ and $`\widehat{a}^{}`$, . For example, one has that
$$\sigma (\widehat{a}^{}\widehat{a})(\mathrm{\Phi },\sqrt{2}\lambda )=2\lambda ^2\frac{1}{2}.$$
### 3.2. The group $`SU(2)`$
Let $`SU(2)`$ be the group of the unitary $`2\times 2`$ complex matrices with determinant $`1`$. It is a unimodular connected simply-connected compact Lie group. The corresponding Lie algebra is
$$su(2)=\{\frac{i}{2}(x\sigma _1+y\sigma _2+z\sigma _3):x,y,z\}$$
where $`\sigma _i`$ are the Pauli matrices
$$\sigma _1=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)\sigma _2=\left(\begin{array}{cc}0& i\\ i& 0\end{array}\right)\sigma _3=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right).$$
In the following we identify $`su(2)`$ with $`^3`$ using the basis $`(\frac{i\sigma _k}{2})_{k=1}^3`$. Let $`V=\{(x,y,z)^3:\sqrt{(x^2+y^2+z^2)}<2\pi \}`$, it is known that $`V`$ is an open neighborhood of $`0`$ such that the exponential map restricted to $`V`$ is a diffeomorphism from $`V`$ onto the open set $`\mathrm{exp}(V)`$ and the complement of $`\mathrm{exp}(V)`$ is negligible with respect to the Haar measure of $`SU(2)`$. Moreover one can check that
$$d(exp)_{(x,y,z)}=4\frac{\mathrm{sin}^2(\frac{\sqrt{x^2+y^2+z^2}}{2})}{x^2+y^2+z^2}.$$
If we choose the Haar measure on $`SU(2)`$ in such a way that the constant $`D`$ in Eq. (2) is $`1`$ one has that
(17)
$$_H1𝑑h=_V|d(exp)_{(x,y,z)}|𝑑x𝑑y𝑑z=16\pi ^2,$$
(usually the Haar measure on compact groups is normalized to $`1`$).
Given $`j`$ such that $`2j`$, let $`D^j`$ be the irreducible representation of $`SU(2)`$ acting on $`=^{2j+1}`$. Since the group is compact, $`D^j`$ is square-integrable and the space of the Hilbert-Schmidt operators coincides with the space of all operators $`(^{2j+1})`$.
Since the measure of $`SU(2)`$ is normalized according to Eq. (17), it is well known that the formal degree is $`d=\frac{2j+1}{16\pi ^2}`$, see for example .
For all $`nS^2`$, define $`J_n`$ as the hermitian matrix such that
$$D^j(\mathrm{exp}(tn))=e^{itJ_n}t.$$
Then, the quorum defined by $`D^j`$ is the set of spin operators $`\{J_n:nS^2\}`$ labelled by the space $`S^2`$ with the measure $`\frac{dn}{4\pi }`$, being $`dn`$ the area element of the sphere. It is known that the (simple) eigenvalues of each $`J_n`$ are $`\lambda =j,\mathrm{},j`$ and there exists a unitary operator $`W_n`$, unique up to a phase, such that
$$Q_n=W_{n}^{}{}_{}{}^{1}J_zW_n,$$
where $`J_z=J_{(0,0,1)}.`$
Let now $`A(^{2j+1})`$, then, according to Eq. (12) and taking into account that $`C_2=4\pi `$, the corresponding estimator is
$$\sigma (A)(n,\lambda )=\frac{2j+1}{\pi }_0^{2\pi }e^{i\lambda t}\mathrm{Tr}(Ae^{itJ_n})\mathrm{sin}^2(\frac{t}{2})𝑑t,$$
where $`nS^2`$ and $`\lambda =j,\mathrm{},j`$. Equation (13) becomes
$$\mathrm{Tr}(TA)=\underset{\lambda =j}{\overset{j}{}}_{S^2}\sigma (A)(n,\lambda )|W_ne_\lambda ,TW_ne_\lambda |^2\frac{dn}{4\pi },$$
where $`(e_\lambda )_{\lambda =j}^j`$ is a basis of eigenvectors of $`J_z`$.
This example is realized experimentally by a Stern-Gerlach machine. The quantum system is the spin degree of freedom of an elementary particle with spin $`j`$ and the number $`|W_ne_\lambda ,TW_ne_\lambda |^2`$ is the probability to obtain the value $`\lambda `$ measuring the spin along the axis $`n`$, chosen randomly according to the measure $`\frac{dn}{4\pi }`$.
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# Gauge factor of thick film resistors: outcomes of the variable range hopping model
## I introduction
Thick film resistors (TFRs) are composite materials in which metallic grains (RuO<sub>2</sub>, Bi<sub>2</sub>Ru<sub>2</sub>O<sub>7</sub>, etc.) are embedded in an insulating glassy matrix. The characteristic transport properties of these materials render the TFRs particularly suitable as cryogenic thermometers and piezoresistive sensors. Despite this successfull practical aspect, on the theoretical side the situation is less satisfying. The microscopic mechanism of transport in TFRs is in fact far from to be understood and several theoretical models have been proposed, none of them however being able to provide a completely satisfactory and definite answer.
Among the proposed transport mechanisms in TFRs, the most representative ones are those based on tunneling within a network of interconnected metallic clusters separated by thin glassy layers , hopping between isolated metallic grains and phonon-assisted variable-range-hopping mechanism . The first model should be consistent with a highly dishomogenous structure in which the metallic phase is organized mainly in large segments separated by thin layers of glass while the last two models are more suitable for homogeneous dispersions of separated metallic grains. The actual situation seems to lay somewhere in between these two structures. Recent extended x-ray absorption fine structure experiments in fact have revealed a bimodal distribution of the metallic particle sizes in RuO<sub>2</sub>-based TFRs. According to these measurements, the metallic phase is organized in large RuO<sub>2</sub> grains or clusters with linear size ranging from $`200`$ Åto $`6000`$ Åand much smaller RuO<sub>2</sub> grains with sizes of order $`20`$-$`70`$ Å. Where the actual current path takes place is still under debate. However, It is plausible that transport in TFRs takes place mainly along paths connecting the small metallic grains, the large RuO<sub>2</sub> clusters being in fact too far away (a distance of order $`100`$ Åas inferred from transmission electron microscopy ) from each other for an electron to directly hop from a large cluster to another large cluster.
TFRs have also a characteristic temperature dependence of the resitance $`R(T)`$. At around room temperature, $`R(T)`$ has a minimum and it slightly increases at higher $`T`$. At lower temperatures, $`R(T)`$ increases signalling a non-metallic behavior. In this low-temperature region, $`R(T)`$ has been observed to follow a $`\mathrm{exp}(T_0/T)^x`$ behavior with $`x1/4`$ or $`x1/2`$ depending on the room temperature resistance value of the samples. Recent measurements have shown that $`R(T)\mathrm{exp}(T_0/T)^x`$ with $`x=1/4`$ with a cross-over to $`x=1/2`$ as the temperature is lowered for the most resistive samples. Such kind of temperature dependence is typical of Mott variable-range-hopping (VRH) mechanism of transport affected by Coulomb interaction between grains (Coulomb gap effect ). Hence, the transport properties of TFRs at low temperatures seem to be qualitatively described by ordinary VRH as other disordered systems like $`n`$-type CdSe . However, the values of $`T_0`$ and $`T_0^{}`$ that best fit the experimental data would lead to values of the optimal hopping length $`R_h`$ of the same order of the localization length $`\xi `$, whereas the VRH mechanism should be observed only for $`R_h\xi `$ . The validity of the VRH mechanism is therefore not clear and additional informations are required to analyse the validity of the VRH approach in describing transport in TFRs.
Here we propose that a quantity which can be helpful and easily measurable is the strain sensitivity, or gauge factor, and its temperature dependence. The gauge factor (GF) relates the variation of the total resistance $`R(T)`$ with an applied strain $`\epsilon `$ and it is defined as follows:
$$\mathrm{GF}=\frac{\delta R(T)}{\epsilon R(T)},$$
(1)
where $`\delta R(T)`$ is the variation of $`R(T)`$ under the applied strain $`\epsilon `$. TFRs have high values of GF, typically between $`2`$ and $`35`$ at room temperature , and it is precisely this property that renders these systems particularly suitable for piezoresistive sensor applications. Here we are interested in studying how the VRH mechanism affects GF and its temperature dependence. We find that if the low temperature transport is due to the VRH mechanism then the gauge factor acquires a characteristic temperature dependence which could be experimentally determined.
In the next section we use the asymptotic regimes predicted by the VRH model to extract qualitative behaviors for $`GF`$. In Sec. III, a simple model is introduced with the aim of studying the effect of finite inter-grain distances and to test the validity of the qualitative behaviors of GF.
## II Qualitative results
In this section we analyze the transport properties of TFRs from the point of view of VRH model. To this end, we must assume drastic simplifications and we model therefore the TFRs structure as given by an homogeneous dispersion of small metallic grains (of typical sizes $`20`$-$`70`$ Å) in the amorphous matrix and we neglect any possible effect of the larger metallic clusters on the bulk resistivity. As we discuss below, the assumption of homogeneity probably leads to a strong under-estimation of the room temperature GF value. However, instead of absolute values, here we are mainly interested in relative changes of GF as the temperature is varied.
An additional simplification we assume is to disregard the high temperature rise of $`R(T)`$ beyond the temperature of minimum. Although the origin of this feature is still debated, it should nevertheless been given by electron-phonon intra-grain scattering, thermal expansion effects or maybe a combination of both. Within these approximations, at sufficiently high temperatures transport is governed by tunneling between adjacent grains which are separated by a mean distance $`d`$. In this high temperature regime therefore
$$R(T)R_0\mathrm{exp}(2d/\xi )$$
(2)
where $`\xi `$ is the localization length and it is related to the barrier potential $`V`$ between adjacent grains as $`\xi =\mathrm{}/\sqrt{2mV}`$, where $`m`$ is the electron mass. An applied strain modifies $`d`$ and the resulting longitudinal GF becomes
$$\mathrm{GF}\mathrm{GF}_0=2d/\xi .$$
(3)
In the above expression it is implicitly assumed that the electrons hop along paths mostly parallel to the direction in which strain is applied. However, the actual microscopic paths are the result of hopping processes also along directions perpendicular to the applied strain. In this way, the longitudinal gauge factor is something less than $`2d/\xi `$ and the amount of reduction depends on specific material properties. However, as shown in appendix A, such a reduction can be estimated by introducing a phenomenological parameter $`\chi `$ which measures the percentage of hops along directions perpendicular to the strain, leading to the following expression for the longitudinal gauge factor:
$$\mathrm{GF}_0=\frac{1\chi /(1\nu )}{1+\chi }2d/\xi ,$$
(4)
where $`\nu `$ is the Poisson ratio which varies between $`0.2`$ and $`0.4`$ for typical TFRs . As shown in appendix A, $`0<\chi <1/2`$ and consequently $`0.1(2d/\xi )<\mathrm{GF}_0<2d/\xi `$. Since $`\xi `$ is of order of $`d`$, these values of GF<sub>0</sub> are much less than those measured in TFRs ($`2<\mathrm{GF}_0<35`$). We think that a strong enhancement of GF<sub>0</sub> could be achieved when the assumption of homogeneity is relaxed. In fact, a strongly non-homogeneous distribution of metal grains in the glassy matrix could lead to a local strain distribution very different from the averaged one. In particular, the paths along which the current flows can enter regions of concentrated strain leading therefore to an enhanced total GF<sub>0</sub>. Studies along this direction are currently under developement.
Having obtained the high temperature resistance and strain sensitivity, we analize now the situation at low temperatures. According to the VRH model, at sufficiently low temperatures hopping to adjacent grains is no longer favourable and $`R(T)`$ acquires a temperature dependence. The question is whether the VRH mechanism of transport affects the low temperature regime of GF as well and if this variation can be experimentally tested.
In the VRH model, the transport properties are governed by the probability $`P_{ij}`$ that an electron hops from grain $`i`$ to grain $`j`$ :
$$P_{ij}\mathrm{exp}\left[2r_{ij}/\xi \frac{E(r_{ij})}{K_\mathrm{B}T}\right],$$
(5)
where $`r_{ij}`$ is the distance between grains $`i`$ and $`j`$ and $`E(r_{ij})`$ is the energy threshold the electron experiences in hopping the distance $`r_{ij}`$. Assuming that a finite fraction of the grains is charged, the energy $`E(r_{ij})`$ is made of two contributions :
$$E(r_{ij})=\frac{\beta }{gr_{ij}^3}+\frac{\gamma e^2}{\kappa r_{ij}},$$
(6)
where $`g`$ is the electron density of states, $`e`$ is the electron charge, $`\kappa `$ is the dielectric constant, $`\beta `$ and $`\gamma `$ are dimensionless constants. In the second right hand of equation (6), the first term describes the energy needed to hop to a grain a distance $`r_{ij}`$. This term is proportional to $`r_{ij}^3`$ because the probability an electron at site $`i`$ has to find a site $`j`$ with energy much nearer to its own scales as $`r_{ij}^3`$. The second term in Eq. (6) describes the Coulomb energy due to finite charging of the grains, as it is expected by the fluctuations in the energy levels in the grains. This term is responsible for the opening of the Coulomb gap.
At sufficiently low temperatures, the resistance is governed by $`1/P_{\mathrm{max}}`$ where $`P_{\mathrm{max}}`$ is the maximum of $`P_{ij}`$. Hence, the optimization of the exponential in the hopping probability (5) leads to a temperature dependence of $`R(T)`$ characterized by two cross-over temperatures:
$$T_0=\frac{2048\beta }{27g\xi ^3K_B},$$
(7)
which stems from the first term in the second hand side of Eq.(6), and
$$T_0^{}=\frac{8\gamma e^2}{\kappa \xi K_B},$$
(8)
which is given by the Coulomb interaction. Usually $`T_0^{}<T_0`$ and $`R(T)\mathrm{exp}(T_0/T)^{1/4}`$ for $`T_0^{}<T<T_0`$ while $`R(T)\mathrm{exp}(T_0^{}/T)^{1/2}`$ for $`T<T_0^{}`$.
Depending on the particular temperature range, the strain sensitivity of $`R(T)`$ is governed by the strain sensitivity of $`T_0`$ or $`T_0^{}`$. Let us first focus on the temperature range in which $`R(T)\mathrm{exp}(T_0/T)^{1/4}`$. An applied strain $`\epsilon `$ leads to following relation between the resistance variation $`\delta R`$ and $`\delta T_0`$:
$$\frac{\delta R}{R}=\frac{1}{4}\left(\frac{T_0}{T}\right)^{1/4}\frac{\delta T_0}{T_0}.$$
(9)
$`T_0`$ is inversely proportional to the electron density of states $`g`$ which is proportional to the density of grains $`n`$. Assuming that an applied strain affects the density of metallic grains leaving the volume of the grains unchanged , we obtain $`\delta T_0/T_0=\delta g/g=\epsilon (12\nu )`$ and the gauge factor $`\mathrm{GF}=\delta R/\epsilon R`$ becomes:
$$\mathrm{GF}=\frac{12\nu }{4}\left(\frac{T_0}{T}\right)^{1/4}.$$
(10)
GF therefore increases as $`T^{1/4}`$ as the temperature is lowered provided $`T`$ is in the range for which $`R(T)\mathrm{exp}(T_0/T)^{1/4}`$ holds true.
By further lowering the temperature, transport becomes affected by the Coulomb interaction and $`R(T)`$ crosses towards the $`\mathrm{exp}(T_0^{}/T)^{1/2}`$ regime. Hence, under applied strain, $`\delta R/R`$ reduces to:
$$\frac{\delta R}{R}=\frac{1}{2}\left(\frac{T_0^{}}{T}\right)^{1/2}\frac{\delta T_0^{}}{T_0^{}}.$$
(11)
The variation $`\delta T_0^{}`$ is driven by the strain dependence of the total dielectric constant $`\kappa `$ and, in full generality, it can be expressed as:
$$\frac{\delta T_0^{}}{T_0^{}}=\frac{\delta \kappa }{\kappa }=\lambda \epsilon (12\nu ),$$
(12)
where $`\lambda `$ is a dimensionless parameter and its value depends on the dielectric constants of the metallic and glassy phases. In appendix B we provide an explicit expression of $`\lambda `$ based on the Maxwell-Garnett formula for the dielectric response of small metallic particles suspended in a dielectric. By using (11) and (12), the gauge factor in the Coulomb regime reduces to:
$$\mathrm{GF}=\frac{12\nu }{2}\lambda \left(\frac{T_0^{}}{T}\right)^{1/2}.$$
(13)
The Coulomb effect is therefore reflected in a cross-over from a $`T^{1/4}`$ to a $`T^{1/2}`$ dependence of GF as the temperature is sufficiently lowered.
## III Simple model for the GF temperature dependence
Summarizing the main results obtained in the previous section, the temperature dependence of the resistance and the intrinsic GF as $`T`$ is lowered is characterized, within the VRH theory, by three well distinguishable trends:
$$\begin{array}{ccc}R\mathrm{exp}(2d/\xi )\hfill & \hfill & \mathrm{GF}2d/\xi \hfill \\ R\mathrm{exp}\left(\frac{T_0}{T}\right)^{1/4}\hfill & \hfill & \mathrm{GF}\left(\frac{T_0}{T}\right)^{1/4}\hfill \\ R\mathrm{exp}\left(\frac{T_0^{}}{T}\right)^{1/2}\hfill & \hfill & \mathrm{GF}\left(\frac{T_0^{}}{T}\right)^{1/2}\hfill \end{array}$$
(14)
Note that for each behavior listed above, GF is proportional to $`\mathrm{ln}\rho (T)`$. The last two behaviors originate from the optimization of the inter-grain hopping probability $`P_{ij}`$ of Eq.(5). It is possible in fact to define a universal functional form for $`R(T)`$ which contains both $`\mathrm{exp}(T_0/T)^{1/4}`$ and $`\mathrm{exp}(T_0^{}/T)^{1/2}`$ as limiting values and which fits well transport data of Indium doped CdSe samples. However, for sufficiently high temperatures, the optimization of $`P_{ij}`$ requires optimal hopping distances of order or less than the mean distance $`d`$ between adjacent grains. It would therefore be more correct to formulate the transport problem in such a way that the value of $`d`$ is explicitly included. Here we provide a simple version of such a formulation capable of describing automatically the limiting behaviors listed in Eq.(14).
A simple general equation for the conductance $`G`$ can be defined as follows. A single hopping process between grains at distance $`r_{ij}`$ and with a hopping probability $`P_{ij}`$ can be regarded as a resistive element with resistance $`R_{ij}`$ proportional to $`1/P_{ij}`$. Hence, the total resistance can be constructed by interpreting paths characterized by different hopping distances as resistors in a parallel geometry. The resulting conductance should therefore be given by a weighted summation of the hopping probabilities $`P_{ij}`$ for all values of $`r_{ij}d`$. In this way, by employing a continuous approximation, we define the conductance $`G(T)`$ as:
$$G(T)=a𝑑rf(r)\mathrm{exp}\left[2r/\xi \frac{E(r)}{K_\mathrm{B}T}\right],$$
(15)
where $`E(r)`$ is given in Eq.(6) and $`a`$ is a constant introduced for dimensional purposes. In the above equation, $`f(r)`$ is a weight function which takes into account the different path geometries. For our purposes, the detailed structure of $`f(r)`$ is unnecessary and we retain only its main feature which is given by a lower cut-off representing the mean distance $`d`$ between adjacent grains. Hence, we approximate $`f(r)`$ by the step function $`\theta (rd)`$. In this way, equation (15) reduces to:
$$G(T)=a_d^+\mathrm{}𝑑r\mathrm{exp}\left(2r/\xi \frac{E(r)}{K_\mathrm{B}T}\right).$$
(16)
For temperatures larger than
$$T_s=\frac{\xi E(d)}{2dK_\mathrm{B}},$$
(17)
equation (16) gives the proper high temperature limit (2), while for $`TT_s`$ the integral can be estimated by optimizing the exponential, leading therefore to the VRH regime. Concerning the gauge factor resulting from (16), it should be noted that an applied strain $`\epsilon `$ affects both the interval of integration and the energy threshold function $`E(r)`$ through the electron density of states $`g`$ and the dielectric constant $`\kappa `$. Therefore, by using the strain derivatives of $`g`$ and $`\kappa `$ introduced in the previous section, GF can be expressed as the sum of three different contributions, $`\mathrm{GF}=\mathrm{GF}_d+\mathrm{GF}_g+\mathrm{GF}_\kappa `$ where:
$$\mathrm{GF}_d=\frac{2d}{\xi }\left[\frac{1\chi /(1\nu )}{1+\chi }\right]\frac{G(T)}{G_0}e^{E(d)/K_\mathrm{B}T}$$
(18)
$$\mathrm{GF}_g=\frac{12\nu }{K_\mathrm{B}T}\frac{a}{G(T)}_d^+\mathrm{}𝑑r\frac{\beta }{gr^3}\mathrm{exp}\left(2r/\xi \frac{E(r)}{K_\mathrm{B}T}\right)$$
(19)
$$\mathrm{GF}_\kappa =\frac{12\nu }{K_\mathrm{B}T}\frac{\lambda a}{G(T)}_d^+\mathrm{}𝑑r\frac{\gamma e^2}{\kappa r}\mathrm{exp}\left(2r/\xi \frac{E(r)}{K_\mathrm{B}T}\right)$$
(20)
In equation (18) $`G_0=a\xi \mathrm{exp}(2d/\xi )/2`$ is the asymptotic conductance in the high temperatures limit.
In Fig. 1a we plot the resistance $`R(T)`$ obtained numerically from Eq.(16) in the absence of Coulomb repulsion ($`\gamma =0`$) for $`2d/\xi =0.2`$, $`2`$, and $`5`$. The resistance data are in units of the high temperature limit $`R_0=1/G_0`$ and the temperature is given in units of $`T_0`$ \[see Eq.(7)\]. The curves show two distinct regimes as a function of temperature due to the finite values of the averaged inter-grain distance $`d`$. By using Eq.(17), the cross-over between these two regimes takes place roughly at $`TT_s=0.105T_0/(2d/\xi )^4`$, so that as $`2d/\xi `$ increases the cross-over temprature $`T_s`$ is reduced. For $`T>T_s`$ \[$`(T_0/T)^{1/4}<1.755(2d/\xi )`$\] the resistance shows an activated regime while at sufficiently low temperatures \[$`(T_0/T)^{1/4}>1.755(2d/\xi )`$\] it becomes of VRH type. The latter behavior is signalled by the dashed lines which fit the resistance with $`R(T)/R_0\mathrm{exp}(\stackrel{~}{T}_0/T)^{1/4}`$, where $`\stackrel{~}{T}_00.75T_0`$. The gauge factors GF associated to the resistance curves of Fig. 1a are plotted in Fig. 1b for $`\nu =0.3`$ and $`\chi =1/4`$. A $`2d/\xi `$ dependent cross-over is obtained for GF as well. For temperatures higher than $`T_s`$, GF is independent of $`T`$ and approaches the limit $`\mathrm{GF}_0`$ given in Eq.(4). On the contrary, by lowering $`T`$, GF acquires a temperature dependence and for $`TT_s`$ increases as $`1/T^{1/4}`$ in agreement with Eq.(10). From the results of Fig. 1b, the $`1/T^{1/4}`$ dependence of the gauge factor holds true in a wide range of temperatures only for quite small values of $`\mathrm{GF}_0`$.
The effect of Coulomb repulsion on $`R(T)`$ and GF is displayed in Fig. 2 where the case $`2d/\xi =2`$ is investigated for different values of the Coulomb parameter $`T_0^{}`$ defined in Eq.(8). In Fig. 2a, for $`T_0^{}/T_00`$ the resistance approach the low temperature asymptotic regime $`R(T)/R_0\mathrm{exp}(\stackrel{~}{T}_0^{}/T)^{1/2}`$ with $`\stackrel{~}{T}_0^{}0.9T_0^{}`$ (dashed lines). The Coulomb effect is reflected also in the temperature dependence of GF, Fig. 2b, where the curves have been calculated by using $`\lambda =1`$, $`\nu =0.3`$ and $`\chi =1/4`$. The set of curves with $`T_0^{}0`$ show the low temperature asymptotic behavior $`\mathrm{GF}(\stackrel{~}{T}_0^{}/T)^{1/2}`$ with values of $`\stackrel{~}{T}_0^{}`$ somewhat larger than those obtained from the resistivity data ($`\stackrel{~}{T}_0^{}1.05T_0^{}`$). Finally the temperature dependence of GF for different values of $`\lambda `$ is plotted in Fig. 3 for the case $`2d/\xi =2`$ and $`T_0^{}/T_0=0.05`$.
## IV discussion and conclusions
The numerical solutions of equations (16)-(20) confirm the qualitative results given in the previous section, i.e., as the resistance crosses over the VRH regime, the gauge factor acquires a temperature dependence which is qualitatively summarized in Eq.(14).
The model here presented is very simplified and susceptible of various improvements. For example, from statistical arguments, it would be certainly more correct to evaluate the temperature dependence of $`K(T)`$ within the percolation theory for transport . This could be achieved by explicitly taking into account in the evaluation of the critical path the lower cut-off $`d`$ for the hopping distance and the strain dependence of $`g`$ and $`\kappa `$. As already stressed before, an additional simplifying assumption we have used is that of homogeneity. However the structure of TFRs is very complex, being made by a glass with dispersed metallic grains with bi-modal distribution of sizes. In this situation, under an applied strain, the stress distribution inside the sample could strongly modify GF as compared to the homogeneous case and to settle quantitatively this modification is certainly a very important issue.
A final open question regards the general applicability of the VRH concept when electrons hop between grains of finite size. In fact, in reality, the discrete electron energy levels of the grains are smeared by finite life-time effects given by the electron-phonon and electron-electron interactions together with the inter-grain tunneling coupling. At sufficiently low temperatures, the life-time effects are small and the discreteness of the energy levels requires phonon assisted hopping for electron transport. However, at sufficiently high temperatures, the energy smearing due to the finite electron life-time can be of the order of the level spacing. In this situation it can be easily shown that electrons can directly tunnel between grains without the assistance of phonons. Of course, a crucial role is played by the grain size which governs the energy level spacing. To our knowledge, the problem of how the VRH between grains is affected by the inclusion of finite life-time effects has never been addressed and its investigation is certainly of the most valuable interest.
In conclusion, the results here reported provide a clear prediction of the temperature dependence of the gauge factor in the hypothesis that transport is well described by the classical VRH theory. Hence, the simultaneous detection of the temperature dependences of the sheet resistance and the gauge factor can provide an useful experimental tool to test the validity of the VRH mechanism in TFRs. We note that some experimental data show that GF for commercial TFRs is largely independent of $`T`$ from room down to cryogenics temperatures . Our analysis would indicate therefore that these experimental results are not compatible with the VRH mechanism of transport. However, it is important to stress that the VRH model predicts a strong temperature dependence of GF only when the resistance crosses over a $`\mathrm{exp}(T_0/T)^x`$ type of regime ($`x=1/2`$ or $`1/4`$), a test that has not been performed in Ref.. Instead, a firm experimental test requires measurements of the temperature dependences of both $`R(T)`$ and GF.
###### Acknowledgements.
We would like to thank Th. Mäder and M. Prudenziati for interesting discussions.
## A
In this appendix we estimate the reduction of the gauge factor due to the contributions of hopping processes perpendicular to the direction on which strain is applied. To this end, we first consider an operative definition of the gauge factor. Let us consider a cantilever beam in the $`x`$-$`y`$ plane with axis lying on the $`x`$ direction on top of which is deposited the TFR. The thicknesses of the cantilever bean and of the TFR are measured in the $`z`$ direction. According to elasticity theory, bending of the cantilever beam produces a finite strain $`\epsilon _{xx}=\epsilon `$ in the $`x`$ direction and no strain in $`y`$ ($`\epsilon _{yy}=0`$). Moreover, if $`\nu `$ is the Poisson ratio, the strain of the TFR along the $`z`$ direction is $`\epsilon _{zz}=\epsilon \nu /(1\nu )`$.
To measure the longitudinal GF of the TFR, a potential difference must be applied along the $`x`$ direction. If the microscopic paths are made of discrete hopping processes exclusively along the direction of the field, the longitudinal GF would be equal to $`2d/\xi `$ just as in Eq.(2). Conversely, the transversal GF would be zero because in this case the field is applied in the $`y`$ direction for which $`\epsilon _{yy}=0`$. Let us consider now the situation in which there is a finite percentage $`\chi `$ of hops along directions perpendicular to that of the applied field. The resulting resistance $`\stackrel{~}{R}_x`$, when the field is applied along the $`x`$ direction, becomes $`\stackrel{~}{R}_x=(1\chi )R_x+\chi (R_y+R_z)`$, where $`R_x`$, $`R_y`$ and $`R_z`$ are resistances for hops in the $`x`$, $`y`$ and $`z`$ directions, respectively. Although $`R_x`$, $`R_y`$ and $`R_z`$ have the same average value, they differ when strain is applied. In fact:
$$\frac{\delta R_x}{\epsilon _{xx}R_x}=\frac{2d}{\xi };\frac{\delta R_y}{\epsilon _{xx}R_y}=0;\frac{\delta R_z}{\epsilon _{xx}R_z}=\frac{\nu }{1\nu }\frac{2d}{\xi }.$$
(A1)
Hence the longitudinal gauge factor $`\mathrm{GF}_L`$ becomes:
$`\mathrm{GF}_L{\displaystyle \frac{\delta \stackrel{~}{R}_x}{\epsilon _{xx}\stackrel{~}{R}_x}}`$ $`=`$ $`{\displaystyle \frac{(1\chi )\delta R_x+\chi \delta R_z}{\epsilon _{xx}[(1\chi )R_x+\chi (R_y+R_z)]}}`$ (A2)
$``$ $`{\displaystyle \frac{1\chi /(1\nu )}{1+\chi }}{\displaystyle \frac{2d}{\xi }}.`$ (A3)
By following the same reasoning, the transversal gauge factor $`\mathrm{GF}_T`$ can be estimated by realizing that the resistance for a field applied in the $`y`$ direction is $`\stackrel{~}{R}_y=(1\chi )R_y+\chi (R_x+R_z)`$ and consequentely:
$`\mathrm{GF}_T{\displaystyle \frac{\delta \stackrel{~}{R}_y}{\epsilon _{xx}\stackrel{~}{R}_y}}`$ $`=`$ $`{\displaystyle \frac{\chi (\delta R_x+\delta R_z)}{\epsilon _{xx}[(1\chi )R_y+\chi (R_x+R_z)]}}`$ (A4)
$``$ $`{\displaystyle \frac{\chi [1\nu /(1\nu )]}{1+\chi }}{\displaystyle \frac{2d}{\xi }}.`$ (A5)
If the Poisson ratio $`\nu `$ is known, the phenomenological parameter $`\chi `$ can be estimated by measuring the ratio between the longitudinal and the transversal GFs which, according to (A2) and (A4), is given by:
$$\frac{\mathrm{GF}_L}{\mathrm{GF}_T}=\frac{1\chi \nu }{\chi (12\nu )}.$$
(A6)
Experimentally this ratio is found to be larger than the unity and consequently the above expression leads to $`\chi <1/2`$.
## B
Here, we evaluate the strain dependence of the dielectric constant for TFRs. We model the structure as given by spherical metallic particles dispersed in a glassy matrix. The metallic and glassy phases occupy volumes $`V_\mathrm{m}`$ and $`V_\mathrm{g}`$, respectively, so that the volume fraction of the metallic phase is $`\varphi =V_\mathrm{m}/(V_\mathrm{m}+V_\mathrm{g})V_\mathrm{m}/V_\mathrm{g}`$ for low metallic concentrations. Within the dipole approximation, the total dielectric constant $`\kappa `$ is given by the Maxwell-Garnett formula :
$$\kappa =\kappa _\mathrm{g}\left[1+\frac{3\varphi }{1\varphi }\right],$$
(B1)
where $`\kappa _\mathrm{g}`$ is the dielectric constant of the glass. For sufficiently small values of $`\varphi `$, the variation $`\delta \kappa `$ given by an applied strain $`\epsilon `$ is:
$$\delta \kappa =\kappa \frac{\delta \kappa _\mathrm{g}}{\kappa _\mathrm{g}}3\kappa _\mathrm{g}\varphi \epsilon (12\nu ),$$
(B2)
where we have considered the metallic particles as perfectly rigid bodies compared to the glass so that $`\delta \varphi =\varphi \epsilon (12\nu )`$ .
To find the variation $`\delta \kappa _\mathrm{g}`$ we should know an explicit expression for $`\kappa _\mathrm{g}`$ which is however a difficult problem. Nevertheless, experiments suggest that dielectric response of glasses is under several aspects quite similar to that of ionic solids and, in particular, it has been shown that the variation of the glass polarizability with pressure resembles that of the ionic crystalline compounds . By starting with this latter observation we could argue that, for our purposes, the dielectric constant for simple glasses is roughly given by:
$$\kappa _\mathrm{g}=\frac{1+8\pi \alpha /3\mathrm{\Omega }}{14\pi \alpha /3\mathrm{\Omega }},$$
(B3)
where $`\mathrm{\Omega }`$ has dimension of volume and $`\alpha 1/K`$ is the ionic polarization with $`K`$ being the elastic constant. An expression similar to (B3) is also recovered in theoretical analyses based on coarse-graining approach (there, $`\mathrm{\Omega }`$ is a course-graining volume). An applied strain modifies $`\kappa _\mathrm{g}`$ via the volume $`\mathrm{\Omega }`$ and the elastic constant $`K`$. The strain dependence of the latter quantity can be guessed by observing that $`K=V_{\mathrm{ion}}^{\prime \prime }(r_0)`$ where $`V_{\mathrm{ion}}(r)`$ is the interaction energy of two ions a distance $`r`$ apart and $`r_0`$ is the equilibrium distance. A pedagogical model for $`V(r)`$ is $`V(r)=a/r+\mathrm{exp}(br)+\mathrm{const}.`$ , so that one obtains $`K=V_{\mathrm{ion}}^{\prime \prime }(r_0)1/r_0^2`$. This result suggests that $`\alpha /\mathrm{\Omega }`$ scales as $`1/[\mathrm{distance}]`$ and consequently from Eq.(B3) one obtains in a straightforward way:
$$\delta \kappa _\mathrm{g}\frac{(\kappa _\mathrm{g}1)(\kappa _\mathrm{g}+2)}{3}\epsilon (12\nu ).$$
(B4)
The above expression plugged into Eq.(B2) leads finally to:
$`{\displaystyle \frac{\delta \kappa }{\kappa }}`$ $``$ $`\left[{\displaystyle \frac{(\kappa _\mathrm{g}1)(\kappa _\mathrm{g}+2)}{3\kappa _\mathrm{g}}}+3{\displaystyle \frac{\kappa _\mathrm{g}}{\kappa }}\varphi \right]\epsilon (12\nu )`$ (B5)
$``$ $`\lambda \epsilon (12\nu ),`$ (B6)
where the last equality defines the dimensionless quantity $`\lambda `$ used in the main text.
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# On complexes equivalent to 𝕊³-bundles over 𝕊⁴
## Introduction
$`𝕊^3`$-bundles over $`𝕊^4`$ have played an important role in topology and geometry since Milnor showed that the total spaces of such bundles with Euler class $`\pm 1`$ are manifolds homeomorphic to $`𝕊^7`$ but not always diffeomorphic to it. In 1974, Gromoll and Meyer exhibited one of these spheres (a generator in the group of homotopy 7-spheres) as a double coset manifold i.e. a quotient of $`\mathrm{Sp}(2)`$ hence showing that it admits a metric of nonnegative curvature (cf. ). Until recently, this was the only exotic sphere known to admit a metric of nonnegative sectional curvature. Then in , K. Grove and W. Ziller constructed metrics of nonnegative curvature on the total space of $`𝕊^3`$-bundles over $`𝕊^4`$. They also asked for a classification of these bundles up to homotopy equivalence, homeomorphism and diffeomorphism. These questions have been addressed in many papers such as , , and more recently in . In this paper we attempt to fill the gap in the previous papers; we consider the problem of determining when a given CW complex is homotopy equivalent to such a bundle. The problem was motivated by : the Berger space, $`\mathrm{Sp}(2)/\mathrm{Sp}(1)`$, is a 7-manifold that has the cohomology ring of an $`𝕊^3`$-bundle over $`𝕊^4`$, but does it admit the structure of such a bundle? The fact that it cannot be a principal $`𝕊^3`$-bundle over $`𝕊^4`$ is straightforward and is proved in .
Let $`X`$ be a simply connected CW complex with integral cohomology groups given by
(1) $`H^i(X)`$ $`=\text{if }i=0,7`$
$`=_n\text{if }i=4`$
where $`n`$ is some fixed integer. We say that $`X`$ is oriented with fundamental class $`[X]`$ if a generator $`[X]H_7(X)`$ is specified. For oriented $`X`$ we define the linking form as,
$`b:`$ $`H^4(X)H^4(X)_n`$
$`xy\beta ^1(x),y[X]`$
where $`\beta :H^3(X,_n)H^4(X)`$ is the Bockstein isomorphism. The homomorphism $`[X]:H^4(X)H_3(X)`$ is obtained by capping with the fundamental class and $`,:H^3(X,_n)H_3(X)_n`$ is the Kronecker pairing. We now state our main theorem:
###### Theorem 1.
Let $`X`$ be a simply connected CW complex as above. Then $`X`$ is homotopy equivalent to an $`𝕊^3`$-bundle over $`𝕊^4`$ if and only if the following two conditions hold:
* The secondary cohomology operation $`\mathrm{\Theta }`$ is trivial, where
$$\mathrm{\Theta }:H^4(X,𝔽_2)H^7(X,𝔽_2)$$
corresponds to the relation $`\mathrm{Sq}^2\mathrm{Sq}^2=\mathrm{Sq}^3\mathrm{Sq}^1`$ in the mod 2 Steenrod algebra.
* The linking form $`b:H^4(X)H^4(X)_n`$ is equivalent to a standard form for some choice of orientation on $`X`$ i.e. there exists an isomorphism $`\psi :_nH^4(X)`$ such that $`b(\psi (x),\psi (y))=xy`$.
Using the method outlined on page 32 of , it is easy to show that if $`X`$ smoothable, then $`\mathrm{\Theta }`$ is trivial and we have,
###### Corollary 2.
Let $`M`$ be a simply connected 7-manifold with integral cohomology groups given by
$`H^i(M)`$ $`=\text{if }i=0,7`$
$`=_n\text{if }i=4.`$
Then $`M`$ is homotopy equivalent to an $`𝕊^3`$-bundle over $`𝕊^4`$ if and only if its linking form is equivalent to a standard form for some choice of orientation on $`M`$.
The previous corollary can be strengthened in some cases using the results of D. Wilkens. In his paper , Wilkens classified simply connected manifolds, $`M^7`$, with integral cohomology as above up to the addition of homotopy 7-spheres (and hence up to $`PL`$-homeomorphism type). From the discussion in the Appendix, we have:
###### Theorem 3.
Let $`M`$ be a simply connected manifold with integral cohomology as in (1) and $`H^4(M)_n`$. If either $`n`$ is odd or $`\frac{n}{2}`$ is odd, then $`M`$ is $`PL`$-homeomorphic to an $`𝕊^3`$-bundle over $`𝕊^4`$ if and only if its linking form is equivalent to a standard form.
The above result can be further strengthened to include all integers $`n`$ using results in a forthcoming paper by B. Botvinnik and C. Escher (cf. ). The seven dimensional Berger manifold is a curious space. It is described as the homogeneous space, $`M=\mathrm{Sp}(2)/\mathrm{Sp}(1)=\mathrm{SO}(5)/\mathrm{SO}(3)`$, where the embedding of $`\mathrm{Sp}(1)`$ in $`\mathrm{Sp}(2)`$ is maximal. It is an isotropy irreducible space and has the cohomology ring as in (1) with $`H^4(M)=_{10}`$ (see Section 4). It admits a normal homogeneous metric of positive sectional curvature (cf. ). In the following question was asked:
###### Question.
Does the Berger space, $`M=\mathrm{Sp}(2)/\mathrm{Sp}(1)`$, admit the structure of an $`𝕊^3`$-bundle over $`𝕊^4`$?
In Section 4 it is shown that the linking form for the Berger space is equivalent to a standard form. So applying Corollary 2 we see that the Berger space is indeed homotopy equivalent to an $`𝕊^3`$-bundle over $`𝕊^4`$. Since $`|H^4|=10`$ for this space, we apply Theorem 3 to get,
###### Corollary 4.
The Berger space, $`M=\mathrm{Sp}(2)/\mathrm{Sp}(1)`$, is $`PL`$-homeomorphic to an $`𝕊^3`$-bundle over $`𝕊^4`$.
It remains open whether the Berger space is in fact diffeomorphic to such a bundle. This involves computing the Eells-Kuiper invariant, $`\mu `$, for this manifold (cf. ). The $`\mu `$ invariant for a 7-manifold is computed by exhibiting the space as a spin boundary; we are unable to do this for the Berger space.
Another application of Theorem 2 is the case when $`n=|H^4(X)|=p^m`$ where $`p`$ is a prime of the form $`p=4k+3`$. Since $`1`$ is not a square in the ring $`_{p^m}`$, any non-degenerate form, $`\alpha :_{p^m}_{p^m}_{p^m}`$, is equivalent to a standard form (up to sign). This shows:
###### Theorem 5.
Let $`X`$ be a Poincaré duality complex with integral cohomology as in (1) and $`H^4(X)=_{p^m}`$ where $`p`$ is a prime of the form $`p=4k+3`$. Then $`X`$ is homotopy equivalent to an $`𝕊^3`$-bundle over $`𝕊^4`$.
The proof of Theorem 1 is organized as follows: In Section 1 we show that the conditions (I) and (II) are necessary for $`X`$ to support the structure of an $`𝕊^3`$ fibration over $`𝕊^4`$. In Section 2 we establish sufficiency of the conditions and exhibit $`X`$ as the total space of an $`𝕊^3`$ fibration over $`𝕊^4`$. In Section 3 we prove that any $`𝕊^3`$ fibration over $`𝕊^4`$ is equivalent to a linear $`𝕊^3`$-bundle over $`𝕊^4`$. In Section 4 we calculate the cohomology of the Berger space and show that its linking form is equivalent to a standard form. Finally in the Appendix we discuss the results of Wilkens, Sasao and James–Whitehead that allow us to prove Theorem 3.
It is a pleasure to thank Wolfgang Ziller for many interesting and insightful discussions.
## 1. Necessity of conditions in Theorem 1
Let $`X`$ be a simply connected CW complex with integral cohomology as in (1). It follows that the 4-skeleton of $`X`$ can be chosen to be equivalent to the space $`P^4(n)`$ defined as the cofiber of the self map of degree $`n`$ on $`𝕊^3`$. The space $`X`$ is then equivalent to the cofiber of some map $`f`$ given by,
$$\begin{array}{ccccc}𝕊^6& \stackrel{f}{}& P^4(n)& & X\end{array}$$
Assume now that $`X`$ supports the structure of the total space in a fibration:
(2)
$$\begin{array}{ccccc}𝕊^3& & X& & 𝕊^4\end{array}$$
Consider the commutative diagram:
where $`\pi :X𝕊^4`$ is the projection map and $`g=\pi \text{ }\text{ }i`$. An easy argument using the Serre spectral sequence for the fibration (2) shows that $`g^{}:H^4(𝕊^4)H^4(P^4(n))`$ is an epimorphism. From the theory of secondary cohomology operations (cf. ) it follows that the secondary operation $`\mathrm{\Theta }`$ is trivial where,
$$\mathrm{\Theta }:H^4(X,𝔽_2)H^7(X,𝔽_2)$$
corresponds to the relation $`\mathrm{Sq}^2\mathrm{Sq}^2=\mathrm{Sq}^3\mathrm{Sq}^1`$ in the mod 2 Steenrod algebra.
It remains to check condition (II) in the theorem. Consider the Serre spectral sequence for the fibration (2) converging to $`H^{}(X)`$. We have:
(3) $`E_2^{p,q}`$ $`=E_4^{p,q}=H^p(𝕊^4)H^q(𝕊^3)`$
$`d_4(y_3)=ny_4`$
where $`y_3`$ and $`y_4`$ are suitably chosen generators of $`H^3(𝕊^3)`$ and $`H^4(𝕊^4)`$ respectively. Similarly, for the spectral sequence in $`_n`$-coefficients converging to $`H^{}(X,_n)`$, we have:
$$E_2^{p,q}=E_{\mathrm{}}^{p,q}=H^p(𝕊^4,_n)H^q(𝕊^3,_n)$$
Note that in both spectral sequences there are no extension problems, hence we may identify $`H^{}(X)`$ or $`H^{}(X,_n)`$ with the fifth stage in the respective spectral sequences.
Let $`[X]H_7(X)`$ be the orientation class defined by
(4)
$$y_4y_3,[X]=1$$
It follows from (3) that the Bockstein isomorphism, $`\beta :H^3(X,_n)H^4(X)`$ is given by
(5)
$$\beta ([y_3])=y_4$$
where we henceforth use the notation $`[y]`$ to denote the $`mod(n)`$ reduction of an integral class $`y`$. Using (4) and (5), we get
$$\beta ^1(y_4),y_4[X]=[y_3],y_4[X]=[y_3][y_4],[X]1mod(n)$$
which is simply the statement that the linking form is equivalent to a standard form.
## 2. Construction of a $`𝕊^3`$ fibration
The purpose of this section is to show that any CW complex $`X`$ satisfying conditions (I) and (II) of Theorem 1 is equivalent to the total space of a fibration with base $`𝕊^4`$ and fiber equivalent to $`𝕊^3`$.
Given such a complex, recall that $`X`$ fits into a cofiber sequence
$$\begin{array}{ccccc}𝕊^6& \stackrel{f}{}& P^4(n)& & X\end{array}$$
Let $`p:P^4(n)𝕊^4`$ be any map inducing an epimorphism in cohomology. Such a map always exists since $`P^4(n)`$ is a four dimensional complex. Condition (I) ensures that the composite, $`p\text{ }\text{ }f`$, is null homotopic. Thus we get an extension $`\stackrel{~}{\pi }`$ making the following diagram commute.
We can further extend the above diagram to:
(6)
where the upper horizontal map has degree 1 and $`[n]`$ denotes the self map of degree $`n`$. Now the extension $`\stackrel{~}{\pi }`$ is not unique; the set of extensions admits a transitive action of the group $`\pi _7(𝕊^4)`$, which we describe below. Given $`\alpha \pi _7(𝕊^4)`$ we define $`\alpha \stackrel{~}{\pi }`$ by
In terms of the diagram (6) it is not hard to verfiy that
(7)
$$\mathrm{H}(f(\alpha \stackrel{~}{\pi }))=n^2\mathrm{H}(\alpha )+\mathrm{H}(f(\stackrel{~}{\pi }))$$
where $`\mathrm{H}(g)`$ denotes the Hopf invariant of the map $`g`$.
Let $`F`$ be the homotopy fiber of $`\stackrel{~}{\pi }`$. We will calculate the cohomology of $`F`$ using the Serre spectral sequence for the fibration,
(8)
$$\mathrm{\Omega }𝕊^4FX$$
Recall that $`H^{}(\mathrm{\Omega }𝕊^4)=z_{3k},k=0,1,2,\mathrm{}`$. Consider the diagram of fibrations:
(9) $`\begin{array}{ccc}\mathrm{\Omega }𝕊^4& =& \mathrm{\Omega }𝕊^4\\ & & & & \\ F& & \\ & & & & \\ X& \stackrel{\stackrel{~}{\pi }}{}& 𝕊^4\end{array}`$
The next proposition is an easy consequence of the naturality of the Serre spectral sequence with respect to maps of fibrations.
###### Proposition 2.1.
In the Serre spectral sequence for (7), we have
$$d_4(z_{3k})=y_4z_{3k3}$$
where $`y_4H^4(X)`$ is a generator.
It follows that the classes, $`nz_{3k}E_4^{0,3k}`$, survive to the next stage. It remains to calculate $`d_7(nz_{3k})`$. Let $`G`$ be the homotopy fiber of $`f(\stackrel{~}{\pi })`$. Using (6) we get a diagram of fibrations:
(10) $`\begin{array}{ccc}\mathrm{\Omega }𝕊^4& \stackrel{\mathrm{\Omega }[n]}{}& \mathrm{\Omega }𝕊^4\\ & & & & \\ F& & G\\ & & & & \\ X& & 𝕊^7\end{array}`$
Notice that in the Serre spectral sequence for $`\mathrm{\Omega }𝕊^4G𝕊^7`$, we have the identity $`d_7(z_6)=\mathrm{H}(f(\stackrel{~}{\pi }))y_7`$ where $`y_7H^7(𝕊^7)`$ is a generator. Moreover since $`(\mathrm{\Omega }[n])^{}(z_6)=n^2z_6`$, using (10) we have,
###### Proposition 2.2.
The Hopf invariant, $`\mathrm{H}(f(\stackrel{~}{\pi }))`$, is a multiple of $`n`$ i.e. $`\mathrm{H}(f(\stackrel{~}{\pi }))=\lambda n`$, and in the Serre spectral sequence for (8), we have
$$d_7(nz_{3k})=\lambda y_7z_{3k6}$$
From Propositions 2.1 and 2.2, we deduce that
$`H^i(F)`$ $`=y_3i=3,y_3=nz_3`$
$`=_\lambda y_{7+3k}i=7+3k,k=0,1,2,\mathrm{}`$
Using the universal coefficients theorem, the homology is:
$`H_i(F)`$ $`=x_3i=3,`$
$`=_\lambda x_{6+3k}i=6+3k,k=0,1,2,\mathrm{}`$
where $`x_3`$ and $`y_3`$ are dual to each other. Now considering the Serre spectral sequence for the fibration, $`FX𝕊^4`$ converging to $`H_{}(X)`$ we have:
$`E_{p,q}^2=E_{p,q}^4`$ $`=H_p(𝕊^4)H_q(F)`$
$`d_4(x_4)`$ $`=nx_3`$
where $`x_4H_4(𝕊^4)`$ is a suitably chosen generator. Since $`H_6(X)=0`$, the map, $`d_4:E_{4,3}^4E_{0,6}^4=_\lambda `$ must be an epimorphism. Hence the class $`\lambda x_4x_3E_{4,3}^{\mathrm{}}`$ represents an orientation $`[X]H_7(X)`$.
In the dual picture for the cohomology Serre spectral sequence converging to $`H^{}(X)`$, we have:
(11) $`E_2^{p,q}=E_4^{p,q}`$ $`=H^p(𝕊^4)H^q(F)`$
$`d_4(y_3)`$ $`=ny_4`$
where $`y_4H^4(𝕊^4)`$ is the class dual to $`x_4`$. The classes $`y_3`$ and $`y_4`$ are permanent cycles in the spectral sequence converging to $`H^{}(X,_n)`$.
From the definition of $`[X]`$ we have,
(12)
$$[y_3],y_4[X]=[y_3][y_4],[X]\lambda mod(n)$$
As in the previous section, (11) and (12) imply
$$\beta ^1(y_4),y_4[X]\lambda mod(n)$$
Since the linking form is assumed to be equivalent to a standard form, it follows that $`\lambda \pm \tau ^2mod(n)`$ for some $`\tau (_n)^{}`$. Let $`m`$ be an integer so that $`m\tau ^1mod(n)`$. Define $`\pi :X𝕊^4`$ as the composite, $`\pi =[m]\text{ }\text{ }\stackrel{~}{\pi }`$. We now have a commutative diagram:
$$\begin{array}{ccc}X& & 𝕊^7\\ \pi & & f\left(\pi \right)=\left[m\right]\text{ }\text{ }f\left(\stackrel{~}{\pi }\right)& & \\ 𝕊^4& \stackrel{[n]}{}& 𝕊^4\end{array}$$
Note that $`\mathrm{H}(f(\pi ))\pm nmod(n^2)`$. Using (7) we may further assume $`\mathrm{H}(f(\pi ))=\pm n`$. It follows now from Proposition 2.2 that the homotopy fiber of $`\pi `$ is equivalent to $`𝕊^3`$ since it is a simply connected homology 3-sphere. We have therefore succeeded in writing $`X`$ as the total space of an $`𝕊^3`$ fibration over $`𝕊^4`$ as required.
## 3. From fibrations to bundles
In this section we show that any $`𝕊^3`$ fibration over $`𝕊^4`$ is equivalent to an $`𝕊^3`$-bundle over $`𝕊^4`$. The argument is fairly standard, but we outline it for completeness.
Let $`\xi `$ be an $`𝕊^3`$ fibration over $`𝕊^4`$. Restricting $`\xi `$ to the hemispheres, $`D_+`$ and $`D_{}`$, we get trivial fibrations, $`\xi _+`$ and $`\xi _{}`$, respectively. The fibration $`\xi `$ may then be obtained by gluing $`\xi _+`$ and $`\xi _{}`$ along their common boundary, $`𝕊^3`$, by a map lying in a homotopy class,
$$𝕊^3SG(3)$$
where $`SG(3)`$ denotes the monoid of orientation preserving self maps of $`𝕊^3`$.
Recall that (linear) $`𝕊^3`$ bundles are classified by homotopy classes of maps, $`𝕊^3\mathrm{SO}(4)`$. The classical $`J`$ homomorphism identifies $`\mathrm{SO}(4)`$ with the submonoid of $`SG(3)`$ of linear actions. Therefore to show that any $`𝕊^3`$ fibration over $`𝕊^4`$ is equivalent to an $`𝕊^3`$-bundle over $`𝕊^4`$, it suffices to show that the map,
$$J_{}:\pi _3(\mathrm{SO}(4))\pi _3(SG(3))$$
is an epimorphism.
For a fixed basepoint of $`𝕊^3`$, one has an evaluation map that evaluates the effect of a self map of $`𝕊^3`$ on the basepoint. It is easy to see that this map has a section and hence it follows that we have a map of short exact sequences:
$$\begin{array}{ccccccccc}0& & \pi _3(\mathrm{SO}(3))& & \pi _3(\mathrm{SO}(4))& \stackrel{\text{ev}_{}}{}& \pi _3(𝕊^3)& & 0\\ & & J_{}& & J_{}& & & & \\ 0& & \pi _3(SG_{}(3))& & \pi _3(SG(3))& \stackrel{\text{ev}_{}}{}& \pi _3(𝕊^3)& & 0\end{array}$$
where $`SG_{}(3)`$ are basepoint preserving elements of $`SG(3)`$.
It suffices to show that $`J_{}:\pi _3(\mathrm{SO}(3))\pi _3(SG_{}(3))`$ is an epimorphism. Then the result will follow by the 5-lemma. Consider the stabilization of $`J`$:
$$\begin{array}{ccc}\pi _3(\mathrm{SO}(3))& \stackrel{J_{}}{}& \pi _3(SG_{}(3))\\ \times 2& & \times 2& & \\ \pi _3(\mathrm{SO})& \stackrel{J_{}^s}{}& \pi _3^s(S^0)\end{array}$$
It is well known that $`\pi _3(SO(3))==\pi _3(SO)`$. Furthermore, it is also known that $`\pi _3(SG_{}(3))=_{12}`$ and $`\pi _3^s(S^0)=_{24}`$. By , $`J_{}^s`$ is an epimorphism and hence $`J_{}`$ is an epimorphism as well. This completes the proof of Theorem 1.
## 4. The Berger space
We briefly recall the construction of the Berger space. Consider $`^5`$ represented as the space of $`3\times 3`$ traceless, symmetric matrices. Then the conjugation action of $`\mathrm{SO}(3)`$ on this space affords a (maximal) representation into $`\mathrm{SO}(5)`$. The quotient space, $`M^7=\mathrm{SO}(5)/\mathrm{SO}(3)`$ may also be written as $`\mathrm{Sp}(2)/\mathrm{Sp}(1)`$ for a maximal embedding of $`\mathrm{Sp}(1)`$ into $`\mathrm{Sp}(2)`$. Berger showed in that this space admits a normal homogeneous metric of positive sectional curvature. In it was shown that this space cannot be a principal $`𝕊^3`$-bundle over $`𝕊^4`$. We shall address the question of whether it is equivalent to an $`𝕊^3`$ fiber bundle over $`𝕊^4`$.
The cohomology of this space is well known. However, we outline the calculation as we will need the setup to compute the linking form. In terms of the standard maximal tori we have a commutative diagram:
(13) $`\begin{array}{ccc}\mathrm{Sp}(1)& \stackrel{\psi }{}& \mathrm{Sp}(2)\\ & & & & \\ S^1& \stackrel{\psi _{1,3}}{}& S^1\times S^1\end{array}`$
where $`\psi _{1,m}(z)=(z,z^m)`$.
Let $`B\psi :B_{\mathrm{Sp}(1)}B_{\mathrm{Sp}(2)}`$ be the map on the level of classifying spaces. It follows from (13) that in cohomology we have:
(14) $`B\psi ^{}(p_1)`$ $`=10p_1`$
$`B\psi ^{}(p_2)`$ $`=9p_1^2`$
where $`H^{}(B_{\mathrm{Sp}(2)})=[p_1,p_2]`$ and $`H^{}(B_{\mathrm{Sp}(1)})=[p_1]`$.
The homogeneous space, $`M=\mathrm{Sp}(2)/\psi (\mathrm{Sp}(1))`$ is the concrete description of the Berger space. To calculate its cohomology, consider the fibration
(15)
$$\mathrm{Sp}(2)MB_{\mathrm{Sp}(1)}$$
We have a pullback diagram:
$$\begin{array}{ccc}\mathrm{Sp}(2)& =& \mathrm{Sp}(2)\\ & & & & \\ M& & E_{\mathrm{Sp}(2)}\\ & & & & \\ B_{\mathrm{Sp}(1)}& \stackrel{B\psi }{}& B_{\mathrm{Sp}(2)}\end{array}$$
Recall that $`H^{}(\mathrm{Sp}(2))=\text{E}(y_3,y_7)`$, where $`y_3`$ and $`y_7`$ transgress to $`p_1`$ and $`p_2`$ respectively in the Serre spectral sequence for the universal fibration. Using (14) and the pullback diagram above, we have:
###### Proposition 4.1.
In the Serre spectral sequence for (15) converging to $`H^{}(M)`$, we have,
$`d_4(y_3)`$ $`=10p_1`$
$`d_8(y_7)`$ $`=9p_1^2`$
It follows immediately from Proposition 4.1 that
(16) $`H^i(M)`$ $`=i=0,7`$
$`=_{10}i=4`$
We would like to know whether $`M`$ is homotopy equivalent to an $`𝕊^3`$-bundle over $`𝕊^4`$. By Corollary 2, it suffices to compute the linking form for $`M`$.
Let $`S_{1,3}\mathrm{Sp}(2)`$ denote the $`\psi `$-image of the standard maximal torus in $`\mathrm{Sp}(1)`$. We then have a fibration:
(17)
$$𝕊^2\mathrm{Sp}(2)/S_{1,3}M$$
An easy spectral sequence argument then yields
###### Proposition 4.2.
$`H^{}(M)`$ maps isomorphically onto $`H^{}(\mathrm{Sp}(2)/S_{1,3})`$ in degrees 0,4 and 7. The corresponding maps in homology are isomorphisms as well in degrees 0,3 and 7.
Fix an orientation $`[M]H_7(M)`$. We identify $`[M]`$ with a class $`[M]`$ in $`H_7(\mathrm{Sp}(2)/S_{1,3})`$ using Proposition 4.2. It is clear that the linking form on $`M`$ is equivalent to the form,
$`\alpha :`$ $`H^4(\mathrm{Sp}(2)/S_{1,3})H^4(\mathrm{Sp}(2)/S_{1,3})_{10}`$
$`xy\beta ^1(x),y[M]`$
It will be easier to calculate $`\alpha `$ on $`\mathrm{Sp}(2)/S_{1,3}`$ than the linking form on $`M`$.
Let $`\mathrm{\Delta }=\mathrm{Sp}(1)\times \mathrm{Sp}(1)\mathrm{Sp}(2)`$, be the standard diagonal embedding of $`\mathrm{Sp}(1)\times \mathrm{Sp}(1)`$. Consider the fibration:
(18)
$$\mathrm{\Delta }/S_{1,3}\mathrm{Sp}(2)/S_{1,3}\mathrm{Sp}(2)/\mathrm{\Delta }$$
The homogeneous spaces $`\mathrm{\Delta }/S_{1,3}`$ and $`\mathrm{Sp}(2)/\mathrm{\Delta }`$ can be identified with the spaces $`𝕊^2\times 𝕊^3`$ and $`𝕊^4`$ respectively. Hence their cohomologies have the structure of exterior algebras:
$$H^{}(\mathrm{\Delta }/S_{1,3})=\text{E}(y_2,y_3),H^{}(\mathrm{Sp}(2)/\mathrm{\Delta })=\text{E}(y_4),$$
where $`y_4`$ is chosen so that in the Serre spectral sequence for (18), we have
(19)
$$d_4(y_3)=10y_4$$
From Proposition 4.2, this is the only non-trivial differential. Since there are no extension problems, we identify $`H^{}(\mathrm{Sp}(2)/S_{1,3})`$ with the $`E_{\mathrm{}}`$ term. We can do the same for the Serre spectral sequence in $`_{10}`$ coefficients, converging to $`H^{}(\mathrm{Sp}(2)/S_{1,3},_{10})`$. Note that $`y_4y_3H^7(\mathrm{Sp}(2)/S_{1,3})`$ is a generator and hence,
(20)
$$[y_3],y_4[M]=y_4y_3,[M]\pm 1mod(10)$$
As in Section 1, we deduce from (19) and (20) that $`\alpha (y_4,y_4)\pm 1mod(10)`$. So if $`[M]`$ is chosen suitably, then the linking form for the Berger space is equivalent to a standard form.
## Appendix A Wilkens’ $`\beta `$ invariant for $`𝕊^3`$-bundles over $`𝕊^4`$
Equivalence classes of $`𝕊^3`$-bundles over $`𝕊^4`$ with structure group $`\mathrm{SO}(4)`$ are in one–one correspondence with $`\pi _3(\mathrm{SO}(4))`$. One can construct generators $`\rho ,\sigma \pi _3(\mathrm{SO}(4))`$ as follows:
$$\rho (u)v=uvu^1,\sigma (u)=uv$$
where $`u,v`$ represent quaternions of norm 1. We now adopt the convention of . With the above choice of generators, the pair $`(m,n)`$ will represent the bundle $`\xi _{m,n}`$ corresponding to $`m\rho +n\sigma `$. The total space of the associated $`𝕊^3`$-bundle over $`𝕊^4`$ will be denoted by $`M_{m,n}`$. The cohomology of $`M_{m,n}`$ can be computed without too much difficulty using the Serre spectral sequence (cf. );
$`H^0(M_{m,n})H^7(M_{m,n})=`$
$`H^4(M_{m,n})=_nx_4`$
where $`x_4`$ is the pullback of the generator of $`H^4(𝕊^4)`$.
In , D. Wilkens studied the class of 2-connected, 7-manifolds. For such manifolds he considered the set of invariants $`(H^4(M)^{},b,\beta )`$. Here $`G^{}`$ refers to the torsion part of the group $`G`$, $`b`$ is the linking form and $`\beta H^4(M)`$ is the spin characteristic class of the tangent bundle such that $`2\beta =p_1`$.
Since we are interested in $`𝕊^3`$-bundles over $`𝕊^4`$, we restrict ourselves to the case when $`H^4(M)=_n`$, a finite, cyclic group. Wilkens showed that if $`n`$ is odd, then the data $`(H^4(M),b,\beta )`$ uniquely classifies the manifold $`M`$ up to oriented $`PL`$-homeomorphism type. When $`n`$ is even, there are at most two inequivalent manifolds with the same data. For the manifolds $`M_{m,n}`$ we know that $`H^4(M_{m,n})=_nx_4`$ and the linking form is standard. Hence the Wilkens invariants for $`M_{m,n}`$ are $`(H^4(M_{m,n}),\beta )`$. From it follows that the characteristic class $`\beta `$ for $`M_{m,n}`$ is $`\pm 2mx_4`$. This is because the tangent bundle $`TM_{m,n}`$ is stably equivalent to the pullback of the vector bundle corresponding to $`\xi _{m,n}`$.
In , , James and Whitehead have shown that $`M_{m,n}`$ is oriented homotopy equivalent to $`M_{m^{},n}`$ if and only if $`m\pm m^{}mod(n,12)`$ (see also ). When $`n=2\text{odd}`$, the above condition implies that the manifolds $`M_{m,n}`$ and $`M_{m+\frac{n}{2},n}`$ represent distinct oriented homotopy types. Since they have the same Wilkens data, they realize the two distinct (oriented) $`PL`$-homeomorphism types suggested by \[13, Theorem 2\]. Hence in the case $`\frac{n}{2}`$ is an odd integer, the two distinct possibilities for the data $`(H^4,b,\beta )`$ are realized by $`𝕊^3`$-bundles over $`𝕊^4`$.
###### Proof of Theorem 3.
: Let $`M`$ be a simply connected 7-manifold with integral cohomology as in (1) and linking form equivalent to a standard form. Pick an orientation on $`M`$ and an isomorphism, $`\psi :_nH^4`$ so that $`b(\psi (a),\psi (b))=ab`$. Let $`x=\psi (1)`$; then $`b(rx,sx)=rs`$.
For any manifold with the above cohomology, $`\beta mod(2)=w_4`$, the fourth Stiefel–Whitney class of the tangent bundle. An easy calculation using the Wu classes shows that for such manifolds, the total Stiefel–Whitney class is trivial and in particular $`w_4=0`$. Hence, $`\beta =2mx`$, an even class. So the Wilkens data for $`M`$ is equivalent to the Wilkens data for the manifold $`M_{m,n}`$. From the previous discussion, it follows that when $`n`$ is odd or when $`n`$ is twice an odd number, $`M`$ is $`PL`$-homeomorphic to $`M_{m,n}`$.∎
###### Remark.
Given a manifold that is $`PL`$-homeomorphic to an $`𝕊^3`$-bundle over $`𝕊^4`$, the obstruction to diffeomorphism is measured by the $`\mu `$ invariant of Eells and Kuiper (cf. ). Computing this invariant requires exhibiting the manifold in question as the boundary of an eight dimensional spin manifold.
###### Remark.
It seems reasonable to expect that the methods outlined here can be used to prove the analog of Theorem 1 for $`𝕊^7`$-bundles over $`𝕊^8`$.
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# Nonaxial octupole deformations in light 𝑁=𝑍 nuclei at high spins
## 1 Introduction
One of the most important concepts in the many–body theory of finite Fermi systems is the mean field approach. In fact, many phenomena observed in nuclei can be explained by means of a spontaneous symmetry breaking mechanism which leads to a mean field solution that does not obey symmetries of the original many–body Hamiltonian . A nucleus with a partially filled highest shell spontaneously deforms in such a way that a lower energy minimum is achieved, and thus a non–spherical equilibrium shape is attained. Coriolis forces may cause a similar effect for a nucleus which has a spherical shape in a ground state. In particular, a strong coupling of normal and intruder states near the Fermi surface at large rotational frequencies can lead to a superdeformed or to an octupole deformed shape for certain combinations of protons and neutrons . The most practical method for the analysis of nuclear shapes is a phenomenological macroscopic–microscopic method (MMM) which combines the liquid drop model describing the bulk properties of nuclear matter and the Strutinsky Shell Correction method providing the description of quantum shell effects of phenomenological potentials . The more fundamental approach is based on self–consistent Hartree-Fock (HF) calculations once a particular choice of nucleon-nucleon interaction has been made. A commonly used inter-nucleon force is that of Skyrme. While the HF+Skyrme approach describes major features of nuclei quite well, it does not completely take into account the pairing correlations. In addition, various sets of parameters for Skyrme forces may not provide a definite answer in some cases. The Hartree-Fock-Bogoluibov (HFB) method with the Gogny forces resolves these problems quite effectively. Furthermore, it has the same prediction power as the MMM .
In recent years a lot of experimental and theoretical efforts has been devoted to the analysis of superdeformed (SD) bands in different mass region . On the other hand, the study of octupole degrees of freedom is also a topical subject in the nuclear structure physics . It turns out that octupole deformations are significant for superdeformed nuclei as well as for normal-deformed (ND) nuclei. Most of these studies were restricted to the axial octupole deformation proportional to the $`Y_{30}`$ term in the octupole family. The reason for this is primarily computational complications arising as a result of the extra degree of freedom introduced by non-axial octupole deformations. However, during the last few years some remarkable results have been reported. It was found that $`Y_{31}`$ and $`Y_{32}`$ components resulted in the lowest energy octupole vibrations for oblate and prolate superdeformed shapes and therefore may have consequences for octupole instability . Calculations using the MMM with the Woods–Saxon potential predicted the importance of the banana-type $`Y_{31}`$ deformation for highly deformed nuclei , while the $`Y_{32}`$ deformation has been found to be important in the $`{}_{}{}^{222}Ra`$ nucleus . The self-consistent HF+Skyrme calculations suggest the softness of the oblate states in $`A80`$ nuclei against the $`Y_{33}`$ deformation in the ground states. The self-consistent cranking HF+Skyrme approach predicts that the $`Y_{31}`$ deformation is important for a correct description of the yrast band of $`{}_{}{}^{32}S`$. Moreover, the study of octupole deformations sheds light on the interplay between a classical chaos and a quantum spectrum of finite Fermi systems . Guided by thorough investigation of non-axial octupole deformations in a harmonic oscillator model , we aim to analyze how symmetries break at high spins in the cranking HFB approach with the Gogny interaction. In accordance with the octupole instability suggested at particle number $`N=`$16 and 28 , we choose two $`N=Z`$ nuclei, $`{}_{}{}^{32}S`$ and $`{}_{}{}^{56}Ni`$. In Sec.2 we review the main features of our model. The discussion of main results is presented in Sec.3 with the conclusion following in Sec.4. In the Appendix we introduce a simple two-level model to understand qualitatively the behavior of the Inglis-Beliaev moment of inertia due to strong octupole coupling between two single-particle (s.p.) states.
## 2 The Model
As mentioned above, the present self-consistent cranking HFB calculation has been performed with the effective Gogny D1S interaction . This interaction provides a good description of many nuclear properties over the nuclide chart, e.g., ground state energies, odd-even energy differences, electron scattering data , fission barriers . Also, a good description is obtained for bulk properties of rotating nuclei in actinide region , mercury region and fp-shell region . Our numerical code has been used for the analysis of microscopic dynamics in light rotating nuclei . In these calculations, however, the signature symmetry was conserved (see discussion about different symmetries in rotating nuclei in ). The present calculations are performed without assuming á priori the axial and signature symmetries. In addition, our Hamiltonian includes the Coulomb interaction and the center of mass correction up to the exchange terms exactly.
The $`z`$-axis is taken as a rotational axis in our code. To save the CPU time in numerical calculations, we impose the $`\widehat{P}e^{i\pi \widehat{J}_z}`$ ($`z`$-simplex) and the $`\widehat{P}e^{i\pi \widehat{J}_y}\widehat{\tau }`$ ($`\widehat{S}_y^T`$) symmetries , where $`\widehat{P}`$ is the parity operator, $`e^{i\pi \widehat{J}_i}`$ is the rotation operator around the $`i`$-axis (i=y,z) on angle $`\pi `$ , and $`\widehat{\tau }`$ is the time-reversal operator. Due to the $`z`$-simplex and $`\widehat{S}_y^T`$ symmetries, the mass asymmetry of nucleus is allowed only along the $`x`$-axis. Thus, we solve numerically the following cranking HFB equations:
$`\delta \varphi \left(\omega \right)|\widehat{H}\lambda _p\widehat{Z}\lambda _n\widehat{N}\omega \widehat{J}_z`$
$`+\mu _x\varphi \left(\omega \right)|\widehat{x}|\varphi \left(\omega \right)\widehat{x}|\varphi \left(\omega \right)`$ $`=`$ $`0.`$ (1a)
$`\varphi \left(\omega \right)|\widehat{Z}|\varphi \left(\omega \right)=Z,\varphi \left(\omega \right)|\widehat{N}|\varphi \left(\omega \right)`$ $`=`$ $`N,`$ (1b)
$`\varphi \left(\omega \right)|\widehat{J}_z|\varphi \left(\omega \right)`$ $`=`$ $`I,`$ (1c)
$`\varphi \left(\omega \right)|\widehat{x}|\varphi \left(\omega \right)`$ $`=`$ $`0,`$ (1d)
where the Lagrange multipliers $`\lambda _p`$ and $`\lambda _n`$ are the chemical potentials of proton and neutron fields, respectively (the operators $`\widehat{Z}`$ and $`\widehat{N}`$ are protons and neutrons number operators); the Lagrange multiplier $`\omega `$ is the angular frequency of a collective rotation around the $`z`$-axis and $`\widehat{J}_z`$ is the $`z`$-component of the angular momentum operator $`\widehat{𝑱}`$. To keep the center of mass motion fixed, we also impose the quadrupole constraint operator $`\mu _x\varphi \left(\omega \right)|\widehat{x}|\varphi \left(\omega \right)\widehat{x}`$ to the Routhian $`\widehat{R}=\widehat{H}\lambda _p\widehat{Z}\lambda _n\widehat{N}\omega \widehat{J}_z`$ in the $`x`$-axis direction.
The quasi-particle (q.p.) operators
$$\widehat{\alpha }_i^{}=\underset{k}{}U_{ki}\widehat{c}_k^{}+V_{ki}\widehat{c}_k,$$
(2)
where the state $`|k`$ is a single-particle basis state (see below), obey the equation of motion
$$[\widehat{R},\widehat{\alpha }_i^{}]=ϵ_i(\omega )\widehat{\alpha }_i^{}$$
(3)
which defines the quasi-particle energies $`ϵ_i(\omega )`$ and quasi-particle amplitudes $`U_{ki}`$ and $`V_{ki}`$ as functions of the rotational frequency $`\omega `$.
We define the $`y`$-axis as the quantization axis of deformation for convenience. Consequently, the quadrupole, $`\beta _2`$ and $`\gamma `$, and octupole , $`\beta _{3m}`$, deformation parameters are defined as
$`\beta _2\mathrm{cos}\gamma `$ $``$ $`{\displaystyle \frac{4\pi }{5}}{\displaystyle \frac{r^2Y_{20}(\theta ,\phi )}{AR_0^2}}`$
$`=`$ $`\sqrt{{\displaystyle \frac{\pi }{5}}}{\displaystyle \frac{(3y^2r^2)}{AR_0^2}},`$
$`\beta _2\mathrm{sin}\gamma `$ $``$ $`{\displaystyle \frac{4\pi }{5}}{\displaystyle \frac{r^2(Y_{22}(\theta ,\phi )+Y_{22}(\theta ,\phi ))}{\sqrt{2}AR_0^2}}`$
$`=`$ $`\sqrt{{\displaystyle \frac{3\pi }{5}}}{\displaystyle \frac{(x^2z^2)}{AR_0^2}},`$
$`\beta _{30}`$ $``$ $`{\displaystyle \frac{4\pi }{3}}{\displaystyle \frac{r^3Y_{30}(\theta ,\phi )}{AR_0^3}}`$
$`=`$ $`\sqrt{{\displaystyle \frac{7\pi }{9}}}{\displaystyle \frac{y(5y^23r^2)}{AR_0^3}},`$
$`\beta _{31}`$ $``$ $`{\displaystyle \frac{4\pi }{3}}{\displaystyle \frac{r^3(Y_{31}(\theta ,\phi )Y_{31}(\theta ,\phi ))}{\sqrt{2}AR_0^3}}`$
$`=`$ $`\sqrt{{\displaystyle \frac{21\pi }{18}}}{\displaystyle \frac{x(5y^2r^2)}{AR_0^3}},`$
$`\beta _{32}`$ $``$ $`{\displaystyle \frac{4\pi }{3}}{\displaystyle \frac{r^3(Y_{32}(\theta ,\phi )+Y_{32}(\theta ,\phi ))}{\sqrt{2}AR_0^3}}`$
$`=`$ $`\sqrt{{\displaystyle \frac{105\pi }{9}}}{\displaystyle \frac{y(x^2z^2)}{AR_0^3}},`$
$`\beta _{33}`$ $``$ $`{\displaystyle \frac{4\pi }{3}}{\displaystyle \frac{r^3(Y_{33}(\theta ,\phi )Y_{33}(\theta ,\phi ))}{\sqrt{2}AR_0^3}}`$
$`=`$ $`\sqrt{{\displaystyle \frac{35\pi }{18}}}{\displaystyle \frac{x(x^23z^2)}{AR_0^3}},`$
where $`(r,\theta ,\phi )`$ are spherical coordinates related to the Cartesian coordinates in the rotating frame $`(x,y,z)`$ as
$$(x,y,z)(r\mathrm{sin}\theta \mathrm{cos}\phi ,r\mathrm{cos}\theta ,r\mathrm{sin}\theta \mathrm{sin}\phi ).$$
(10)
All deformations proportional to the odd power of the $`y`$ coordinate are forbidden in our code due to the $`\widehat{S}_y^T`$ symmetry. Since $`\beta _{30}`$ and $`\beta _{32}`$ are proportional to the odd $`y^n`$ terms, we use $`\beta _{31}`$ and $`\beta _{33}`$ to represent the degree of the non-axial octupole deformation, when the $`y`$-axis is the largest axis of a prolate nucleus. Notice that the octupole deformation with a mass asymmetry like a pear shape is also represented by the combination of $`\beta _{31}`$ and $`\beta _{33}`$ when the $`x`$-axis is the largest axis of a prolate nucleus <sup>1</sup><sup>1</sup>1From these $`\beta _{31}`$ and $`\beta _{33}`$ values, for an example, we can obtain the value of $`\beta _{30}`$ for the $`x`$-quantization axis as
$$\beta _{30}^x=\frac{\sqrt{6}}{4}\beta _{31}\frac{\sqrt{10}}{4}\beta _{33}.$$
(11) .
The s.p. wave functions have been expanded in a three-dimensional harmonic oscillator basis up to the principal quantum number $`N=8`$ for $`{}_{}{}^{32}S`$ and up to $`N=10`$ for $`{}_{}{}^{56}Ni`$. The basis has been symmetrized with respect to the $`z`$-simplex operation, and eigenfunctions are eigenstates of the $`\widehat{S}_y^T`$ operator (the Goodman basis ). Since the ground state shapes of the chosen nuclei are a normal-deformed one for $`{}_{}{}^{32}S`$ and a spherical one for $`{}_{}{}^{56}Ni`$, we use a spherical Cartesian basis with the same range parameters of the Hermite polynomials for all axes. The range parameters have been optimized to reproduce the largest binding energy of each ground state.
To understand the relation between s.p. degrees of freedom and collective effects due to the rotation, we calculate three moments of inertia, the static moment of inertia $`𝒥^{(1)}\left(\omega \right)I/\omega `$, the dynamical moment of inertia $`𝒥^{(2)}\left(\omega \right)dI/d\omega `$, and the Inglis-Beliaev (IB) moment of inertia
$$𝒥_{\mathrm{IB}}\left(\omega \right)=2\underset{i>j}{}\frac{|J_{ij}\left(\omega \right)|^2}{ϵ_i\left(\omega \right)+ϵ_j\left(\omega \right)},$$
(12)
which is the leading order term of the nuclear moment of inertia. Here $`J_{ij}\left(\omega \right)`$ is a matrix element of the angular momentum operator $`\widehat{J}_z`$
$$J_{ij}\left(\omega \right)=\varphi \left(\omega \right)|[\widehat{\alpha }_j\widehat{\alpha }_i,\widehat{J}_z]|\varphi \left(\omega \right).$$
(13)
It is well known that the IB moment of inertia is too small to reproduce the absolute value of the moment of inertia. Comparing $`𝒥_{\mathrm{IB}}\left(\omega \right)`$ and $`𝒥^{(1)}\left(\omega \right)`$ or $`𝒥^{(2)}\left(\omega \right)`$, we can find the validity of using the $`𝒥_{\mathrm{IB}}\left(\omega \right)`$ for the analysis of structural changes in the nuclear mean field. In addition, we analyze different contributions to the IB moment of inertia, Eq.(12),
$$𝒥_{\mathrm{IB}}\left(\omega \right)=\underset{\tau ,s}{}𝒥_{\tau ,s}\left(\omega \right)𝒥_{p,+1}\left(\omega \right)+𝒥_{p,1}\left(\omega \right)+𝒥_{n,+1}\left(\omega \right)+𝒥_{n,1}\left(\omega \right),$$
(14)
where $`(\tau ,s)`$ denote the isospin $`\tau =p,n`$ and the $`z`$-simplex quantum number $`s=\pm 1`$ which characterize different subspaces. The comparison of different components of $`𝒥_{\mathrm{IB}}\left(\omega \right)`$ could provide the information in which subspace the q.p. degrees of freedom mainly affect the bulk properties of the nuclei.
The cranking HFB equations (2) are solved in an iterative way. As the convergence condition for each HFB state, we impose the condition
$$\underset{i}{}\left|ϵ_i^{(n)}\left(\omega \right)ϵ_i^{(n1)}\left(\omega \right)\right|.1\text{[KeV]},$$
(15)
where $`n`$ is the number of iterations in the course of solving the cranking HFB equations.
## 3 Discussion of results
### 3.1 $`{}_{}{}^{32}S`$
In Figs.1-3 results for the yrast line states, angular momenta and deformations as functions of the rotational frequency are presented. The binding energy of the ground state in $`{}_{}{}^{32}S`$ is well reproduced in our calculations (Fig.1). We found that the contribution from the pairing interaction terms to the total binding energy is almost negligible. The energy gaps between q.p. energies near the Fermi energies are large, about 4 MeV, both in the proton and neutron fields. Therefore, the calculations for the ND and SD bands have been done using the cranking HF approximation with the same parameter set as for the HFB calculations. While there are differences in details between our calculations and those with Skyrme forces , we obtained similar main results for the ND and SD bands. The SD band becomes the yrast one for $`I12\mathrm{}`$. The excitation energy of the SD minimum relative to the ground state is about 10 MeV (Fig. 1). Except for the binding energy, these results are also consistent with those of SLy4-HF calculations , in which $`\beta `$-deformations are $`\beta =.16`$ and $`\beta =.7`$ in the ground state and SD minimum state, respectively. In comparison with the results of the Skyrme III calculations , however, finite values of the non-axial octupole deformation $`Y_{31}`$ in the ND yrast band are obtained at larger rotational frequency $`\omega 1.5`$ \[MeV/ $`\mathrm{}`$\]. As is seen in Fig.4(a), the value of the non-axial octupole deformation $`|\beta _{31}|`$ increases suddenly at $`\omega 1.5`$ \[MeV/ $`\mathrm{}`$\], where the pseudo level-crossing occurs between s.p. orbits in the subspaces with $`s=\pm 1`$ (Fig.5(a)), both in the proton and neutron fields. These orbits can be associated with the principal quantum numbers $`N=2`$ and $`N=3`$ at $`\omega =0`$\[MeV/$`\mathrm{}`$\]. As a result, the $`Y_{31}`$ deformation becomes favorable at $`\omega 2.2`$ \[MeV/ $`\mathrm{}`$\], for $`{}_{}{}^{32}S`$ with a finite value $`\beta _{31}=.13`$. The density distribution projected on the plane perpendicular to the rotation axis is shown on Fig.6(a) for the octupole band at $`\omega 2.2`$ \[MeV$`/\mathrm{}`$\] where the $`\beta _{31}`$ deformation has a maximal value. High resolution $`\gamma `$ spectroscopy could test the validity of different predictions based on different nuclear interactions.
From the analysis of Figs.4(a), 7(a) it follows that there is a correlation between the behavior of the dynamical moment of inertia $`𝒥^{(2)}\left(\omega \right)`$ and the occurrence of the octupole deformation. Similar to the behavior of the $`|\beta _{31}|`$, the dynamical moment of inertia $`𝒥^{(2)}\left(\omega \right)`$ begins to increase at $`\omega 1.5`$ \[MeV/$`\mathrm{}`$\], and both quantities show local maxima around $`\omega 2`$ \[MeV/$`\mathrm{}`$\]. Due to the strong octupole interaction the quasi-crossing between s.p. levels is smooth (Fig.5(a)). The static moment of inertia $`𝒥^{(1)}\left(\omega \right)`$ is less sensitive to the quasi-crossing of s.p. levels, and it changes smoothly with the increase of the rotational frequency $`\omega `$. However, the IB moment of inertia may reflect the structural changes in a mean field like the dynamical moment of inertia $`𝒥^{(2)}\left(\omega \right)`$. In fact, it also begins to increase at $`\omega 1.5`$ \[MeV/$`\mathrm{}`$\] and shows a local maximum at $`\omega 2.2`$ \[MeV/$`\mathrm{}`$\]. The correlation between the IB moment of inertia and the magnitude of the octupole deformation $`|\beta _{31}|`$ may be understood within a simple two-level model presented in the Appendix A. According to this model, the two rotation-aligned single-particle states with a strong octupole coupling near the Fermi-surface may give a large contribution to the value of the IB moment of inertia. From Fig.8(a) follows that the magnitudes of $`𝒥_{p,+1}\left(\omega \right)`$ and $`𝒥_{n,+1}\left(\omega \right)`$ (see Eq.(14)) are smaller than those from the $`s=1`$ subspaces which play the dominant roles in the collective rotation at low spins, in the region $`\omega <1.5`$\[MeV/$`\mathrm{}`$\]. However, at the $`\omega >1.5`$\[MeV/$`\mathrm{}`$\], $`𝒥_{p,+1}\left(\omega \right)`$ and $`𝒥_{n,+1}\left(\omega \right)`$ begin to increase rapidly, and they show local maxima at $`\omega 2.2`$\[MeV/$`\mathrm{}`$\] where $`|\beta _{31}|`$ also shows the maximum. Since $`𝒥_{p,1}\left(\omega \right)`$ and $`𝒥_{n,1}\left(\omega \right)`$ are almost unchanged in this region, the occurrence of the $`Y_{31}`$ deformation is mainly due to the s.p. quasi-crossing in the $`s=+1`$ subspaces. To justify this statement we calculate the quantity $`\mathrm{cos}^2\psi \left(\omega \right)`$ (see Appendix). Using the identity of the intruder orbits $`\mu `$ which have odd parity at $`\omega =0`$\[MeV/$`\mathrm{}`$\], we evaluate the $`\mathrm{cos}^2\psi \left(\omega \right)`$ approximately as:
$$\mathrm{cos}^2\psi \left(\omega \right)\underset{n_x+n_y+n_z=odd}{}|n_x,n_y,n_z;\sigma |\mu \left(\omega \right)|^2,$$
(16)
where a s.p. vector $`|n_x,n_y,n_z;\sigma `$ is a component of the Goodman basis in the $`s=+1`$ subspace, and $`n_x`$, $`n_y`$ and $`n_z`$ are the number of quanta on the $`x`$, $`y`$ and $`z`$ axes, respectively. As is shown in Fig.9(a), the evaluated value of $`\mathrm{cos}^2\psi \left(\omega \right)=1/2`$ or equivalently $`|\psi \left(\omega \right)|=\pi /4`$ at $`\omega 2.2`$ \[MeV$`/\mathrm{}`$\]. According to the analysis within the two-level model (see Appendix), at this value of $`|\psi \left(\omega \right)|`$ the mixing of s.p. states with opposite parities is the strongest one.
Similar to the results of , our calculations show rather shallow minima for the $`\beta _{31}`$ deformation. In Fig.10(a) we present results for the fluctuations of octupole and quadrupole deformations
$$\mathrm{\Delta }Q_{lm}\frac{\sqrt{\varphi \left(\omega \right)|\widehat{Q}_{lm}^2|\varphi \left(\omega \right)\varphi \left(\omega \right)|\widehat{Q}_{lm}|\varphi \left(\omega \right)^2}}{|\varphi \left(\omega \right)|\widehat{Q}_{lm}|\varphi \left(\omega \right)|}$$
(17)
(see also Fig.11(a)). The contribution of the ground state octupole correlations may improve the mean field results. This contribution could be estimated, for example, in the random phase approximation like according to the prescription , however, this is beyond the scope of the present paper.
As is seen in Fig.7(a), both moments of inertia, $`𝒥^{(2)}\left(\omega \right)`$ and $`𝒥_{\mathrm{IB}}\left(\omega \right)`$, begin to increase again at $`\omega 2.5`$ \[MeV$`/\mathrm{}`$\] , which is apparently due to the quasi-crossings in the $`s=1`$ subspaces at this point. There are also level-crossings between s.p. orbits with different simplex numbers (Fig.5(a)), which suggest the instability of the $`z`$-simplex symmetrized state at high spins. We may speculate that the breaking of the $`z`$-simplex symmetry due to the tilted or the chiral rotation may lead to a better description of high spin states in $`{}_{}{}^{32}S`$ at $`\omega 2.5`$ \[MeV$`/\mathrm{}`$\].
### 3.2 $`{}_{}{}^{56}Ni`$
Our calculations reproduce the binding energies and the charge radii of ground states in $`{}_{}{}^{5660}Ni`$ quite well (Fig. 12). In this chain of nuclei we concentrate our attention on $`{}_{}{}^{56}Ni`$.
Though the shape of the ground state in $`{}_{}{}^{56}Ni`$ is spherical, we also found two shape isomers with the ND and SD configuration at $`\omega =0`$ \[MeV/$`\mathrm{}`$\], which are $`5`$ MeV and $`18`$ MeV above the ground state, respectively (see Fig.13). The yrast non-rotating states observed in $`{}_{}{}^{56}Ni`$ are related to vibrational excitations. In fact, their properties can be studied within the self-consistent cranking approach with the rotation around a symmetry axis and we will discuss this subject in a forthcoming paper. The calculations for the ground state, ND and SD bands have been done in a similar manner to those of $`{}_{}{}^{32}S`$ in the HF approximation, since the contribution of the pairing interaction terms to the total binding energies are almost negligible. On Fig.2 the evolution of the angular momentum with the increase of the rotational frequency is shown for SD and ND bands. At low spins the SD minimum is formed by the 4$`p`$-4$`h`$ configuration with respect to the ground state in $`{}_{}{}^{56}Ni`$, $`\pi [(3)^2(4)^2]\nu [(3)^2(4)^2]`$, which has a prolate shape with the quadrupole deformation $`\beta =.6`$ (Fig.3). This is also 4$`h`$ configuration with respect to the SD state in $`{}_{}{}^{60}Zn`$ due to the $`N=30`$ SD gap in the $`Zn`$ isotopes , which is also seen in Fig.14. The MMM calculations also predict the SD shape at large rotational frequencies. We restricted, however, our analysis of the SD band in the region of $`\omega =0.00.4`$\[MeV/$`\mathrm{}`$\], since at larger rotational frequencies the contribution of higher shells $`(N>10)`$ becomes important.
Using the decomposition of the rational harmonic oscillator (RHO) into the isotropic ones, the relation between multi-clusters and mean field results has been discussed in . For SD shapes in the RHO the one symmetric and one asymmetric combinations of spherical oscillators are expected (see also ). Within the RHO the SD states in $`{}_{}{}^{56}Ni`$ correspond to the asymmetric combination of two spherical oscillators with magic numbers $`40`$ and $`16`$. The density distribution of the SD state (see Fig.15) does not show such a multi-cluster structure at $`\omega =0`$ \[MeV/$`\mathrm{}]`$. The possible reason is that our configuration space is too small. However, the SD minimum could be related to the resonance state in the <sup>28</sup>Si+<sup>28</sup>Si collision at high excitation energy $`E_{\mathrm{CM}}^{}=6570`$MeV and high angular momenta $`I=3442\mathrm{}`$ reported in . A thorough study of $`{}_{}{}^{56}Ni`$ as well as $`{}_{}{}^{32}S`$ at high spins may help to understand the link between the tendency for nuclei to create strongly deformed shapes and the tendency to develop the cluster structure.
A prolate deformation $`\beta =.35`$ is found for the ND minimum (Fig. 3) at $`\omega =0.0`$ \[MeV/$`\mathrm{}`$\]. With the increase of the rotational frequency the non-axial octupole deformation $`Y_{31}`$ arises in the ND band. As is seen in Fig.4(b), the value of the non-axial octupole deformation increases rapidly around $`\omega .9`$\[MeV/$`\mathrm{}`$\]. It corresponds to $`I10`$\[$`\mathrm{}`$\], where the pseudo level-crossing occurs between s.p. orbits which have the principal quantum numbers $`N=3`$ and $`N=4`$, respectively, at $`\omega =0`$\[MeV/$`\mathrm{}`$\] (Fig.5(b)). The maximal value of the $`Y_{31}`$ deformation, $`\beta _{31}=.09`$, is approached at $`\omega 1.4`$\[MeV/$`\mathrm{}`$\] (Fig.3(b)). As is seen in Fig. 8(b), $`𝒥_{p,+1}\left(\omega \right)`$ and $`𝒥_{n,+1}\left(\omega \right)`$ increase rapidly around $`\omega .9`$\[MeV/$`\mathrm{}`$\] and approach the maximal value at $`\omega 1.3`$\[MeV/$`\mathrm{}`$\]. Using the two-level model and similar arguments as for the case of $`{}_{}{}^{32}S`$ we may conclude: since $`𝒥_{p,1}\left(\omega \right)`$ and $`𝒥_{n,1}\left(\omega \right)`$ are decreasing in this region of the rotational frequency, it is most likely that the octupole deformation is caused by the quasi-crossing of s.p. levels, which have different parities at $`\omega =0`$\[MeV/$`\mathrm{}`$\], in the $`s=+1`$ subspaces. Both $`\mathrm{cos}^2\psi \left(\omega \right)`$ of the intruder orbits with the positive parity at $`\omega =0`$\[MeV/$`\mathrm{}`$\] in the proton and neutron fields,
$$\mathrm{cos}^2\psi \left(\omega \right)\underset{n_x+n_y+n_z=even}{}|n_x,n_y,n_z;\sigma |\mu \left(\omega \right)|^2,$$
(18)
have maximal values at $`\omega 1.3`$\[MeV/$`\mathrm{}`$\] and $`\omega 1.35`$\[MeV/$`\mathrm{}`$\], where the components of the IB moment of inertia, $`𝒥_{p,s=+1}\left(\omega \right)`$ and $`𝒥_{n,s=+1}\left(\omega \right)`$, show maxima, respectively (Fig. 9 (b)).
The dynamical and IB moments of inertia also increase at $`\omega .9`$ \[MeV/$`\mathrm{}`$\] and show maxima around $`\omega 1.2`$ \[MeV/$`\mathrm{}`$\] (Fig.7(b)). However, the behavior of the $`𝒥^{(2)}\left(\omega \right)`$ and $`𝒥_{\mathrm{IB}}\left(\omega \right)`$ moment of inertia is less correlated at larger values of the rotational frequency. The pseudo-crossing is much sharper in $`{}_{}{}^{56}Ni`$ in comparison with the one in $`{}_{}{}^{32}S`$ (Fig.5), i.e. the octupole interaction is expected to be weaker. Both $`\mathrm{cos}^2\psi \left(\omega \right)`$ decrease rapidly at $`\omega >1.55`$\[MeV/$`\mathrm{}`$\] (see Fig.(9(b)). Due to the termination of the parity mixing of the s.p. levels, the matrix element $`J_{ij}\left(\omega \right)`$ between mixed states, Eq.(13), becomes rapidly small. As a result, the IB moment of inertia, $`𝒥_{\mathrm{IB}}\left(\omega \right)`$, also decreases. On the other hand, the dynamical moment of inertia, $`𝒥^{(2)}\left(\omega \right)`$, reflects changes of the self-consistent mean field and, in particular, the modification of the two-body interaction due to the rotation. Due to the two-body interaction, the high intruder unoccupied state ($`s=+`$) contributes to the sudden increase of the dynamical moment of inertia observed at high rotational frequencies $`\omega >1.5[`$MeV$`/\mathrm{}]`$.
The fluctuations of octupole deformations $`\mathrm{\Delta }Q_{31}`$ and $`\mathrm{\Delta }Q_{33}`$ are always larger than 1 in all region of the ND band (Fig. 10(b)), in spite of the presence of the octupole minimum (Fig. 11(b)). The octupole minimum in $`{}_{}{}^{56}Ni`$ is much more shallow in comparison with the one of $`{}_{}{}^{32}S`$ and it could explain large fluctuations of the octupole deformations in the considered case.
## 4 Summary
It was expected from various calculations based on the MMM that octupole deformations would arise in rotating nuclei . Using the cranking HF(B) approach with the effective Gogny interaction, we found that the non-axial octupole deformation associated with the $`Y_{31}`$ term in the octupole family becomes important in the yrast band of $`{}_{}{}^{32}S`$ at the angular momenta $`I5[\mathrm{}]`$. The primary mechanism behind the occurrence of the octupole deformation is related to the strong mixing via octupole interaction of s.p. orbits with a positive simplex quantum number. Similar phenomena have been observed in the nucleus $`{}_{}{}^{56}Ni`$, where we predict the octupole softness in the excited ND band at high spins $`I10[\mathrm{}]`$.
Finally, an exploration of the octupole phenomenon certainly could deepen our understanding of different aspects of a spontaneous symmetry breaking mechanism in finite Fermi systems like nuclei. In particular, the breaking of the intrinsic reflection symmetry could be related to unexpected strong electric dipole and octupole transitons in rotational bands and to a formation of multi-cluster structures in strongly deformed nuclei. Measurements using new generations of modern detectors can test the predictions made within the HF(B) approach and lead to new insights regarding effective nucleon-nucleon interactions and their properties.
## Acknowledgements
We are thankful to F. Sakata, Y. Hashimoto and Y. Kanada-En’yo for their support in numerical calculations. We are also grateful to K. Matsuyanagi for illuminating communications. This work was supported in part by the RFFI under Grant 00-02-17194.
## Appendix A A two-level model
Let us consider the situation such that the high-lying intruder s.p. orbit $`\mu `$ is coming down to the highest-lying occupied orbit $`\nu `$ due to the Coriolis term $`\omega \widehat{J}_z`$; they have the same simplex quantum number. For the sake of simplicity, we restrict our discussion within the $`2\times 2`$ subspace. We assume that i) our states are eigenstates of the parity operator $`P`$ at $`\omega =0`$; ii) there are no changes in the subspace expanded by the simplex partners of $`\mu `$ and $`\nu `$ at $`\omega 0`$.
At $`0\omega \omega _c`$ the eigenstates can be written as
$`|\nu \left(\omega \right)`$ $`=`$ $`\mathrm{cos}\psi \left(\omega \right)|\nu \left(\omega =0\right)+\mathrm{sin}\psi \left(\omega \right)|\mu \left(\omega =0\right),`$
$`|\mu \left(\omega \right)`$ $`=`$ $`\mathrm{cos}\psi \left(\omega \right)|\mu \left(\omega =0\right)\mathrm{sin}\psi \left(\omega \right)|\nu \left(\omega =0\right),`$ (A.1)
where $`|\psi \left(\omega \right)|`$ is a monotonic function of $`\omega `$ <sup>2</sup><sup>2</sup>2A sign of $`\psi \left(\omega \right)`$ is not important for the absolute value of the octupole deformation., which satisfies the following conditions;
$`0|\psi \left(\omega \right)|{\displaystyle \frac{\pi }{2}},`$ (A.2)
$`\psi \left(\omega =0\right)=0,|\psi (\omega _c)|={\displaystyle \frac{\pi }{2}}.`$
Here, $`\omega _c`$ is the largest rotational frequency. Using the unitary transformation in Eq.(A), at $`\omega 0`$ the density matrix has the following form in this subspace
$`𝝆`$$`\left(\omega \right)`$ $``$ $`\left(\begin{array}{cc}\rho _{\nu \nu }\left(\omega \right)& \rho _{\nu \mu }\left(\omega \right)\\ \rho _{\mu \nu }\left(\omega \right)& \rho _{\mu \mu }\left(\omega \right)\end{array}\right)`$ (A.5)
$`=`$ $`\left(\begin{array}{cc}\mathrm{cos}\psi \left(\omega \right)& \mathrm{sin}\psi \left(\omega \right)\\ \mathrm{sin}\psi \left(\omega \right)& \mathrm{cos}\psi \left(\omega \right)\end{array}\right)\left(\begin{array}{cc}1& 0\\ 0& 0\end{array}\right)`$ (A.13)
$`\left(\begin{array}{cc}\mathrm{cos}\psi \left(\omega \right)& \mathrm{sin}\psi \left(\omega \right)\\ \mathrm{sin}\psi \left(\omega \right)& \mathrm{cos}\psi \left(\omega \right)\end{array}\right)`$
$`=`$ $`\left(\begin{array}{cc}\mathrm{cos}^2\psi \left(\omega \right)& \mathrm{sin}2\psi \left(\omega \right)/2\\ \mathrm{sin}2\psi \left(\omega \right)/2& \mathrm{sin}^2\psi \left(\omega \right)\end{array}\right).`$ (A.16)
At $`\omega =0`$, the matrix elements of the octupole operator $`\widehat{Q}_3`$ have the following structure
$$𝒒\left(\begin{array}{cc}q_{\nu \nu }& q_{\nu \mu }\\ q_{\mu \nu }& q_{\mu \mu }\end{array}\right)=\left(\begin{array}{cc}0& q\\ q& 0\end{array}\right),$$
(A.17)
where $`q`$ is a finite real number. At $`\omega 0`$ the expectation value of $`\widehat{Q}_3`$ is
$$|\widehat{Q}_3|=|\mathrm{tr}(𝒒𝝆\left(\omega \right))|=|q\mathrm{sin}2\psi \left(\omega \right)|,$$
due to the assumption ii).
At $`\omega =0`$ matrix elements of the angular momentum operator $`\widehat{J}_z`$ are defined as
$$𝒋\left(\begin{array}{cc}j_{\nu \nu }& j_{\nu \mu }\\ j_{\mu \nu }& j_{\mu \mu }\end{array}\right)=\left(\begin{array}{cc}j& 0\\ 0& j^{}\end{array}\right).$$
(A.18)
At $`\omega 0`$, using the unitary transformation in Eq.(A), we obtain for the matrix elements of $`𝒋`$
$`\left(\begin{array}{cc}j_{\nu \nu }\left(\omega \right)& j_{\nu \mu }\left(\omega \right)\\ j_{\mu \nu }\left(\omega \right)& j_{\mu \mu }\left(\omega \right)\end{array}\right)=`$ (A.21)
$`\left(\begin{array}{cc}j\mathrm{cos}^2\psi \left(\omega \right)+j^{}\mathrm{sin}^2\psi \left(\omega \right)& (j^{}j)\mathrm{sin}2\psi \left(\omega \right)/2\\ (j^{}j)\mathrm{sin}2\psi \left(\omega \right)/2& j^{}\mathrm{cos}^2\psi \left(\omega \right)+j\mathrm{sin}^2\psi \left(\omega \right)\end{array}\right),`$ (A.24)
From Eqs. (A) and (A.21), it follows that $`|\widehat{Q}_3|`$ and $`|j_{\mu \nu }\left(\omega \right)|^2=(j^{}j)^2\mathrm{sin}^22\psi \left(\omega \right)/4`$ have a maximum at $`|\psi \left(\omega \right)|=\pi /4`$, where there is a strong mixing of the unperturbed eigenfunctions. Therefore, the $`\mu \nu `$ component of the Inglis-Beliaev moment of inertia
$`{\displaystyle \frac{2|j_{\mu \nu }\left(\omega \right)|^2}{ϵ_\mu \left(\omega \right)+ϵ_\nu \left(\omega \right)}}={\displaystyle \frac{(j^{}j)^2\mathrm{sin}^22\psi \left(\omega \right)}{2(ϵ_\mu \left(\omega \right)+ϵ_\nu \left(\omega \right))}}`$ (A.25)
can be large enough to affect the total value of the moment of inertia due to the smallness of its denominator $`ϵ_\mu \left(\omega \right)+ϵ_\nu \left(\omega \right)`$ at this rotational frequency.
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# Variational dynamics in open spacetimes
## 1 Introduction
The idea that surface terms can be important when the Lagrangian method is applied to cosmology has been studied earlier in the context of spatially homogeneous but anisotropic models \- . In this case, a surface term appears when the Lagrangian is varied with respect to a spatially homogeneous metric perturbation, and the assumption of spatial homogeneity prevents the vanishing of this term when it is evaluated on an arbitrarily distant compact two-surface. In most other cases where the variational approach is applied to cosmology, surface terms are made to vanish trivially by evaluating only variations with respect to variables which vanish outside a bounded domain. The justification for this approach seems to be that one recovers the ‘correct’ field equations, which are the standard Euler-Lagrange equations. In a cosmological context, this way of reasoning can be questioned, both from a theoretical and an observational point of view. From observations, it is not a priori clear which are the correct equations of motion describing the dynamics of fields at length scales larger than the observable universe, and in different cosmological models. From a theoretical point of view, the relation between extremal action fields and classical physics has a natural foundation in quantum field theory. However, field configurations which vanish outside a bounded domain do not play a central role in quantum field theory, and this assumption may be questioned in a cosmological context when the spacetime itself is not bounded.
In this paper we will study this situation by means of an idealized model, which consist of a Klein-Gordon field in both a spatially flat and a spatially open FLRW universe. Our motivation for studying a scalar field stems from the aim to keep our equations simple, and the possible importance of these fields in the description of the early universe. We will concentrate on the open FLRW geometry, since this spacetime has some specific properties which allow surface terms to become important.
One of these properties is that eigenfunctions of the spatial Laplacian occur in two types. First, there are eigenfunctions with eigenvalues exceeding $`\frac{1}{6}`$ times the spatial curvature, and these eigenfunctions are complete in the space of square integrable functions . Second, there are eigenfunctions of the spatial Laplacian with eigenvalues between zero and $`\frac{1}{6}`$ times the spatial curvature. This last type of eigenfunctions cannot be square integrated, and they are responsible for long-range correlations in a spatially open universe . In spite of the fact that these perturbations cannot be square integrated, they may naturally occur in an open universe which is created in an exponentially expanding false vacuum , or they may be generated during preheating .
Another important property of the open Friedmann-Lemaître-Robertson-Walker (FLRW) geometry is that a spatial volume and the surface of its boundary grow at the same rate when the infinite volume limit is taken. The combination of large boundary surfaces and the presence of long-range correlations in open spacetimes appears to have an effect on the growth of surface terms at spatial infinity in these spacetimes.
Besides the theoretical reasons which make the open FLRW spacetime an interesting object to study, the open FLRW geometry has gained relevance as a model for the observed universe, with observations favoring a relatively small value of the density parameter . Furthermore, progress has been made in describing the creation of an open FLRW universe from an exponentially expanding false vacuum (see, e.g., ), and the theory of perturbations in open FLRW spacetimes has been worked out in greater detail \- .
This paper is structured as follows. In section 2 we discuss the physical relevance of action-extremizing field configurations, and we show that surface terms can contribute to the variation of the action for square-integrable perturbations. In section 4 we decompose the scalar field perturbations in terms of eigenfunctions of the spatial Laplacian, and we discuss the occurrence of supercurvature modes. The dynamics of the extremal action configurations is considered in section 5, and we recover the usual equation of motion for each perturbation component, with an additional source term, which can be expressed in terms of a surface integral which is evaluated at spatial infinity. We show that this source term can be neglected in the case where there are only subcurvature excitations of the scalar field, but it appears to diverge in the case where there are supercurvature perturbations. Due to this divergence, extremality of the action can only be defined in the restricted phase-space of field perturbations for which surface terms are finite. Depending on how one restricts the phase-space of field perturbations, a nontrivial source term contributes to the equation of motion for the extremal action configurations. In section 6, we consider the quantum correlation function of the scalar field. In the case where there are supercurvature perturbations, it is shown that the action functional is sensitive to degrees of freedom of the scalar field which have zero $`L_2`$-norm. It therefore appears that the correlation function is not well defined, unless one adopts nontrivial constraints on the phase-space of the scalar field, or one needs to include the zero-norm degrees of freedom in the integration over paths.
## 2 The extremal action principle
In this section we briefly review the variational approach to classical field theory. We then use arguments from quantum field theory to motivate a modified form of the variational method in a cosmological context. Surprisingly, it appears that non-local interactions at a classical level can emerge from the underlying quantum theory with a standard expression for the Lagrangian. While our explicit calculations involve only the simple case of a scalar field, our arguments are relevant in a more general field theoretical context, including general relativity. We will come back to this point at the very end of this paper.
One way of describing the dynamics of a classical field is by formulating a field equation. A specific solution of the field equation is determined by the boundary or periodicity conditions which apply to the system. It is of interest to note that the dynamics of the fields which can be observed in nature are described by field equations which act locally, while mathematical consistency does not require this. Hence, the dynamics of classical fields has a local aspect, in the sense that the field equations involve only the field variables and derivatives thereof at each point. Further, a particular classical field configuration is subject to global constraints, which act in the form of boundary or periodicity conditions. The work on this paper started as an attempt to establish whether the local aspects of the dynamics of fields, which is apparent from the structure of the field equations, are fundamental in nature.
In order to gain a deeper insight in the global aspects of the dynamics of fields, a Lagrangian approach appears to be most suitable. In this approach, an action functional is constructed from the field variables over the entire spacetime. The dynamics of the field then follows by requiring that the action is extremal in the space of field configurations. Establishing extremality of the action amounts to showing that the action does not vary at first order, for arbitrary infinitesimal perturbations of the field variables. It is essential to note that the field perturbations which are used to ‘test’ the extremality of the action in the classical description, are purely a mathematical construct. Further, we stress that the choice of the action functional is motivated with the aim to recover the field equations, and hence the Lagrangian description has the same physical content as the field equation.
At this point, let us formulate more precisely the question whether the local form of the interactions in nature is fundamental. On the one hand, it is well known that there exist conserved quantities which are related to global symmetries of the action . Although the existence of conserved quantities suggests an underlying global aspect of the dynamics of classical fields, this global aspect is in fact a consequence of applying Gauss’s theorem to a four-divergence, which vanishes locally as at each point in our spacetime as a consequence of the field equation. On the other hand, there is the question whether there can be a non-local coupling between physical fields, which acts at the level of the field equation. In particular, one would like to know whether non-local interactions at a classical level can emerge from an underlying quantum theory for which the Lagrangian has the usual local form. In this paper we will focus on this last question.
Let us now consider in some detail how classical field theory arises as a limit of an underlying quantum field theory. According to the Feynman path-integral approach to quantum field theory, the expectation value of an operator $`O`$ which acts on a field $`\psi `$, is given by the formal expression,
$$O=Z^1\text{d}[\psi ]O[\psi ]e^{\mathrm{𝑖𝑆}[\psi ]/\mathrm{}},$$
(1)
where $`S[\psi ]`$ is the action functional, $`Z`$ is a normalization constant, and $`\text{d}[\psi ]`$ is a measure on the space of field configurations (see, e.g., ). The integral is evaluated over all field configurations (paths) which are continuous and which satisfy certain initial or periodicity conditions. One should note that there is considerable difficulty involved in making the path-integral well defined, which is due to the fact that typical paths which contribute to the integral are non-differentiable. In our derivation, where we consider a free field, the different degrees of freedom decouple, and one can ignore those degrees of freedom which vary with infinite frequency.
As $`\mathrm{}`$ approaches zero in expression (1), the oscillatory behavior of the integrand suggests that the integral is dominated by those field configurations $`\psi `$ which are in some sense near to a field configuration $`\psi _0`$ which extremizes the action. Since $`\mathrm{}`$ is close to zero when expressed in terms of macroscopic units of time and energy, one therefore expects that classical physics is accurately described by an action extremizing field configuration $`\psi _0`$. The essential difference between this classical limit, and the classical theory which we discussed previously, is the fact that in the former case there are physical field perturbations which probe the phase-space nearby an action extremizing configuration, while in the latter case these field perturbations are purely a mathematical construct. As we will show in the following, this difference can give rise to an essentially different expression which describes the dynamics of the classical field.
As is well known, extremality of the action for $`\psi _0`$ implies that this configuration satisfies the classical field equations, provided that a surface term vanishes for all paths. In a classical variational treatment, surface terms are set to zero trivially by considering only paths which have compact support. However, this restriction on the type of paths does not occur in the sum over paths (1), and it seems natural to consider field configurations $`\psi _0`$ which extremize the action for the most general class of paths for which extremality of the action can be defined.
Indeed, one should note that in a classical treatment of cosmological perturbations one does not normally assume that perturbations must have compact support. However, if one accepts that classical perturbations do not have compact support, then it seems rather unnatural to require that quantum fluctuations about the classical field configurations have compact support. If this would be the case, then there would be a finite distance beyond which there are still classical perturbations while quantum fluctuations vanish. This appears to contradict the Copernican principle, which is commonly adopted in cosmology.
Considering the relation between classical and quantum physics, it should be mentioned that the path-integral approach does not only explain more than a classical approach (i.e., testable quantum effects), but one also needs to assume more than in classical physics (e.g., the existence of a classical regime , as well as various infinite subtractions ). One might therefore feel that the validity of the path-integral approach is as questionable as the classical variational approach, when it is applied to cosmological situations where it has not been tested. When seen in this light, the classical assumption that field-perturbations are restricted to have compact support is not proven to be wrong, but rather, it represents one possible choice in a more general class of boundary or asymptotic conditions. Whichever point of view one favors, it seems interesting to investigate the implications of relaxing the assumption that field-perturbations must have compact support. We will discuss these implications in the following.
## 3 Scalar field in FLRW geometry
The line element of the FLRW geometry is given by,
$$\text{d}s^2=\text{d}t^2+a^2(t)\left[\text{d}\chi ^2+c^2\mathrm{sinh}^2c\chi (\text{d}\theta ^2+\mathrm{sin}^2\theta \text{d}\varphi ^2)\right],$$
(2)
where $`c=𝐑^+`$ for the spatially open geometry, while the spatially flat and closed geometry are obtained by taking the limit $`c0`$ or by choosing $`ci\times 𝐑^+`$ respectively. We will refer to the geometry with the line element (2) as $``$, while a spatial hypersurface of constant time $`t`$ is referred to as $`\mathrm{\Sigma }`$.
It follows directly from expression (2) that the surface of a spatial sphere of constant radius $`\chi _0`$ grows as fast as the three-volume inside the sphere, when one considers the limit where $`\chi _0\mathrm{}`$. One may therefore expect that surface terms can be equally important as volume terms when we take the infinite volume limit in an open universe. This situation is essentially different from the situation in a spatially flat spacetime, where the surface of a spatial sphere of constant radius $`\chi `$ grows by one power of $`\chi `$ less fast then the three-volume which is contained inside the sphere.
We will consider a scalar field $`\psi `$, which is described by the Lagrangian density
$$[\psi ]=\frac{1}{2}\sqrt{g}(g^{\mu \nu }_\mu \psi _\nu \psi +m^2\psi \psi ),$$
(3)
where $`g_{\mu \nu }`$ denotes the FLRW metric (2), and $`g=det(g_{\mu \nu })`$.
We define the action of the $`\psi `$-field as the integral of the Lagrangian density (3) over the entire spacetime,
$$S[\psi ]:=\text{d}^4x[\psi ].$$
(4)
Note that the integral in this expression does not need to converge. This is not necessarily a problem if one is interested in calculating the variation of the action under a change of the field from $`\psi `$ to $`\psi +\delta \psi `$, where $`\delta \psi `$ is a suitably small ‘test-perturbation’. The question arises which restriction one has to impose on the test-perturbations $`\delta \psi `$ so that the first-order variation $`\delta S`$ is well defined. The first-order variation of the action (4) follows by the standard procedure of functional derivation,
$$\delta S=\text{d}^4x\left(\frac{\delta }{\delta \psi }\delta \psi +\frac{\delta }{\delta _\mu \psi }\delta _\mu \psi \right).$$
(5)
By partially integrating equation (5), where the Lagrangian is given by expression (3), we obtain
$$\delta S=\text{d}^4x\sqrt{g}\left[\psi _{;\mu }^{;\mu }m^2\psi \right]\delta \psi \text{d}^4x\sqrt{g}(\psi ^{;\mu }\delta \psi )_{;\mu },$$
(6)
where a semicolon denotes the covariant derivative.
Provided that the second term on the right-hand side of equation (6) vanishes for nonzero perturbations $`\delta \psi `$, then the condition $`\delta S=0`$ implies the vanishing of the term in brackets, and hence the field equation holds. This is the case when we consider test-perturbations $`\delta \psi D`$, where $`D`$ is defined as the class of perturbations which are bounded and which have compact support. However, as we mentioned in the beginning of this section, the restriction to test-perturbations $`\delta \psi D`$ does not follow from known physical principles, when the spacetime itself is non-compact. Let us therefore try to determine the largest class of test-perturbations for which the variation of the action is well defined. For a scalar field $`\psi `$, and a Lagrangian which is bi-linear in the field variable, it is clear that square integrability of $`\delta \psi `$ is a necessary condition for the existence of the variation of the action (6), i.e., we require $`\delta \psi L_2()`$. It is not a priori clear whether $`\delta \psi L_2()`$ is a sufficient condition for the existence of the variation of the action (6), and it may be necessary to restrict the type of test-perturbations further to ensure that $`\delta S`$ exists. Assuming that we are able to determine the largest class of test-perturbations $`\delta \psi `$ for which $`\delta S`$ exists, then it remains a question whether there exist field configurations $`\psi _0`$ such that $`\delta S`$ vanishes for all perturbations $`\delta \psi `$ about $`\psi _0`$.
Let us first address the question whether the restriction $`\delta \psi L_2()`$ is sufficient to ensure the existence of $`\delta S`$. The answer to this question is negative, which we show by an example where the contribution of surface terms to $`\delta S`$ diverges, while $`\psi `$ is a solution of the field equation and $`\delta \psi L_2()`$. Since we will focus on surface effects at spatial infinity, we require that $`\delta \psi `$ can be square integrated over a spatial hypersurface of constant time in the geometry (2), i.e., $`\delta \psi L_2(\mathrm{\Sigma })`$, while we do not specify the time dependence of $`\delta \psi `$. It is clear from the expression of the line element (2) that a square integrable test-perturbation $`\delta \psi `$ must approach zero faster than $`1/\chi `$ in the spatially flat case, and faster than $`e^\chi `$ in the spatially open case. A specific example of a square integrable test-perturbation is given by
$$\delta \psi =(1+\chi )^{(1+\alpha )}_\chi \psi \text{and}\delta \psi =e^{(1+\alpha )\chi }_\chi \psi ,$$
(7)
in the spatially flat and open case respectively, and $`\alpha 𝐑^+`$. By substituting expressions (7) for $`\delta \psi `$ into equation (6), and using (2), we find
$$\delta S=4\pi a^2(t)\text{d}t\text{d}\mathrm{\Omega }\underset{\chi \mathrm{}}{lim}F(\chi )(_\chi \psi )^2,$$
(8)
where $`\text{d}\mathrm{\Omega }`$ denotes the volume element on the unit two-sphere, and $`F(\chi )=\chi ^{1\alpha }`$ in the spatially flat case, and $`F(\chi )=e^{(1\alpha )\chi }`$ in the spatially open case. Indeed, expression (8) diverges for some values of $`\alpha (0,1]`$, provided that the term $`_\chi \psi `$ does not approach to zero as fast as $`F^{1/2}(\chi )`$ in the limit where $`\chi \mathrm{}`$. The variation of the action (8) can therefore be arbitrarily large, for $`\delta \psi L_2(\mathrm{\Sigma })`$.
Let us now address the question whether there exist configurations of the $`\psi `$-field which extremize the action for all $`\delta \psi L_2(\mathrm{\Sigma })`$, in the cosmologically interesting case where $`\psi `$ and $`_\chi \psi `$ do not vanish at spatial infinity. We show that the answer to this question is negative. We will therefore use a result which is derived in the following, which states that a field configuration which extremizes the action for all $`\delta \psi L_2(\mathrm{\Sigma })`$ must be a solution of the field equation. We combine this with the result which was derived earlier in this section, which shows that a solution of the field equation for which $`_\chi \psi `$ does not approach to zero at spatial infinity, does not extremize the action for all $`\delta \psi L_2(\mathrm{\Sigma })`$. Hence, it follows that action extremizing configurations do not exist for $`\delta \psi L_2(\mathrm{\Sigma })`$ and $`_\chi \psi `$ not approaching to zero at infinity.
In deriving the proof above, we assumed that a field configuration which extremizes the action for all $`\delta \psi L_2(\mathrm{\Sigma })`$ must be a solution of the field equation. In order to proof this, let us recall that for $`\delta \psi D`$, i.e., the class of test-perturbations which are bounded and which have compact support, extremality of the action implies that the field equation holds and vice-versa. Configurations which do not satisfy the field equation can therefore not extremize the action for all $`\delta \psi D`$, and since $`DL_2(\mathrm{\Sigma })`$ these configurations do not extremize the action for all $`\delta \psi L_2(\mathrm{\Sigma })`$. Hence it follows that a field configuration which extremizes the action for all $`\delta \psi L_2(\mathrm{\Sigma })`$ must be a solution of the field equation, which proves our assumption.
The observation that action extremizing configurations do not in general exist for $`\delta \psi L_2(\mathrm{\Sigma })`$, implies that the usual identification between classical physics and action extremizing configurations becomes ambiguous when we allow for perturbations which do not fall off sufficiently fast at infinity. There are several ways by which one could try to resolve the problem which is posed by the non-existence of extremal action configurations for test-perturbations $`\delta \psi L_2(\mathrm{\Sigma })`$. We will discuss these possible solutions in the following.
First, let us recall that the restriction $`\delta \psi L_2(\mathrm{\Sigma })`$ was found to be necessary to ensure finiteness of $`\delta S`$, but due to the contribution of a surface term to $`\delta S`$ this restriction is not sufficient. This observation suggests that the class of test-perturbations $`\delta \psi `$ should be restricted further, such that $`\delta S`$ is finite for all $`\delta \psi `$. Although finiteness of $`\delta S`$ is easily achieved by requiring that the test-perturbations $`\delta \psi `$ fall off sufficiently fast, this does not imply that extremal action configurations exist in the space of test-perturbations for which $`\delta S`$ is finite. The reason for this is that the existence of extremal action configurations requires that the surface term contribution to $`\delta S`$ vanishes completely, which is clearly a stronger restriction on $`\delta \psi `$ than the condition that $`\delta S`$ is finite. Although one could restrict $`\delta \psi `$ to ensure that the surface term contribution to $`\delta S`$ vanishes completely, this would be rather add-hoc since it is not shown that this is the only possible restriction on the class of test-perturbations for which extremal action configurations exist.
Instead of restricting the class of test-perturbations, one could also attempt to remove the contribution of surface terms to $`\delta S`$ by modifying the Lagrangian density (3). Let us therefore note that the choice of the Lagrangian density is motivated by the fact that one recovers the Klein-Gordon equation, provided that the variation of the action and the surface term in equation (6) vanish. In the classical variational approach, where surface terms are made to vanish by assuming boundary conditions on $`\delta \psi `$, one therefore has the freedom to add a term to the Lagrangian density which has the form of a four-divergence, since the variation of this term equals a vanishing surface term. In this section we questioned the assumption that the surface term in equation (6) vanishes in perturbed flat and open FLRW spacetimes. However, it is conceivable that one can add a four-divergence term to the Lagrangian (3) such that its variation cancels the surface term in equation (6). Indeed, in the context of Hamiltonian cosmology, as well as in quantum cosmology, it appears to be natural to add a surface term to the Einstein-Hilbert action which has the property that its variation cancels an identical term which arises from the variation of the Einstein-Hilbert action \- .
Let us now consider whether the same possibility exists in the case where we are dealing with a scalar field. We therefore add a generic surface term to the action (4), which has the form
$$S_B[\psi ]=\frac{1}{2}\text{d}^4x\sqrt{g}B_{;\mu }^\mu ,$$
(9)
where $`B^\mu =B^\mu [\psi ]`$, and then we consider whether the variation of this surface term may cancel the surface term in equation (6). The variation of $`B^\mu `$ follows by the method of functional derivation, i.e., treating $`\psi `$ and $`_\nu \psi `$ as independent variables:
$$\delta B^\mu =\frac{\delta B^\mu }{\delta \psi }\delta \psi +\frac{\delta B^\mu }{\delta _\nu \psi }\delta _\nu \psi ,$$
(10)
where we used that $`B^\mu `$ cannot depend on higher than first-order derivatives of $`\psi `$. It is clear that any dependence of $`B^\mu `$ on higher than first-order derivatives of $`\psi `$ contributes terms to the variation of the action which are proportional to the variation of higher than first-order derivatives of $`\delta \psi `$. These terms cannot cancel against the surface term in equation (6), which contains at most first-order derivatives of $`\delta \psi `$, although a cancellation was required.
The requirement that the surface term in equation (6) cancels the surface term which arises from the variation of $`S_B`$ results in the conditions
$$\frac{\delta B^\mu }{\delta \psi }=\psi ^{;\mu },\text{and}\frac{\delta B^\mu }{\delta _\nu \psi }=0,$$
(11)
for all $`\mu ,\nu `$, and we used expression (10). The first condition in equation (11) constrains $`B^\mu `$ to be of the form $`B^\mu =\psi ^{;\mu }\psi +`$ c<sub>1</sub>, where $`c_1`$ is a functional which does not depend on $`\psi `$, while the second condition constrains $`B^\mu `$ to be a functional which does not depend on $`_\mu \psi `$. Clearly, both requirements are exclusive, and there exists no functional $`B^\mu `$ such that the variation of $`S_B`$, (9), cancels the surface term in equation (6). Note, however, that the precise form of the surface term in equation (6) does change by adding a term of the form (9) to the action. Hence, the contribution of a surface term to the variation of the scalar field action (4) appears to be generic, although its precise form is ambiguous. In the following calculation we will retain the surface term which appears in equation (6), which means that we assume $`S_B`$ to vanish.
Having considered the possibility to adopt further restrictions on the type of test-perturbations, as well as modifying the action by adding a surface term contribution, we have not found an argument which shows us that we can neglect the contribution of a surface term to the variation of the action. However, taking the surface term in equation (6) seriously confronts us with the problem that field configurations which extremize the action in the space of test-perturbations for which $`\delta S`$ is well defined, do not in general exist. It should be noted, however, that the non-existence of action extremizing field configurations does not need to be a problem if one could show that those test-perturbations for which the action functional is not extremal, have a zero phase-space measure in the space of fields $`\psi `$. Indeed, it is clear that paths of the form (8), which yield large surface terms at spatial infinity, are highly special in the sense that the asymptotic behavior of these paths is correlated with the field $`\psi `$ about which we expand. Therefore, one expects that these paths occupy a very small amount of phase-space in the space of field configurations in which extremality of the action is considered, and their relevance for the dynamics of the $`\psi `$-field may be negligible. Note, however, that precisely the same argument applies to the case where $`\delta \psi D`$, since in this case $`\delta \psi `$ is specified to be exactly equal to zero for arbitrarily large radii $`\chi `$. In order to make these considerations quantitative, it is necessary to introduce a measure on the phase-space of the $`\psi `$-field. We will address this problem in the following sections.
## 4 Perturbations in open FLRW
In order to obtain a quantitative description of the space of field configurations of the scalar field $`\psi `$, it is useful to decompose $`\psi `$ and test-perturbations $`\delta \psi `$ in terms of eigenfunctions of the spatial Laplacian which are complete in the space $`L_2`$ of functions which are square integrable on the hypersurfaces $`\mathrm{\Sigma }(t)`$. The reason why it is convenient to use eigenfunctions of the spatial Laplacian, is that this operator is present in the expression for the variation of the action (6). When we ignore the surface term, it is therefore clear that each eigenfunction only couples to itself, and the dynamics of each mode is independent of the dynamics of all other modes.
Let $`Q(x)`$ be a solution of the Helmholtz equation, i.e.,
$$Q_{;i}^{;i}+(k/a)^2Q=0,$$
(12)
where $`;i`$ denotes the covariant derivative with respect to the coordinate $`x^i\{r,\theta ,\varphi \}`$ in the geometry (2), $`a=a(t)`$ denotes the scale factor, and $`k𝐑^+`$. In the following, we concentrate on the spatially open geometry (2), while we consider the spatially flat spacetime as a limiting case of the spatially open geometry. A basis of solutions of equation (12), which are complete in the space of $`L_2`$ functions on $`\mathrm{\Sigma }(t)`$, and which factorize in terms of an angular and a radially dependent part, is given by
$$Z_{qlm}=\mathrm{\Pi }_{ql}(\chi )Y_{lm}(\theta ,\varphi ),$$
(13)
where $`Y_{lm}`$ are the standard spherical harmonics on the unit two-sphere, and the radially dependent functions $`\mathrm{\Pi }_{ql}(\chi )`$ are solutions of the equation
$$\frac{1}{g_2}\frac{}{\chi }g_2\frac{}{\chi }\mathrm{\Pi }_{ql}(\chi )=\left(k^2\frac{l(l+1)}{g_2}\right)\mathrm{\Pi }_{ql}(\chi ),$$
(14)
where $`g_2=c^2\mathrm{sinh}^2c\chi `$. Equation (14) has solutions of the form,
$$\mathrm{\Pi }_{ql}(\chi )=N_{ql}(\mathrm{sinh}c\chi )^l\left(\frac{1}{\mathrm{sinh}c\chi }\frac{\text{d}}{\text{d}\chi }\right)^{l+1}\mathrm{cos}(qc\chi ),$$
(15)
where $`q`$ is defined by $`q^2=k^2/c^21`$, and
$$N_{ql}:=\sqrt{\frac{2}{\pi }}\left[\underset{n=0}{\overset{l}{}}(n^2+q^2)\right]^{1/2}$$
(16)
is a normalization factor . Notice that the $`q=0`$ mode solves the Helmholtz equation (12) with a nonzero eigenvalue equal to $`c^2/a^2`$, which equals $`\frac{1}{6}`$ times the spatial curvature in the geometry (2).
The radial solutions for the spatially flat geometry are obtained by taking the limit $`c0`$ in expression (15), keeping $`k`$ fixed,
$$\underset{c0}{lim}\mathrm{\Pi }_{ql}(\chi )=\sqrt{\frac{2}{\pi }}kj_l(k\chi ),$$
(17)
where $`j_l`$ denotes the spherical Bessel function . From now on, we assume that the spacetime is open, such that $`c𝐑^+`$, and without loss of generality we may set $`c=1`$ in expression (2) by absorbing a factor $`c`$ in the definition of the comoving radial coordinate $`\chi `$ and by absorbing a factor $`c^1`$ in the definition of the scale factor $`a(t)`$.
It follows from expression (15) that the radial functions $`\mathrm{\Pi }_{ql}`$ can be written as the product of an oscillating factor $`\mathrm{cos}q\chi `$ or $`\mathrm{sin}q\chi `$, and a factor which approaches to zero exponentially as $`\mathrm{sinh}^1\chi `$ in the limit where $`\chi \mathrm{}`$. Since the modes $`Z_{qlm}`$ with $`q𝐑^+`$ vary at comoving length scales which are typically smaller than the curvature scale which we have set equal to one in the FLRW geometry (2), these modes are called subcurvature modes.
There exist solutions of the Helmholtz equation (12) for which $`k^2(0,1]`$, which corresponds to imaginary values of $`qi\times (0,1]`$. The explicit expression for these modes is still given by equation (15), where the factor $`\mathrm{cos}(q\chi )`$ is replaced by $`\mathrm{cosh}(|q|\chi )`$. The modes $`Z_{qlm}`$ with $`qi\times (0,1]`$ approach to zero as a constant times $`\mathrm{exp}((|q|1)\chi )`$ in the limit where $`\chi \mathrm{}`$, and since they vary at length scales greater than the curvature scale one calls them supercurvature modes.
We define the spatial integration operation by
$$f:=\underset{ϵ0}{lim}f(ϵ)$$
(18)
where
$$f(ϵ):=\text{d}\mathrm{\Omega }_0^{1/ϵ}\text{d}\chi \mathrm{sinh}^2(\chi )f.$$
(19)
and $`\text{d}\mathrm{\Omega }^2`$ denotes the volume element on the unit two-sphere. The subcurvature modes $`Z_{qlm}(q𝐑^+)`$ are orthonormal with respect to spatial integration,
$$Z_{qlm}Z_{q^{}l^{}m^{}}=\delta (qq^{})\delta _{ll^{}}\delta _{mm^{}},$$
(20)
and they are known to be complete in the space $`L_2(\mathrm{\Sigma })`$ , which consists of equivalence classes of functions $`f`$ for which $`|f|^2`$ exists, where we identify functions $`f`$ which differ only on a set of Lebesgue measure zero.
For the supercurvature modes, the indefinite integral over the radius in expression (20) does not exist, so that these modes cannot be normalized in the $`L_2(\mathrm{\Sigma })`$ sense. Furthermore, expression (20) diverges when only one of the modes $`Z`$ corresponds to a supercurvature mode, and $`l=l^{}`$ and $`m=m^{}`$. Therefore, the supercurvature modes cannot be decomposed in terms of the subcurvature modes. Mechanisms which may be responsible for the generation of supercurvature perturbations in open spacetimes have been investigated in .
The $`\psi `$-field may be expanded in terms of the modes $`Z_{qlm}`$,
$$\psi (x,t)=\psi ^{}(x,t)+\psi ^+(x,t),$$
(21)
where
$$\psi ^{}(x,t):=\underset{lm}{}_0^{\mathrm{}}\text{d}q\psi _{qlm}(t)Z_{qlm}(x),$$
(22)
$$\psi ^+(x,t):=\underset{lm}{}_0^i\text{d}\overline{q}\psi _{\overline{q}lm}(t)Z_{\overline{q}lm}(x),$$
(23)
where $`x=\{\chi ,\theta ,\varphi \}`$, and the integration over $`\overline{q}`$ runs along the imaginary axis in the complex $`\overline{q}`$-plane.
An important class of perturbations, which is believed to occur in the early universe, corresponds to the case where the coefficient of each independent mode is chosen according to a Gaussian probability distribution (see, e.g., \- ). For this type of perturbation, which is called a ‘Gaussian perturbation’ or ‘random-field’, there are no correlations between the coefficients $`\psi _{qlm}`$ for different values of $`q,l`$, and $`m`$. The statistical properties of a random-field are determined by the variance of the Gaussian probability distribution, which we call $`\sigma `$. In the generic case, where $`\sigma `$ depends on $`q,l,`$ and $`m`$, one cannot determine the variances $`\sigma (q,l,m)`$ from a single realization of a random-field, which is determined by the set of coefficients $`\psi _{qlm}`$. Instead, one would need an infinite ensemble of random-fields, in order to deduce the statistical properties, i.e., the variances $`\sigma (q,l,m)`$, according to which these random-fields are generated.
Let us now define the ensemble average of a functional as the weighted sum of this functional over all random-fields in an ensemble, where the weight factor is given by the probability for each specific random-field to occur. This allows us to define the two-point correlation function of the $`\psi `$-field as the ensemble average of $`\psi (x)`$ times $`\psi (x^{})`$. A random-field $`\psi (x)`$ is said to be statistically homogeneous and isotropic when the two-point correlation function is invariant under the group of isometries on $`\mathrm{\Sigma }`$, i.e., the group of rotations and spatial translations. Clearly, the two-point correlation function of a statistically homogeneous and isotropic random field can only be a function of a distance measure which is invariant under the group of isometries on $`\mathrm{\Sigma }`$, and we can take this distance measure to be the length $`d(x,x^{})`$ of a geodesic which relates the points $`x`$ and $`x^{}`$. In can be shown that statistical homogeneity and isotropy of a random-field $`\psi (x)`$ holds if and only if the variances $`\sigma `$ do not depend on the labels $`l`$ and $`m`$ .
Although it seems rather artificial to introduce the concept of an ensemble in the context of cosmology, since we can only observe one universe, a physical interpretation of the ensemble average is provided by the property of ergodicity. In the context of random-fields, ergodicity is defined as the equivalence of ensemble averaging and spatial averaging, where the spatial average of the two-point correlation function is defined by summing $`\psi (x)`$ times $`\psi (x^{})`$ over random sets of points $`x`$ and $`x^{}`$ for which the geodesic distance $`d(x,x^{})`$ has a specific value. In the case where $`\mathrm{\Sigma }`$ is a Euclidean three-space, ergodicity can be proven to hold under fairly weak assumptions , but for a hyperbolic three-space no proof seems to be known, while it is usually assumed.
In the following, we will assume a Gaussian statistically homogeneous and isotropic spectrum of subcurvature perturbations. One should note that this type of perturbation cannot be square integrated. This follows by substituting the expansion of $`\psi `$, (22), into the hypersurface integral (18) and using the orthonormality relation (20). The resulting expression contains an indefinite sum over $`l`$ and $`m`$ of the squared coefficient $`\psi _{qlm}`$, and this sum diverges when the variance $`\sigma (q)`$ is nonzero. It is therefore clear that the property of non-square integrability is not specifically related to the presence of supercurvature modes.
## 5 Extremal action dynamics
Let us now calculate the variation of the action (6), which is evaluated over a bounded spatial volume $`V(\chi _0)`$, which we define as those points in the geometry (2) for which $`\chi <\chi _0`$, and then we consider the limit where $`\chi _0\mathrm{}`$. We obtain
$$\delta S=\text{d}t\underset{\chi _0\mathrm{}}{lim}[a^3\text{d}\mathrm{\Omega }^2_0^{\chi _0}\text{d}\chi \mathrm{sinh}^2\chi \delta \psi (\frac{1}{\sqrt{g}}_\mu g^{\mu \nu }\sqrt{g}_\nu m^2)\psi $$
$$a\mathrm{sinh}^2\chi \text{d}\mathrm{\Omega }^2\delta \psi _\chi \psi |_{\chi =\chi _0}],$$
(24)
where $`a=a(t)`$. Using the definition of the integration operation (20), expression (24) can be written in the form,
$$\delta S=\text{d}t[a^3\delta \psi (\frac{1}{\sqrt{g}}_\mu g^{\mu \nu }\sqrt{g}_\nu m^2)\psi $$
$$a\underset{\chi _0\mathrm{}}{lim}\mathrm{sinh}^2\chi _0\text{d}\mathrm{\Omega }\delta \psi _\chi \psi |_{\chi =\chi _0}].$$
(25)
We will consider separately the cases where the expansion of the field $`\psi `$ includes only subcurvature modes, and the case where the expansion includes supercurvature modes as well.
### 5.1 Open spacetime with subcurvature perturbations
Let us first consider the case where the field $`\psi `$ can be expanded in terms of only subcurvature modes, i.e., we assume that $`\psi _{qlm}=0`$ for all $`qi\times (\mathit{0},\mathit{1}]`$, so that only the first term in the expansion of the field (21) is nonzero. Equation (25) can then be evaluated separately for each mode, by substituting the expansion (21) into expression (25), and using the orthonormality relation (20). We obtain
$$\delta S=\text{d}t\text{d}q\underset{l,m}{}\delta \psi _{qlm}(t)[a^3(\frac{1}{\sqrt{g}}_0g^{00}\sqrt{g}_0a^2(t)k^2m^2)\psi _{qlm}(t)$$
(26)
$$\underset{\chi _0\mathrm{}}{lim}a\mathrm{sinh}^2\chi _0\text{d}q^{}\psi _{q^{}lm}\mathrm{\Pi }_{ql}_\chi \mathrm{\Pi }_{q^{}l}|_{\chi =\chi _0}].$$
The requirement that the variation of the action vanishes for nonzero perturbations $`\delta \psi _{qlm}(t)`$ implies an equation of motion for each perturbation component $`\psi _{qlm}(t)`$, namely,
$$\left(\frac{1}{\sqrt{g}}_0g^{00}\sqrt{g}_0a^2k^2m^2\right)\psi _{qlm}(t)=J_{qlm},$$
(27)
where
$$J_{qlm}:=\underset{\chi _0\mathrm{}}{lim}\left[a^2\mathrm{sinh}^2\chi _0\text{d}q^{}\psi _{q^{}lm}\mathrm{\Pi }_{ql}_\chi \mathrm{\Pi }_{q^{}l}|_{\chi =\chi _0}\right].$$
(28)
Note that $`J_{qlm}`$ acts as a source term in equation (27), and this term couples perturbations which have the same angular wave numbers $`l`$ and $`m`$. One would like to know whether the limit in expression (28) exists, and whether or not this term can be neglected. In order to answer this question, we need to evaluate the integral over $`q^{}`$ of the distribution $`\psi _{q^{}lm}`$, which is multiplied by a factor which is of order unity. According to equation (20) and (21), the distribution $`\psi _{q^{}lm}`$ can be defined by,
$$\psi _{q^{}lm}=\underset{ϵ0}{lim}\psi _{q^{}lm}(ϵ),$$
(29)
where
$$\psi _{q^{}lm}(ϵ):=Z_{q^{}lm}\psi (ϵ)$$
(30)
and the limit $`ϵ0`$ should be evaluated after the integration over $`q^{}`$ is performed. When we integrate over a bounded volume, then the modes $`Z_{qlm}`$ are dependent in the sense that their overlap $`Z_{qlm}Z_{q^{}lm}(ϵ)`$ is nonzero and of the order of $`ϵ^1`$ for $`qq^{}`$ of the order of $`ϵ`$. The number of independent modes in a fixed $`q^{}`$-interval therefore tends to diverge as $`ϵ^1`$ in the limit where $`ϵ0`$. In the previous section, we introduced the concept of a Gaussian perturbation. In order to generate a Gaussian perturbation which has an amplitude of order one, the coefficients $`\psi _{qlm}`$ in the expansion of the field (22) need to be uncorrelated for values of $`q`$ differing more than $`ϵ`$, while the amplitude of the coefficients must diverge as $`ϵ^{\frac{1}{2}}`$ when $`ϵ0`$. The asymptotic behavior of the integral over $`q^{}`$ in expression (28) can therefore be estimated as the sum of $`ϵ^1`$ uncorrelated numbers which are of the order of $`ϵ^{1/2}`$, multiplied by a $`q^{}`$-interval which is of the order of $`ϵ`$. In the limit where $`ϵ0`$, the term between brackets in expression (28) will therefore remain of order one, and the expression does not converge. Note, however, that the left-hand side of the equation of motion (27) is proportional to the coefficient $`\psi _{qlm}`$, which diverges as $`ϵ^{1/2}`$ in the limit where $`ϵ0`$. We therefore find that the source term on the right-hand side of equation (27) can be neglected in the infinite volume limit, when the perturbations of the field are Gaussian and of the subcurvature type.
### 5.2 Open spacetime with supercurvature perturbations
Let us now attempt to derive an equation of motion for the $`\psi `$-field, in the case where the expansion of the $`\psi `$-field (21) includes supercurvature perturbations.
We may therefore substitute the expansion of the $`\psi `$-field (21) in the expression for the variation of the action (24), which yields,
$$\delta S=\text{d}t\underset{\chi _0\mathrm{}}{lim}$$
(31)
$$\times [a^3\text{d}\mathrm{\Omega }^2_0^{\chi _0}\text{d}\chi \mathrm{sinh}^2\chi \delta \psi (\frac{1}{\sqrt{g}}_\mu g^{\mu \nu }\sqrt{g}_\nu m^2)(\psi ^{}+\psi ^+)$$
$$a\text{d}\mathrm{\Omega }^2\mathrm{sinh}^2\chi \delta \psi _\chi (\psi ^{}+\psi ^+)|_{\chi =\chi _0}].$$
Using the definition of the integration operation (18), and expression (21), we recover expression (27), with an additional source term which accounts for the coupling between subcurvature and supercurvature perturbations, i.e.,
$$\left(\frac{1}{\sqrt{g}}_0g^{00}\sqrt{g}_0a^2k^2m^2\right)\psi _{qlm}^{}(t)=J_{qlm}+J_{qlm}^+,$$
(32)
where $`q𝐑^+`$, $`J_{qlm}`$ is given by expression (28), and
$$J_{qlm}^+:=\underset{\chi _0\mathrm{}}{lim}_0^i\text{d}\overline{q}[(\frac{1}{\sqrt{g}}_0g^{00}\sqrt{g}_0a^2k^2m^2)\psi _{\overline{q}lm}^+(t)$$
$$\times _0^{\chi _0}\text{d}\chi \mathrm{sinh}^2\chi \mathrm{\Pi }_{ql}\mathrm{\Pi }_{\overline{q}l}+a^2\mathrm{sinh}^2\chi _0\psi _{\overline{q}lm}^+\mathrm{\Pi }_{ql}_\chi \mathrm{\Pi }_{\overline{q}l}|_{\chi =\chi _0}].$$
(33)
Note that both terms which contribute to expression (33) diverge exponentially in the limit where $`\chi _0\mathrm{}`$, and the limit in this expression does not exist, unless the divergent terms cancel. Let us therefore observe that the two terms at the right-hand side of equation (33) diverge exponentially as $`\mathrm{exp}|\overline{q}\chi |`$, (see section 4), and both terms oscillate due to the radial function $`\mathrm{\Pi }_{qlm}`$. A cancellation of the divergent terms in equation (33) requires that both terms oscillate with the same phase. By re-writing equation (33), using,
$$_0^{\chi _0}\text{d}\chi \mathrm{sinh}\chi \mathrm{\Pi }_{ql}\mathrm{\Pi }_{\overline{q}l}=\frac{\mathrm{sinh}^2\chi _0}{q^2\overline{q}^2}\left|\mathrm{\Pi }_{ql}_\chi \mathrm{\Pi }_{\overline{q}l}\mathrm{\Pi }_{\overline{q}l}_\chi \mathrm{\Pi }_{ql}\right|_{\chi =\chi _0},$$
(34)
one finds that $`J_{qlm}^+`$ diverges as the product of an exponential factor $`\mathrm{exp}(|\overline{q}|+1)\chi `$, multiplied by the sum of two terms which oscillate out of phase as $`\mathrm{\Pi }_{ql}`$ and $`_\chi \mathrm{\Pi }_{ql}`$, respectively. Therefore, the right-hand side of equation (32) diverges, and we cannot use this equation to describe the time-evolution of the perturbation component $`\psi _{qlm}(t)`$. Recall that in the absence of supercurvature perturbations, surface terms appeared to give rise to a negligible correction to the equation of motion for each perturbation component $`\psi _{qlm}(t)`$, which followed by requiring that $`\delta S=0`$ for all $`\delta \psi L_2(\mathrm{\Sigma })`$. When supercurvature perturbations are present, equations (31) and (32) show that it is precisely a surface term which contributes a divergent term to the variation of the action for all $`\delta \psi Z_{qlm}`$. In this case, the extremal action condition $`\delta S=0`$ cannot be satisfied for all $`\delta \psi L_2(\mathrm{\Sigma })`$, irrespectively of the equation of motion which the field satisfies. It is however clear that the condition $`\delta S=0`$ must have solutions when test-perturbations are confined to some subspace of $`L_2(\mathrm{\Sigma })`$ for which $`\delta S`$ is well defined. We will determine these subspaces in the following.
According to expressions (25) and (21), the surface term which contributes to $`\delta S`$ behaves asymptotically as $`\delta \psi `$ times a factor $`\mathrm{sinh}^2\chi _\chi \psi ^+`$ in the limit where $`\chi \mathrm{}`$. The contribution of surface terms to the variation of the action (24) will therefore be finite and convergent, provided that $`\mathrm{sinh}^2\chi \delta \psi _\chi \psi ^+`$ converges when $`\chi \mathrm{}`$. Let us now define the class of test-perturbations $`\{\delta \psi \}_c`$ by the requirement that $`\mathrm{sinh}^2\chi \delta \psi _\chi \psi ^+`$ converges to a constant $`c𝐑`$ when $`\chi \mathrm{}`$.
Note that it follows from the definition of $`\{\delta \psi \}_c`$ that $`\{\delta \psi \}_c`$ contains $`D`$, i.e., the class of functions which are bounded and which have compact support. As is well known, the class of functions $`D`$ is infinite dimensional in the sense that there exists a denumerable infinite set of linearly independent basis-functions which is complete in $`D`$ , and therefore $`\{\delta \psi \}_c`$ must be infinite dimensional, for arbitrary $`c𝐑`$. It is therefore not clear whether one class of test-perturbations $`\{\delta \psi \}_c`$ for some specific value of $`c𝐑`$ dominates in terms of the phase-space which is occupied by these test-perturbations. We will make this statement more precise in the following section, where it is shown that that the classes of test-perturbations $`\{\delta \psi \}_c`$, for different values of $`c𝐑`$, are equivalent up to variations with vanishing $`L_2()`$-norm.
Summarizing, we found that the contribution of surface terms to the variation of the action diverges for square integrable field perturbations which do not fall off at a specific rate, depending on the spectrum of supercurvature perturbations. In the presence of supercurvature perturbations, extremality of the action can therefore only be defined with respect to a restricted class of field perturbations. Surface terms contribute a non-trivial source term to the standard Klein-Gordon equation, but the magnitude thereof depends on the choice of the restricted class of test-perturbation with respect to which the action is extremized. The dynamics of the ‘classical’ field configurations therefore remains undetermined, unless one finds a physical argument which constrains the phase-space of the $`\psi `$-field uniquely.
## 6 Quantum correlations
In the previous section we showed that surface terms constrain the phase-space of test-perturbations for which the variation of the Klein-Gordon action is finite, in an open FLRW spacetime with supercurvature perturbations. One may also question whether the nontrivial surface terms which we found have an effect on quantum correlations of the $`\psi `$-field. As is clear from expression (1), the quantum correlation function of the $`\psi `$-field can be expressed as a weighted integral over all continuous field configurations, and the weight factor depends on the source term $`J^+`$, which may be infinite.
The two-point correlation function is given by the formal expression (see, e.g., )
$$\tau (x,x^{}):=Z^1\text{d}[\psi ]\psi (x)\psi (x^{})e^{\mathrm{𝑖𝑆}[\psi ]/\mathrm{}},$$
(35)
where $`x`$ denotes the set of coordinates on $``$.
The standard method to calculate the two-point correlation function is to expand the field $`\psi `$, about some background configuration $`\psi _0`$, in terms of a denumerable complete set of solutions of the four-dimensional Helmholtz equation (see, e.g., for the details involved in this calculation). Since $`L_2()`$ is known to be separable, there exists a denumerable and complete set of solutions, which we call $`\psi _i`$, and we can choose these solutions to be orthonormal in $`L_2()`$. A generic expansion of the field $`\psi `$, about a configuration $`\psi _0`$, takes the form
$$\delta \psi :=\psi \psi _0=\underset{i}{}a_i\psi _i,$$
(36)
where $`a_i𝐑`$. Further, the measure on the space of the field $`\psi `$ can be expressed in terms of the coefficients $`a_i`$, i.e.,
$$\text{d}[\psi ]=\underset{i}{}\mu \text{d}a_i,$$
(37)
where $`\mu `$ is a normalization constant with the dimension of inverse length, and the indefinite product runs over all values of the label $`i`$.
By substituting the expansions of the field (36) and the measure (37) into the expression for the correlation function (35), the path-integral can be evaluated explicitly. Assuming that there are no nontrivial source terms of the kind which we discussed in the previous section, then the standard expression for the two-point correlation function follows in terms of the complete set of modes $`\psi _i`$. We will not repeat this calculation here, which can be found, e.g., in , but instead we will consider what is the effect on the two-point correlation function (35) when there is a nontrivial source term $`J^+[\psi ]`$ which contributes to the variation of the action.
Let us now define the set of functions $`\stackrel{~}{\psi }_i\{\delta \psi \}_0`$, which satisfy the property that the linear span of the modes $`\stackrel{~}{\psi }_i`$ is dense in $`\{\delta \psi \}_0`$, and the modes $`\stackrel{~}{\psi }_i`$ are chosen so that they are orthonormal with respect to the $`L_2()`$-inner product. We would like to show that the modes $`\stackrel{~}{\psi }_i`$ are complete in $`L_2()`$. Note that the class of functions $`D()`$, which are bounded and which have compact support on $``$, is contained in $`\{\delta \psi \}_0`$. But $`D()`$ is known to be dense in $`L_2()`$ with the $`L_2()`$-norm, and therefore the linear span of the modes $`\stackrel{~}{\psi }_i`$ must be dense in $`L_2()`$. At this point, let us note that the set of functions $`L_2()`$, with the $`L_2()`$-inner product, form a Hilbert space $`H`$. It is a standard result that a set of functions $`\{\psi _i\}`$ is complete in $`H`$ when the linear span of the functions $`\psi _i`$ is dense in $`H`$, and vice-versa (see, e.g., ). This observation implies that the modes $`\stackrel{~}{\psi }_i`$ are complete in $`L_2()`$.
We therefore have two complete and orthonormal sets of functions $`\psi _i`$ and $`\stackrel{~}{\psi }_i`$ in $`L_2()`$, and an arbitrary field perturbation $`\delta \psi L_2()`$ can be expressed in terms of the modes $`\stackrel{~}{\psi }_i`$, i.e.,
$$\delta \psi :=\psi \psi _0=\underset{i}{}\stackrel{~}{a}_i\stackrel{~}{\psi }_i.$$
(38)
It is simple to show that the transformation which expresses one set of basis functions in terms of the other must be orthogonal. Let us now express the measure $`\text{d}[\psi ]`$, given by expression (37), in terms of the new set of modes $`\stackrel{~}{\psi }_i`$. We obtain,
$$\text{d}[\psi ]=\underset{i}{}\mu \text{d}\stackrel{~}{a}_i,$$
(39)
where we used that the Jacobian of the transformation relating the coefficients $`a_i`$ and $`\stackrel{~}{a}_i`$ equals one when the transformation is orthogonal.
One could expect that the path-integral, evaluated with the measures (37) and (39), gives rise to the same result, since all we have done is to express one complete basis of modes in terms of the other. This observation is not correct. Note that when the path-integral (35) is performed with the measure (39), then the source term $`J^+[\psi ]`$ vanishes trivially, since the argument $`\psi `$ is a linear combination of the modes $`\stackrel{~}{\psi }_i`$, and therefore $`\psi \{\delta \psi \}_0`$. On the contrary, when the path-integral is performed with the measure (37), then $`\psi `$ is a linear combination of the modes $`\psi _i`$, and $`J^+[\psi ]`$ will generally be nonzero, which follows from the observation that $`J^+[\psi _i]`$ diverges for all $`\psi _i`$, as we showed in the previous section.
Let us try to make precise in which sense the expansion of the field in terms of two complete sets of modes (36) and (38) differs. Since both expansions converge to the same limit $`\delta \psi `$, it follows that the difference between the two expansions can only be a configuration with zero $`L_2()`$-norm. When performing the path-integral (35), using the measures (37) and (39) respectively, we are integrating over paths in $`L_2()`$ which may differ by a zero-norm configuration. These zero-norm configurations are precisely the degrees of freedom which give rise to the nontrivial source term $`J^+[\psi ]`$. In order to show this, let us recall that $`J^+[\psi _i]`$ diverges for all $`\psi _i`$. Since $`J^+[\psi ]`$ is linear in $`\psi `$, and $`J^+[\stackrel{~}{\psi }]=0`$ when $`\stackrel{~}{\psi }`$ is in the linear span of the modes $`\stackrel{~}{\psi }_i`$, it follows that
$$J^+[\psi _i]=J^+[\psi _iP\psi _i],$$
(40)
where $`P\psi _i`$ denotes the projection of $`\psi _i`$ onto the basis of modes $`\stackrel{~}{\psi }_i`$, i.e.,
$$P\psi _i:=\underset{j}{}\stackrel{~}{\psi }_j\psi _i\stackrel{~}{\psi }_j.$$
(41)
But the modes $`\stackrel{~}{\psi }_i`$ where found to be complete in $`L_2()`$, so that $`(1P)\psi _i`$ must have zero $`L_2()`$-norm. The argument of $`J^+`$ on the right-hand side of equation (40) has therefore zero $`L_2()`$-norm, and therefore this must be the degree of freedom which causes the divergence of the source term. Since the action functional depends on zero-norm degrees of freedom through the term $`J^+[\psi ]`$, the expression for the correlation function (35) is under-determined. Recall that the same ambiguity was present when we tried to determine the extremal-action configurations in section 5.2. Although we do not know of a way to resolve this ambiguity, let us consider two different approaches which might work.
First, one can fix the zero-norm degrees of freedom on the basis of a physical or philosophical argument. In practice, this could mean that one sets the source term $`J^+`$ equal to zero by restricting the phase-space of the $`\psi `$-field to a dense subset of $`L_2()`$ for which $`J^+`$ vanishes. In order to make this approach better than just guessing, one needs to establish whether specific restrictions on the phase-space of the $`\psi `$-field lead to different predictions, which can be falsified.
As a different approach, one could change the measure on the space of the $`\psi `$-field in order to accommodate the zero-norm degrees of freedom. Again, the problem is that there is no clear guideline for doing so, unless one can show that different choices of measure lead to different observable predictions.
It is illustrative to consider a similar ambiguity which occurs in the definition of the path-integral, when one is dealing with fluctuations at infinitesimal rather than infinite length scales. This ambiguity is related to the fact that typical paths which contribute to the path-integral are non-differentiable. Since the class of smooth paths ($`C^{\mathrm{}}`$) is dense in the class of continuous paths ($`C^0`$), the difference between a path in $`C^0`$ and the nearest path in $`C^{\mathrm{}}`$ must have zero $`L_2()`$-norm. As we have seen, the measure (37) does not accommodate these degrees of freedom, and the formal expression is ambiguous on the point of the differentiability of the paths over which we integrate. The action functional is however sensitive to the degree of differentiability of the paths, which is made clear by the fact that the action is generally finite for differentiable paths and infinite for non-differentiable paths. One could try to resolve this ambiguity by simply considering paths in $`C^{\mathrm{}}`$, so that the action functional is well defined, but in this case one can show that the field operators in expression (35) commute trivially, and one does not recover quantum physics .
Finally, let us note that similar implications hold for other field theories which are described by an action functional which is non-linear in the field variable. In particular, it is well known that the Einstein field equations can be derived by varying an action functional, which is given by
$$S[g_{\mu \nu }]=\frac{1}{16\pi G}_{}R\sqrt{g},$$
(42)
where $`R`$ denotes the Ricci scalar, and we have ommited a possible contribution from matter fields and a cosmological constant. Similar to the case where we considered a scalar field, a contribution of a surface term to the variation of the action does occur. At first-order in the metric perturbation, the contribution of this surface term is given by ,
$$\delta S[g_{\mu \nu }]=2_{}\left(\delta K+n^ah^{bc}\delta g_{ab;c}\right)\text{d}\mathrm{\Omega },$$
(43)
where $`\delta K`$ denotes the variation of the trace of the extrinsic curvature at the boundary $``$, while $`h^{bc}`$ and $`n^a`$ denote the induced three-metric and the normal to the boundary respectively, and $`\text{d}\mathrm{\Omega }`$ denotes the volume element on $``$. The first term on the right-hand side of equation (43) can be canceled by adding a surface integral of two times the extrinsic curvature $`K`$ to the action functional (42) (see also the discussion in section 3). The second term on the right-hand side of equation (43) vanishes when it is evaluated according to a classical variational approach where we set $`\delta g_{ab}`$ equal to zero at the bounday $``$, but this term could be of interest in cosmological situations when we do not require that perturbations vanish outside a finite volume.
## 7 Conclusion
We revisited the variational principle in a cosmological context. Starting from the path-integral formulation of quantum physics, we argued that there is a correspondence between classical physics and extremal action fields. The phase-space in which extremality of the action is considered, is not constrained in quantum physics, and we showed that there can be a non-trivial contribution arising from surface terms. We made this problem explicit by considering a scalar field in a perturbed open FLRW spacetime. In the case of an open FLRW spacetime with a Gaussian spectrum of subcurvature perturbations, we found no non-trivial correction to the classical equation of motion. In the case where supercurvature perturbations are present, extremality of the action could only be defined after adopting additional restrictions on the phase-space of the scalar field, but the corresponding equations of motion are ambiguous since they depend on how one restricts the phase-space of the field. We showed that the restricted phase-spaces which yield different physical results, differ by perturbations with vanishing $`L_2`$-norm. This ambiguity is present both at a classical level and a quantum level. We briefly discussed a possible strategy to resolve the ambiguity which is due to perturbations with vanishing $`L_2`$-norm.
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# SPECTRA OF DOUBLY HEAVY QUARK BARYONS
## Introduction
The investigation of properties of hadrons containing one or more heavy quarks is very important for understanding the dynamics of quark and gluon interactions. Presently at the LHC, B-factories and the Tevatron with high luminosity, several experiments have been proposed, in which a detailed study of baryons containing two heavy quarks can be performed. In particular, in the forthcoming experiment at CERN, the COMPASS group is going to find the doubly charmed baryons and study their physical properties. In this connection, doubly heavy quark baryons are now becoming one of the most exciting subjects in particle physics. Therefore, theoretical predictions of properties of doubly heavy quark baryons acquire a big significance for the forthcoming experimental study of these particles. So far there have been various approaches by which their mass spectra and other properties can be calculated. One of them is the nonrelativistic quark model which gives relatively accurate results for baryon spectra . The possible quark compositions of doubly heavy quark baryons are $`ccq`$, $`cbq`$ and $`bbq`$, where $`q`$ denotes a light $`u`$, $`d`$ or $`s`$ quark. Note that the baryons containing the top quark(s) are not practical subject here because a top quark is extremely heavy and hence we have no chance to find them as stable hadrons. The doubly heavy quark baryons may be considered as an analogue of the hydrogen molecular ion $`H_2^+`$, which has been treated successfully in the Born-Oppenheimer approximation. The same approximation is expected to be efficient even for doubly heavy quark baryons, though there exist some differences between these baryons and $`H_2^+`$ systems. One of them is, for this case, the appearance of the confining potential in addition to the QCD Coulomb potential. As is well known, the variables of the Schrödinger equation with two-center Coulomb plus confining potential cannot be separated for their kinematical variables, in general. To our knowledge, the two-center potential which allows the separation of variables is only the two-center Coulomb plus harmonic oscillator potential.
In this paper, we treat $`QQq`$ baryons in the nonrelativistic approach by using the solution of the Schrödinger equation with two-center Coulomb plus harmonic oscillator potential, i. e. the well-known method of ethalon equation which is widely used for solving Schrödinger equation with two-center pure Coulomb potential in the physics of $`H_2^+`$ . First, we give a general scheme of treatment $`QQq`$ baryon in the Born-Oppenheimer approximation. Then, the two-center Schrödinger equation with two-center Coulomb plus harmonic oscillator potential is analytically solved with some approximation: the energy term of the light quark moving in the field of two heavy quarks is obtained in the form of asymptotical expansion over the inverse power of the distance between heavy quarks. Finally, we give an analytical formula of the baryon energy spectrum for $`QQq`$.
## Doubly heavy quark baryon in the Born-Oppenheimer approximation
In the Born-Oppenheimer approximation the wave function is split into heavy- and light-quark degrees of freedom
$$\mathrm{\Psi }(R,r)=\underset{n}{}\varphi _n\left(R\right)\psi _n(R,r),$$
where $`R`$ is the distance between two heavy quarks and $`r`$ is the distance between light quark and center-of-mass of the heavy-quark pair. The light quark wave function $`\psi (r,R)`$ and its energy term $`E\left(R\right)`$ can be found from the Schrödinger equation
$$\left[\frac{1}{2m_q}\mathrm{\Delta }+V\left(r_1\right)+V\left(r_2\right)\right]\psi =E\left(R\right)\psi ,$$
where $`r_1`$ and $`r_2`$ are the distances between light and heavy quarks, $`Q_1`$ and $`Q_2`$, respectively. The binding energy of this system is approximated by the equation:
$$\left[\frac{1}{2\overline{M}_{QQ}}\mathrm{\Delta }+V_{QQ}\left(R\right)+E\left(R\right)\right]\varphi =\epsilon \varphi ,$$
where $`\overline{M}_{QQ}`$ is the reduced mass of $`QQ`$.
A quark potential with Coulomb plus harmonic confinement for this baryon is given by
$$V\left(r_{ij}\right)=_{i,j}\frac{1}{4}\lambda _i\lambda _j\left(V_0Ar_{ij}^2+\frac{\alpha _s}{r_{ij}}\right)=\frac{2}{3}_{i,j}\left(V_0Ar_{ij}^2+\frac{\alpha _s}{r_{ij}}\right)$$
In the field of two heavy quarks with this potential, the motion of a light quark can be nonrelativistically described by the following Schrödinger equation,
$$\left[\frac{1}{2}\mathrm{\Delta }\frac{Z}{r_1}\frac{Z}{r_2}+\omega ^2\left(r_1^2+r_2^2\right)\frac{4}{3}V_0\right]\psi =E\left(R\right)\psi ,$$
(1)
where $`Z=2\alpha _s/3`$ and $`\omega ^2=2A/3`$.
In the prolate spheroidal coordinates defined as
$$\xi =\frac{r_1+r_2}{R}\left(1<\xi <\mathrm{}\right),\eta =\frac{r_1r_2}{R}\left(1<\eta <1\right),$$
the potential term in eq.(1) can be written in the form
$$V(r_1,r_2)=\frac{2}{R^2}\frac{a\left(\xi \right)+b\left(\eta \right)}{\xi ^2\eta ^2}+\frac{\omega ^2R^2}{2}\frac{4}{3}V_0,$$
(2)
where
$$a\left(\xi \right)=2ZR\frac{\omega ^2R^4}{4}\xi ^2\left(\xi ^21\right),b\left(\eta \right)=2ZR\frac{\omega ^2R^4}{4}\eta ^2\left(\eta ^21\right).$$
As is well known, the Schrödinger equation with the potential in the form of eq.(2) is separable in the prolate spheroidal coordinates. Then it is convenient to use
$$\psi =\frac{U\left(\xi \right)}{\sqrt{\xi ^21}}\frac{V\left(\eta \right)}{\sqrt{1\eta ^2}}\frac{e^{\pm im\varphi }}{\sqrt{2\pi }},$$
where $`\varphi `$ and $`m`$ are azumuthal angle and azimuthal quantum number, respectively. After substituting this into eq.(1), we obtain from the following ordinary differential equations connected with separation constants $`\lambda `$ and $`m`$:
$$U^{\prime \prime }\left(\xi \right)+\left[\frac{h^2}{4}+\frac{h\left(\alpha \xi \lambda \right)}{\xi ^21}h^4\gamma \xi ^2+\frac{1m^2}{\left(\xi ^21\right)^2}\right]U\left(\xi \right)=0,$$
(3)
$$V^{\prime \prime }\left(\eta \right)+\left[\frac{h^2}{4}+\frac{h\lambda }{1\eta ^2}h^4\gamma \eta ^2+\frac{1m^2}{\left(1\eta ^2\right)^2}\right]V\left(\eta \right)=0,$$
(4)
where $`\alpha =2Z/\sqrt{2E^{}}`$ and $`\gamma =\omega ^2/8E^2`$, and further
$$h=\sqrt{2E^{}}R,$$
(5)
with
$$E^{}=E\frac{\omega ^2R^2}{2}+\frac{4}{3}V_0.$$
Finiteness and continuity of the wave function $`\psi `$ in the whole space lead to the following boundary conditions for the functions $`U`$ and $`V`$:
$$U\left(\xi \right)_{\xi =1}=0,U\left(\xi \right)_\xi \mathrm{}0,$$
(6)
$$V\left(\eta \right)_{\eta =\pm 1}=0.$$
(7)
## Asymptotics of quasi-angular equation
We will approximately solve eqs.(3) and(4) for large $`R`$ by the method of ethalon equation. This method is successfully applied to the solution of nonrelativistic two center Coulomb problem and in the theory of diffraction of waves. Details on the method of ethalon equation are given in , also briefly described in Appendix.
Let us start from the angular equation of (4). As an ethalon equation for eq.(4), we choose the following Whittaker equation :
$$W^{\prime \prime }+\left[\frac{h^4}{4}+\frac{h^2k}{z}+\frac{1m^2}{4z^2}\right]W=0$$
(8)
and seek a solution in the form
$$V=\left[z^{}\left(\eta \right)\right]^{\frac{1}{2}}M_{k,\frac{m}{2}}\left(h^2z\right),$$
(9)
where $`M_{k,\frac{m}{2}}\left(h^2z\right)`$ is the solution (regular at zero) of eq.(8). Substituting (9) into (4) and taking into account (8), we get the following equation for $`z`$:
$`{\displaystyle \frac{z^2}{4}}\gamma \left(x1\right)^2{\displaystyle \frac{1}{h^2}}\left({\displaystyle \frac{1}{4}}+{\displaystyle \frac{kz^2}{z}}{\displaystyle \frac{\lambda }{2x\left(1x/2\right)}}\right)+`$
$`{\displaystyle \frac{\tau }{h^2}}\left({\displaystyle \frac{1}{x^2\left(1x^2\right)}}{\displaystyle \frac{z^2}{z^2}}\right){\displaystyle \frac{1}{2h^2}}\{z,x\}=0,`$ (10)
where
$$\{z,x\}=\frac{3}{2}\left(\frac{z^{\prime \prime }}{z^{}}\right)^2+\frac{z^{\prime \prime \prime }}{z^{}},$$
and $`\tau =\frac{1m^2}{4},`$ $`x=1+\eta `$.
Requirement of coincidence at the transition points
$$z\left(x\right)_{x=0}=0,$$
leads to the following ”quantum condition”
$$\lambda =2kz^{}\left(0\right)+\frac{2\tau }{h^2}\left[\frac{z^{\prime \prime }\left(0\right)}{z^{}\left(0\right)}1\right].$$
(11)
We will seek the solution of eq.(10) and eigenvalues $`\lambda `$ in the form of following asymptotical expansion:
$$z=\underset{k=0}{\overset{\mathrm{}}{}}\frac{z_k}{h^k},\lambda =\underset{k=0}{\overset{\mathrm{}}{}}\frac{\lambda _k}{h^k}.$$
Substitution of these expansions into (10) gives us the recurrence system of differentil equations for $`z`$ :
$$z_0^{}=2\gamma ^{\frac{1}{2}}\left(x1\right),$$
$$z_1^{}=0,$$
$$z_2^{}=\frac{1}{2z_0^{}}+\frac{2kz_0^{}}{z_0}\frac{\left(z_1^{}\right)^2}{2z_0^{}}\frac{2\lambda _0}{z_0^{}x\left(1x/2\right)}\frac{z_2^2}{2},$$
$$..................,$$
and for $`\lambda `$ :
$$\lambda _0=2kz_0^{}\left(0\right),$$
$$\lambda _1=2kz_1^{}\left(0\right),$$
$$\lambda _2=2kz_2^{}\left(0\right)+2\tau \left(\frac{z_0^{\prime \prime }\left(0\right)}{z_1^{}\left(0\right)}1\right),$$
$$................$$
Solving these recurrence equations, we obtain
$$\lambda ^{\left(\eta \right)}=4k\gamma ^{\frac{1}{2}}+\frac{2k\beta 4\tau }{h^2}+O\left(\frac{1}{h^4}\right),$$
(12)
for $`\lambda `$, and
$$z=\gamma ^{\frac{1}{2}}x\left(2x\right)+\frac{1}{h^2}\beta ln\left(1x\right)+O\left(\frac{1}{h^4}\right),$$
(13)
for $`z`$.
¿From boundary conditions one can obtain for quantum number $`k`$:
$$k=q+\frac{m+1}{2},$$
where $`q=0,1,2,\mathrm{}.`$.
## Asymptotics of quasi-radial equation
As an ethalon equation for eq.(3), we take the following equation :
$$W^{\prime \prime }+\left[h^2sh^4y^2\frac{4\tau +3}{4y^2}\right]W=0,$$
(14)
a solution of which is expressed by the confluent hypergeometric functions
$$W=y^ce^{\frac{h^4y^2}{2}}F(\frac{s2c1}{4},c+\frac{1}{2},h^4y^2),$$
where $`c=\frac{1+\sqrt{m^2+3}}{2}`$.
Boundary condition (6) and properties of functions $`F`$ give rise to the following expression for $`s`$:
$$s=4n+\sqrt{m^2+3}+2.$$
Substituting
$$U=\left[y\left(\xi \right)\right]^{\frac{1}{2}}W\left(y\left(\xi \right)\right)$$
into eq.(3), we obtain
$$\frac{y^2y^2}{4}\gamma \xi ^2+\frac{1}{h^2}\left(\frac{1}{4}sy^2\frac{\lambda }{\xi ^21}\right)+\frac{1}{h^3}\frac{\alpha \xi }{\xi ^21}+\frac{4\tau }{h^4\left(\xi ^21\right)^2}\frac{34\tau }{4h^4}\frac{y^2}{y^2}\frac{1}{2h^4}\{y,\xi \}=0.$$
(15)
After substitution of
$$\varphi =\frac{y^2\left(t\right)}{4},$$
this equation can be reduced to the form
$$\varphi ^2\gamma \left(t+1\right)+\frac{1}{h^2}\left(\frac{1}{4}\left(n+\frac{1}{2}\right)\frac{\varphi ^2}{\varphi }\frac{\lambda }{t\left(t+2\right)}\right)+\frac{1}{h^3}\frac{\alpha \left(t+1\right)}{t\left(t+2\right)}+\frac{\tau }{h^4}\left(\frac{\varphi ^2}{\varphi ^2}\frac{4}{t^2\left(t+2\right)^2}\right)[\varphi ,t]=0,$$
where $`t=\xi 1.`$ The quantization condition which follows from $`\varphi \left(x\right)=0`$ is written in the form
$$\lambda =2s\varphi ^{}\left(0\right)+\frac{\alpha }{h}\frac{1}{h^2}\left[\frac{\varphi ^{\prime \prime }}{\varphi ^{}}+1\right]_{t=0}.$$
(16)
Inserting the asymptotical expansions
$$\varphi =\underset{k=0}{\overset{\mathrm{}}{}}\frac{\varphi _k}{h^k},\lambda =\underset{k=0}{\overset{\mathrm{}}{}}\frac{\lambda _k}{h^k}$$
into eq.(Asymptotics of quasi-radial equation) and solving the equations obtained herewith, we get the following result
$`y=2\gamma ^{\frac{1}{4}}\left(t^2+2t\right)^{\frac{1}{2}}+{\displaystyle \frac{1}{h^2}}\delta \gamma ^{\frac{1}{4}}\left(t^2+2t\right)^{\frac{1}{2}}ln\left(t+1\right)+`$
$`{\displaystyle \frac{1}{h^3}}\alpha \gamma ^{\frac{3}{4}}\left(t^2+2t\right)^{\frac{1}{2}}ln{\displaystyle \frac{2\left(t+1\right)}{t+1}}+O\left({\displaystyle \frac{1}{h^4}}\right),`$ (17)
for $`y`$, and
$$\lambda ^{\left(\xi \right)}=2s\gamma ^{\frac{1}{2}}\frac{\alpha }{h}+\frac{4\tau s\delta }{h^2}\frac{s\alpha \gamma ^{\frac{1}{4}}}{2h^3}+O\left(\frac{1}{h^4}\right),$$
(18)
for $`\lambda `$.
## Asymptotical expansion for energy
Asymptotical expansions (12) and (18) give us an expression for the energy term in the form of multipole expansion. In order to obtain this expansion one should insert
$$E^{}=E_0+\frac{E_1}{R}+\frac{E_2}{R^2}+\mathrm{}$$
into eqs.(12) and (18). Equating $`\lambda ^{\left(\eta \right)}`$ to be $`\lambda ^{\left(\xi \right)}`$ and taking into account (5), we get the following equations for coeficients $`E_1,E_2,\mathrm{}`$ :
$$E_1=\frac{1}{6Z}\left[\left(s\omega 2k\omega ^1\right)\left(2E_0\right)^{\frac{5}{2}}+\left(4s^216k^216\tau \right)\left(2E_0\right)^{\frac{3}{2}}\right],$$
$$E_2=\frac{5}{2}E_1^2+2s\omega ^1E_0+E_1\left(2E_0\right)^{\frac{1}{2}}Z^1\left(16\tau ^2+16k^24s^2\right),$$
$$..............$$
Now we need to find $`E_0`$. In order to find this value, we note that for $`R\mathrm{},`$ $`E^{}=E_0`$ and hence we have
$$E=E_0+\frac{\omega ^2R^2}{2}+\frac{4}{3}V_0.$$
(19)
On the other hand, for large $`R`$ we have
$`V(r_1,r_2)={\displaystyle \frac{2Z}{R}}{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}\left({\displaystyle \frac{r}{R}}\right)^lP_l\left(cos\theta \right)+\omega ^2[(r^2+2rRcos\theta +{\displaystyle \frac{R^2}{4}})+`$
$`(r^22rRcos\theta +{\displaystyle \frac{R^2}{4}})]\omega ^2(2r^2+{\displaystyle \frac{R^2}{2}}){\displaystyle \frac{4}{3}}V_0.`$ (20)
Hence, for the energy term with this potential we obtain
$$E=2\omega \left(N+\frac{3}{2}\right)+\frac{\omega ^2R^2}{2}+\frac{4}{3}V_0,$$
(21)
where $`N=n+q+m+1`$ is the principial quantum number. Comparing eqs.(19) and (21), we obtain
$$E_0=2\omega \left(N+\frac{3}{2}\right).$$
Thus, the following asymptotical expansion is obtained for the energy term of light quark in the field of two heavy quarks:
$$E=\frac{4}{3}V_0+\frac{\omega ^2R^2}{2}+E_0+\frac{E_1}{R}+\frac{E_2}{R^2}+\mathrm{}$$
## $`QQq`$ baryon spectra
As mentioned above, the $`QQq`$ binding energy can be finally obtained by solving the Schrödinger equation
$$\left[\frac{1}{2\overline{M}_{QQ}}\mathrm{\Delta }+V_{QQ}\left(R\right)+E\left(R\right)\right]\varphi =\epsilon \varphi .$$
(22)
If one takes $`E\left(R\right)`$ in the form
$$E=\frac{4}{3}V_0+\frac{\omega ^2R^2}{2}+E_0+\frac{E_1}{R},$$
for
$$V_{QQ}\left(R\right)=\omega ^2R^2\frac{Z}{R}\frac{2}{3}V_0,$$
then eq.(22) can be rewritten as
$$\left[\frac{1}{2\overline{M}_{QQ}}\mathrm{\Delta }+\omega ^2R^2\frac{Z^{}}{R}V_0^{}\right]\varphi =\epsilon \varphi ,$$
(23)
where $`Z^{}=ZE_1,`$ $`\omega ^2=\frac{3}{2}\omega ^2`$ and $`V_0^{}=2V_0E_0`$.
To solve this equation, we use the result of where a method for an analytical solution of the Schrödinger equation with potential
$$V\left(R\right)=\frac{Z}{R}+\lambda R^k$$
was offered. Details on this method and its application to our potential are given in Appendix. Application of this method to eq. (23) gives us
$$\epsilon _{Nnl}=2\omega \left(N+\frac{3}{2}\right)+\left[Z^2\omega ^6r_{nl}\right]^{1/5}2V_0,$$
(24)
where $`N`$ is the principial quantum number of the light quark moving in the field of $`QQ`$, $`r_{nl}`$ is defined in Appendix. The formula (24) describes the energy spectrum of the $`QQq`$ baryon. In tables 1, 2 and 3, the mass spectra of $`ccq`$, $`bbq`$ and $`bcq`$ baryons calculated using the formula (24) are given, respectively. The following values of potential parameters are choosen in this calculation $`\alpha _s=0.39,`$ $`\omega ^2=0.174GeV^3,`$ $`V_0=0.05GeV`$ for the potential
$$V=\frac{2}{3}\left(\frac{\alpha _s}{r}+\omega ^2rV_0\right).$$
## Conclusion
In this work we have treated doubly heavy baryons in the Born-Oppenheimer approximation. The following two problems have been solved in the framework of this approximation: (1)the Schrödinger equation for two-center Coulomb plus harmonic oscillator potential and (2)the Schrödinger equation for central symmetric Coulomb plus harmonic oscillator potential. As the final result an analytical formula for the energy spectrum of baryons containing two heavy quarks is derived. Obtained formula is applied for the calculation of mass spectra of doubly heavy quark baryons with various quark compositions. The above analytical results could be useful for further numerical calculations in non-asymptotical region.
## Appendix
## 1 The scaling variational method and its application to Coulomb plus confining potential
Consider the following Hamiltonian
$$H=\frac{1}{2}\mathrm{\Delta }+V\left(r\right),$$
(25)
which obeys the eigenvalue equation
$$H\psi _{nl}=E_{nl}\psi _{nl},<\psi _{nl}\psi _{n^{}l^{}}>=\delta _{nn^{}}\delta _{ll^{}},$$
(26)
$$n,n^{}=1,2,..,l,l^{}=1,2,..,$$
where $`n`$ and $`l`$ denote the principal and angular quantum numbers, respectively. To solve this equation, we start from a set of functions $`\left\{\varphi _{nl}\right\}`$ which are the eigenfunctions of an arbitrary central field Hamiltonian $`H_0`$:
$$H_0\varphi _{nl}=ϵ_{nl}\varphi _{nl},<\varphi _{nl}\varphi _{n^{}l^{}}>=\delta _{nn^{}}\delta _{ll^{}}.$$
(27)
Then, we construct the functionals
$$\epsilon _{nl}\left(\alpha \right)=<\varphi _{nl}^\alpha H\varphi _{nl}^\alpha >,$$
(28)
where
$$\varphi _{nl}^\alpha =\alpha ^{3/2}\varphi _{nl}\left(\alpha r\right).$$
The value for the $`\alpha `$ is determined by
$$\left(\frac{\epsilon }{\alpha }\right)\left(\alpha =a\right)=0.$$
(29)
Let us now apply this method to our problem, i.e. to the Schrödinger equation with potential
$$V\left(r\right)=\frac{Z^{}}{R}+\omega ^2r^2.$$
As $`H_0`$, we choose pure Coulomb Hamiltonian, i.e.
$$H_0=\frac{1}{2}\mathrm{\Delta }\frac{Z^{}}{R}.$$
Then, $`ϵ_{nl}=\frac{Z^2}{2n^2}`$ with $`n=n_r+l+1`$, where $`n_r`$ is the radial quantum number.
According to the above procedure, we have
$`\epsilon _{nl}\left(\alpha \right)=<\varphi _{nl}^\alpha H_0\varphi _{nl}^\alpha >+<\varphi _{nl}^\alpha \omega ^2r^2\varphi _{nl}^\alpha >=`$ (30)
$`{\displaystyle \frac{Z^2\alpha ^3}{2n^2}}+{\displaystyle \frac{\omega ^2\alpha ^2n^2}{2}}\left[5n^2+13l\left(l+1\right)\right],`$
For calculation of the second matrix element, we have used the well-known expression for average value $`\overline{r}^2`$ in Coulomb field which is given in . ¿From $`\epsilon \left(\alpha \right)/\alpha =0`$, we obtain
$$\alpha _0=\left\{\frac{2\omega ^2n^4}{3Z^2}\left[5n^2+13l\left(l+1\right)\right]\right\}^{1/5}.$$
Then, for the energy level one can get the following analytical formula
$$\epsilon _{nl}=\left[Z^2\omega ^6r_{nl}\right]^{1/5},$$
(31)
where
$$r_{nl}=n^2\left[5n^2+13l\left(l+1\right)\right]^3.$$
## 2 The formal procedure of the method of ethalon equation
Let’s consider the following second order differential equation:
$$y^{\prime \prime }\left(x\right)+p^2\left[\lambda q\left(x\right)\right]y\left(x\right)=0$$
(32)
in the interval $`[a,b]`$. Let in this interval eq.(32) has one transition point(poles and zeros of function $`Q(x,\lambda )=q\left(x\right)\lambda `$ are called transition points of this equation).
Equation
$$w^{\prime \prime }\left(z\right)p^2R\left(z\right)w\left(z\right)=0$$
(33)
which has the same or close transition points as eq.(33) is called the ethalon equation for eq.(32).
Solution of eq.(32) we will seek in the form
$$y\left(x\right)=\left[z^{}(x,p)\right]^{1/2}w\left(z(x,p)\right)$$
(34)
where $`w`$ is the solution of eq.(33). Inserting (34) in to eq.(32) and taking into account eq.(33) we obtain the following (nonlinear) differential equation for $`z(x,p)`$:
$$R\left(z\right)z^2Q(x,\lambda )\frac{1}{2p^2}\{z,x\}=0$$
(35)
here
$$\{z,x\}=\frac{3}{2}\left(\frac{z^{\prime \prime }}{z^{}}\right)^2+\frac{z^{\prime \prime \prime }}{z^{}},$$
In the case of eq. (4) we have for $`Q`$
$$Q=\left[\frac{h^2}{4}+\frac{h\lambda }{1\eta ^2}h^4\gamma \eta ^2+\frac{1m^2}{\left(1\eta ^2\right)^2}\right]$$
and for $`R`$ (from eq.(8)
$$R=\left[\frac{h^4}{4}+\frac{h^2k}{z}+\frac{1m^2}{4z^2}\right].$$
So, for $`z`$ one obtains eq.(10).
$$\frac{z^2}{4}\gamma \left(x1\right)^2\frac{1}{h^2}\left(\frac{1}{4}+\frac{kz^2}{z}\frac{\lambda }{2x\left(1x/2\right)}\right)+\frac{\tau }{h^2}\left(\frac{1}{x^2\left(1x^2\right)}\frac{z^2}{z^2}\right)\frac{1}{2h^2}\{z,x\}=0$$
Table 1. The mass spectrum of $`ccq`$ baryon (in GeV) calculated using the formula (24); $`m_q=0.385`$ GeV, $`m_c=1.486`$ GeV, $`n_l`$ and $`n_d`$ are the principial quantum numbers of light quark and $`cc`$ diquark, respectively, $`L`$ is the orbital quantum number of $`cc`$ diquark.
| | $`n_l,n_d,L`$ | mass | $`n_l,n_d,L`$ | mass | $`n_l,n_d,L`$ | mass | $`n_l,n_d,L`$ | mass |
| --- | --- | --- | --- | --- | --- | --- | --- | --- |
| | 1,1,0 | 3.661 | 1,1,1 | 3.613 | 1,2,2 | 3.649 | 1,3,3 | 3.694 |
| | 1,2,0 | 3.730 | 1,2,1 | 3.708 | 1,3,2 | 3.764 | 1,4,3 | 3.825 |
| | 1,3,0 | 3.816 | 1,3,1 | 3.799 | 1,4,2 | 3.872 | 1,5,3 | 3.949 |
| | 1,4,0 | 3.914 | 1,4,1 | 3.901 | 1,5,2 | 3.988 | 1,6,3 | 4.077 |
| | 1,5,0 | 4.024 | 1,5,1 | 4.012 | 1,6,2 | 4.110 | 1,7,3 | 4.210 |
| | 2,1,0 | 3.839 | 2,1,1 | 3.791 | 2,2,2 | 3.828 | 2,3,3 | 3.873 |
| | 2,2,0 | 3.908 | 2,2,1 | 3.887 | 2,3,2 | 3.942 | 2,4,3 | 4.003 |
| | 2,3,0 | 3.994 | 2,3,1 | 3.978 | 2,4,2 | 4.051 | 2,5,3 | 4.127 |
| | 2,4,0 | 4.093 | 2,4,1 | 4.079 | 2,5,2 | 4.166 | 2,6,3 | 4.255 |
| | 2,5,0 | 4.202 | 2,5,1 | 4.190 | 2,6,2 | 4.289 | 2,7,3 | 4.388 |
| | 3,1,0 | 4.018 | 3,1,1 | 3.970 | 3,2,2 | 4.006 | 3,3,3 | 4.051 |
| | 3,2,0 | 4.086 | 3,2,1 | 4.065 | 3,3,2 | 4.120 | 3,4,3 | 4.182 |
| | 3,3,0 | 4.172 | 3,3,1 | 4.156 | 3,4,2 | 4.229 | 3,5,3 | 4.305 |
| | 3,4,0 | 4.271 | 3,4,1 | 4.257 | 3,5,2 | 4.344 | 3,6,3 | 4.433 |
| | 3,5,0 | 4.380 | 3,5,1 | 4.369 | 3,6,2 | 4.467 | 3,7,3 | 4.567 |
Table 2. The mass spectrum of $`bbq`$ baryon (in GeV) calculated using the formula (24); $`m_q=0.385`$ GeV, $`m_b=4.88`$ GeV, $`n_l`$ and $`n_d`$ principial quantum numbers of light quark and $`bb`$ diquark, respectively, $`L`$ is the orbital quantum number of $`bb`$ diquark.
| | $`n_l,n_d,L`$ | mass | $`n_l,n_d,L`$ | mass | $`n_l,n_d,L`$ | mass | $`n_l,n_d,L`$ | mass |
| --- | --- | --- | --- | --- | --- | --- | --- | --- |
| | 1,1,0 | 9.890 | 1,1,1 | 9.874 | 1,2,2 | 9.886 | 1,3,3 | 9.900 |
| | 1,2,0 | 9.911 | 1,2,1 | 9.905 | 1,3,2 | 9.922 | 1,4,3 | 9.942 |
| | 1,3,0 | 9.939 | 1,3,1 | 9.934 | 1,4,2 | 9.957 | 1,5,3 | 9.981 |
| | 1,4,0 | 9.970 | 1,4,1 | 9.966 | 1,5,2 | 9.993 | 1,6,3 | 10.022 |
| | 1,5,0 | 10.005 | 1,5,1 | 10.001 | 1,6,2 | 10.032 | 1,7,3 | 10.064 |
| | 2,1,0 | 10.096 | 2,1,1 | 10.081 | 2,2,2 | 10.093 | 2,3,3 | 10.107 |
| | 2,2,0 | 10.118 | 2,2,1 | 10.112 | 2,3,2 | 10.129 | 2,4,3 | 10.149 |
| | 2,3,0 | 10.146 | 2,3,1 | 10.141 | 2,4,2 | 10.164 | 2,5,3 | 10.188 |
| | 2,4,0 | 10.177 | 2,4,1 | 10.173 | 2,5,2 | 10.200 | 2,6,3 | 10.229 |
| | 2,5,0 | 10.212 | 2,5,1 | 10.208 | 2,6,2 | 10.239 | 2,7,3 | 10.271 |
| | 3,1,0 | 10.303 | 3,1,1 | 10.288 | 3,2,2 | 10.300 | 3,3,3 | 10.314 |
| | 3,2,0 | 10.325 | 3,2,1 | 10.318 | 3,3,2 | 10.336 | 3,4,3 | 10.356 |
| | 3,3,0 | 10.353 | 3,3,1 | 10.347 | 3,4,2 | 10.371 | 3,5,3 | 10.395 |
| | 3,4,0 | 10.384 | 3,4,1 | 10.380 | 3,5,2 | 10.407 | 3,6,3 | 10.436 |
| | 3,5,0 | 10.419 | 3,5,1 | 10.415 | 3,6,2 | 10.446 | 3,7,3 | 10.478 |
Table 3. The mass spectrum of $`bcq`$ baryon (in GeV) calculated using the formula (24); $`m_q=0.385`$ GeV, $`m_b=4.88`$ GeV, $`m_c=1.486`$ GeV, $`n_l`$ and $`n_d`$ principial quantum numbers of light quark and $`bc`$ diquark, respectively, $`L`$ is the orbital quantum number of $`bc`$ diquark.
| | $`n_l,n_d,L`$ | mass | $`n_l,n_d,L`$ | mass | $`n_l,n_d,L`$ | mass | $`n_l,n_d,L`$ | mass |
| --- | --- | --- | --- | --- | --- | --- | --- | --- |
| | 1,2,0 | 7.217 | 1,2,1 | 7.160 | 1,3,2 | 7.178 | 1,4,3 | 7.199 |
| | 1,3,0 | 7.259 | 1,3,1 | 7.206 | 1,4,2 | 7.233 | 1,5,3 | 7.263 |
| | 1,4,0 | 7.307 | 1,4,1 | 7.251 | 1,5,2 | 7.286 | 1,6,3 | 7.324 |
| | 1,5,0 | 7.361 | 1,5,1 | 7.300 | 1,6,2 | 7.343 | 1,7,3 | 7.386 |
| | 2,1,0 | 7.438 | 2,1,1 | 7.355 | 2,2,2 | 7.403 | 2,3,3 | 7.452 |
| | 2,2,0 | 7.471 | 2,2,1 | 7.414 | 2,3,2 | 7.432 | 2,4,3 | 7.454 |
| | 2,3,0 | 7.513 | 2,3,1 | 7.461 | 2,4,2 | 7.488 | 2,5,3 | 7.518 |
| | 2,4,0 | 7.562 | 2,4,1 | 7.505 | 2,5,2 | 7.541 | 2,6,3 | 7.579 |
| | 2,5,0 | 7.615 | 2,5,1 | 7.555 | 2,6,2 | 7.597 | 2,7,3 | 7.641 |
| | 3,1,0 | 7.692 | 3,1,1 | 7.669 | 3,2,2 | 7.687 | 3,3,3 | 7.709 |
| | 3,2,0 | 7.726 | 3,2,1 | 7.716 | 3,3,2 | 7.743 | 3,4,3 | 7.773 |
| | 3,3,0 | 7.768 | 3,3,1 | 7.760 | 3,4,2 | 7.796 | 3,5,3 | 7.833 |
| | 3,4,0 | 7.816 | 3,4,1 | 7.810 | 3,5,2 | 7.852 | 3,6,3 | 7.896 |
| | 3,5,0 | 7.870 | 3,5,1 | 7.864 | 3,6,2 | 7.912 | 3,7,3 | 7.961 |
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# Acknowledgments
## Acknowledgments
We would like to thank M. Spira and P. Uwer for discussions. Z. G. Si wishes to thank the DESY theory group for its hospitality during the final stages of this work.
FIGURE CAPTIONS
Fig. 1. Dimensionless scaling functions $`g_{q\overline{q}}^{(0)}(\eta )`$ (dotted), $`g_{q\overline{q}}^{(1)}(\eta )`$ (full), and $`\stackrel{~}{g}_{q\overline{q}}^{(1)}(\eta )`$ (dashed) that determine the expectation value $`𝒪_1_{q\overline{q}}`$.
Fig. 2. Same as Fig.1, but for $`𝒪_2_{q\overline{q}}`$.
Fig. 3. Same as Fig.1, but for $`𝒪_3_{q\overline{q}}`$.
Fig. 4. Same as Fig.1, but for $`𝒪_4_{q\overline{q}}`$.
Fig. 5. Same as Fig.1, but for $`𝒪_5_{q\overline{q}}`$.
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# Ultrahigh energy neutrino physics
## 1 Introduction
The ultrahigh energy neutrino physics is an interdisciplinary subject which can address many vital problems in astrophysics, particle physics and geophysics, for a review see . Neutrinos can traverse large distances without being disturbed and thus give us important information about distant astronomical objects like active galactic nuclei or gamma ray bursts. The ultrahigh energy neutrinos can have energies up to $`10^{12}\mathrm{GeV}`$, much larger than currently accessible at present colliders, like HERA. At these ultrahigh energies the structure of nucleon is probed at very small values of Bjorken $`x`$. It means that we are possibly entering a region in which partons have large density. Therefore detailed knowledge of parton distribution function at very small $`x`$ is vital to estimate the neutrino-nucleon cross section. This cross section rises strongly with energy and therefore neutrino propagation through the Earth can be affected by the increased interaction with matter at these ultrahigh energies. Thus it is necessary to consider the effect of attenuation of high energy neutrinos while traversing the Earth. This phenomenon could be used as a possible method of the ”Earth tomography” by neutrinos. All these topics will be studied at large neutrino telescopes like AMANDA, NESTOR and ANTARES (for review see ).
In this paper we present a method of calculation the neutrino-nucleon cross section using gluon and quark distributions obtained from a unifed evolution equation for small $`x`$ . This equation embodies both DGLAP and BFKL evolutions on equal footing as well as important subleading effects via consistency constraint . This constraint limits the virtualities of the exchanged gluon momenta in the ladder to their transverse components. We then extrapolate the results to ultrahigh energies and compare the predictions with the other, based on standard global parton analysis ,.
We use the resulting cross section as an input to the transport equation for the neutrino flux penetrating the Earth. We solve this equation for different incident angles as well as different input neutrino fluxes originating from active galactic nuclei, gamma ray bursts and a sample top-down model.
The details of the calculation has been presented elsewhere .
## 2 Sources of ultrahigh energy neutrinos
Neutrinos have the advantage that they are weakly interacting with matter and they are hardly absorbed when travelling along large distances. On the other hand cosmic rays are being absorbed through the following interactions:
* pair production: $`\gamma _{CR}+\gamma _{BR}e^+e^{}`$
* inverse Compton scattering: $`e_{CR}+\gamma _{BR}e+\gamma `$
* photoproduction of pion: $`p_{CR}+\gamma _{BR}p+N\pi `$
* nuclei fragmentation by photo-pion interactions
These reactions cause that the spectrum of the cosmic rays has the GZK cutoff around $`10^{19}\mathrm{eV}`$. Therefore neutrinos are the best candidates for supplying information about distant objects. The possible sources of highly energetic neutrinos are:
* active galactic nuclei
* gamma ray bursts
* top-down models
Active galactic nuclei are the most powerful sources of radiation in the whole Universe. The engine of the AGN is the supermassive black hole with mass bilion times larger than the mass of the sun. Surrounding the black hole is the accretion disc usually accompanied by two jets. The spectrum of the emitted photons spreads from radio waves to TeV energies. The jets are the sources of the most energetic gamma rays, since the particles (electrons and possibly protons) are accelerated in blobs along the jets with Lorenz factor $`\gamma 10`$. Electrons loose energy via synchrotron radiation thus producing very energetic photons. If protons are also accelerated then they can interact with the ambient photons producing pions. Consequently strong flux of neutrinos will be produced .
Gamma ray bursts are also very probable sources of neutrinos. Although the underlying event of GRB’s is not entirely known the present knowledge is consistent that the bursts are produced as a relativistically expanding fireball which intitial radius was around $`100\mathrm{k}\mathrm{m}`$. The original state which is opaque to light expands in a relativistic shock with $`\gamma 300`$ to the point where it becomes optically thin and produces intensive gamma ray spectrum. As in the case of AGN, the acceleration of protons can result in the production of neutrinos because they will photoproduce pions and in the end neutrinos will emerge .
Top-down models are the most speculative scenarios of producing neutrino fluxes at ultrahigh energies. In these models one assumes that the particles are not accelerated but are rather produced as a result of the decay of the supermassive particles $`M_x10^{14}10^{16}\mathrm{GeV}`$. These particles could emerge as a decay of some topological defects: superconducting strings or magnetic monopoles. The top-down models produce quite hard spectrum of neutrinos extending beyond $`10^{12}\mathrm{GeV}`$ and they were proposed as a possible solution to the GZK cutoff .
## 3 Neutrino - nucleon cross section
The dominant interaction for neutrinos is the neutrino-nucleon interaction. It has the largest value of the cross section, and is dominating over the interaction with electrons. There is however one exception: resonant $`W`$ production in $`\overline{\nu _e}e^{}`$ interaction at $`E_\nu =6.4\times 10^5\mathrm{GeV}`$. At this energy this process domintates by 2 orders of magnitude over other contributions. The neutrino - nucleon interaction can be visualised in Fig. 1.
Here $`l`$ is the incoming neutrino of four momentum $`k`$ and $`l^{}`$ is the outgoing neutrino or charged lepton with four momentum $`k^{}`$. $`N`$ is the target nucleon with four momentum $`p`$ and $`X`$ is the arbitrary hadronic final state. This deep inelastic scattering process can be described using two standard variables: $`q^2=Q^2<0`$ is the four momentum squared of the exchanged vector boson, and $`x=\frac{Q^2}{2pq}`$ is the standard Bjorken scaling variable. At high energies (up to $`10^{12}GeV`$) , very small values of $`x`$ can be probed because,
$$x\frac{M_V^2}{2M_N\nu }10^8$$
(1)
where $`M_V`$ is the mass of the heavy vector boson, $`M_N`$ is the nucleon mass and $`\nu `$ is the energy of the exchanged vector boson. These are values of $`x`$ which are not accessible at present electron-proton colliders (for example HERA goes down to $`10^5`$ at fairly low $`Q^2`$). Therefore a detailed knowledge of parton distributions at these small values of $`x`$ are required. We propose to use the unifed BFKL/DGLAP evolution equation . We shall start at first with the pure leading order BFKL equation for the unintegrated gluon distribution function $`f(x,k^2)`$ in the followng form:
$$f(x,k^2)=f^{(0)}(x,k^2)+\overline{\alpha }_Sk^2_x^1\frac{dz}{z}\frac{dk^2}{k^2}\left\{\frac{f(x/z,k^2)f(x/z,k^2)}{|k^2k^2|}+\frac{f(x/z,k^2)}{[4k^4+k^4]^{\frac{1}{2}}}\right\}$$
(2)
where $`\overline{\alpha }_S=N_c\alpha _S/\pi `$ and $`k=k_T,k^{}=k_T^{}`$ denote the transverse momenta of the gluons, see Fig. 2. The term in the integrand containing $`f(x/z,k^2)`$ corresponds to real gluon emission, whereas the terms involving $`f(x/z,k^2)`$ represent the virtual contributions and lead to the Reggeization of the $`t`$-channel exchanged gluons. The inhomogeneous driving term $`f^{(0)}`$ is the input function will be specified later. This is the leading order in $`\mathrm{ln}(1/x)`$ equation which gives the very well known intercept for the gluon $`\lambda =\overline{\alpha _s}4\mathrm{ln}2`$. The next to leading contribution to the t-channel exchange at high energies has been already calculated . It yields however very large correction to the intercept making it physically unreliable. It occurred, that the resummation of the subleading effects should be performed in order to get physically reliable results . Different forms of resummation has been already proposed in the literature. It appears that the imposition of the consistency constraint which limits the phase-space available for the real emission term provides with the partial resummation of the subleading effects. This constraint arises from the fact that in the high energy limit the virtualities of the exchanged momenta are dominated by their transverse parts. By imposition of this constraint one obtains nice physical picture in which the subleading effects are resummed by the limitation of the phase space. The result for the gluon intercept yields reasonable value which is stable, i.e. does not become negative. Second improvement to the equation (2) is the inclusion of the DGLAP terms which are important for large values of $`x`$ and the overall normalisation of the resulting gluon distribution function. We do also include the quark driving term in eq. 2. The resulting evolution equation for the gluon has the following form:
$`f(x,k^2)=\stackrel{~}{f}^{(0)}(x,k^2)+`$
$`+\overline{\alpha }_S(k^2)k^2{\displaystyle _x^1}{\displaystyle \frac{dz}{z}}{\displaystyle _{k_0^2}}{\displaystyle \frac{dk^2}{k^2}}\left\{{\displaystyle \frac{f({\displaystyle \frac{x}{z}},k^2)\mathrm{\Theta }\left({\displaystyle \frac{k^2}{z}}k^2\right)f({\displaystyle \frac{x}{z}},k^2)}{|k^2k^2|}}+{\displaystyle \frac{f({\displaystyle \frac{x}{z}},k^2)}{[4k^4+k^4]^{\frac{1}{2}}}}\right\}`$ (3)
$`+\overline{\alpha }_S(k^2){\displaystyle _x^1}{\displaystyle \frac{dz}{z}}\left({\displaystyle \frac{z}{6}}P_{gg}(z)1\right){\displaystyle _{k_0^2}^{k^2}}{\displaystyle \frac{dk^2}{k^2}}f({\displaystyle \frac{x}{z}},k^2)+{\displaystyle \frac{\alpha _S(k^2)}{2\pi }}{\displaystyle _x^1}𝑑zP_{gq}(z)\mathrm{\Sigma }({\displaystyle \frac{x}{z}},k^2).`$
We specify the driving term in the following form
$$\stackrel{~}{f}^{(0)}(x,k^2)=\frac{\alpha _S(k^2)}{2\pi }_x^1𝑑zP_{gg}(z)\frac{x}{z}g(\frac{x}{z},k_0^2)$$
(4)
Let us note that the inhomogenious term has been entirely specifed in terms of the standard integrated gluon distribution function at the infrared cutoff $`k_0^2`$ and that the BFKL/DGLAP evolution only takes place above this cutoff. The last term is the contribution of the quark distribution to the gluon evolution,
$$\mathrm{\Sigma }=\underset{q}{}x(q+\overline{q})=\underset{q}{}(S_q+V_q)$$
(5)
where $`S`$ and $`V`$ denote the sea and valence quark momentum distributions. The gluon, in turn, helps to drive the sea quark distribution through the $`gq\overline{q}`$ transition. Thus equation (3) has to be solved simultaneously with an equivalent equation for $`\mathrm{\Sigma }(x,k^2)`$. We use the $`k_T`$ factorisation formula , see Fig. 2 as a basis for our evolution equation for the quark distribution.
$$S_q(x,Q^2)=_x^1\frac{dz}{z}\frac{dk^2}{k^2}S_{\mathrm{box}}^q(z,k^2,Q^2)f(\frac{x}{z},k^2)$$
(6)
where $`S^{\mathrm{box}}`$ describes the quark box (and crossed-box) contribution shown in Fig. 2 and can be interpreted as a partonic structure function.
It has been shown (see for example , ) that the $`\mathrm{ln}(1/x)`$ effects are also resummed and play important role in the $`k_T`$ factorisation prescription. In order to get the complete set of evolution equations we also have to add quark-quark splittings (which are small numerically anyway) and the valence quarks, which we take from the set of parametrisations. Thus our complete equation for the singlet quark distribution reads as follows,
$`\mathrm{\Sigma }(x,k^2)`$ $`=`$ $`S_{\mathrm{non}\mathrm{p}}(x)+{\displaystyle \underset{q}{}}{\displaystyle _x^a}{\displaystyle \frac{dz}{z}}S_q^{\mathrm{box}}(z,k^2=0,k^2){\displaystyle \frac{x}{z}}g({\displaystyle \frac{x}{z}},k_0^2)`$
$`+{\displaystyle \underset{q}{}}{\displaystyle _{k_0^2}^{\mathrm{}}}{\displaystyle \frac{dk^2}{k^2}}{\displaystyle _x^1}{\displaystyle \frac{dz}{z}}S_q^{\mathrm{box}}(z,k^2,k^2)f({\displaystyle \frac{x}{z}},k^2)+V(x,k^2)`$
$`+{\displaystyle _{k_0^2}^{k^2}}{\displaystyle \frac{dk^2}{k^2}}{\displaystyle \frac{\alpha _S(k^2)}{2\pi }}{\displaystyle _x^1}𝑑zP_{qq}(z)S_{uds}({\displaystyle \frac{x}{z}},k^2)`$
where $`a=(1+4m_q^2/Q^2)^1`$ and $`V=x(u_v+d_v)`$. Equations (3) and (3) form a set of coupled integral equations for the unkown functions $`f(x,k^2)`$ and $`\mathrm{\Sigma }(x,k^2)`$. We solve them assuming simple parametric form of the inputs:
$`xg(x,k_0^2)`$ $`=`$ $`N(1x)^\beta `$
$`S_{\mathrm{non}\mathrm{p}}(x)`$ $`=`$ $`C_p(1x)^8x^{0.08}`$ (8)
We have used this set of coupled equations to calculate $`F_2`$ structure function at HERA, and we have therefore fixed the values of the free parameters. We then used the resulting parton distribution functions in order to calculate the neutrino - nucleon cross section at high energies. We calculate the cross sections from the usual formula:
$`{\displaystyle \frac{d^2\sigma ^{\nu ,\overline{\nu }}}{dxdy}}`$ $`=`$ $`{\displaystyle \frac{G_FME}{\pi }}\left({\displaystyle \frac{M_i^2}{Q^2+M_i^2}}\right)^2\{{\displaystyle \frac{1+(1y)^2}{2}}F_2^\nu (x,Q^2)`$
$`{\displaystyle \frac{y^2}{2}}F_L^\nu (x,Q^2)\underset{¯}{+}y(1{\displaystyle \frac{y}{2}})xF_3^\nu (x,Q^2)\}`$
where $`G_F`$ is the Fermi coupling constant, $`M`$ is the proton mass, $`E`$ is the laboratory energy of the neutrino and $`y=Q^2/xs`$. The mass $`M_i`$ is either $`M_W`$ or $`M_Z`$ according to whether we are calculating charged current (CC) or neutral current (NC) neutrino interactions.
Functions $`F_2^\nu `$, $`F_L^\nu `$ and $`xF_3^\nu `$ are of course usual structure functions which can be calculated from the parton distribution functions.
In Fig. 3 we show the plot of $`\sigma (\nu N)`$ as a function of energy of the neutrino $`E_\nu `$. One observes strong rise, nearly $`8`$ decades with increasing energy from $`10`$ to $`10^{12}\mathrm{GeV}`$. The shape of the curves up to $`10^5\mathrm{GeV}`$ is determined by the valence quarks whereas beyond $`10^6\mathrm{GeV}`$ everything is driven by the sea quarks. One also notes that the charged current contribution is dominating over the neutral current by factor 3.
## 4 Transport equation
Neutrinos at ultrahigh energies can be quite strongly attenuated when traversing the Earth. Apart from standard absorption neutrinos can undergo regeneration due to neutral current interactions. Charged current interactions remove neutrinos from the flux, but neutral current interaction cause neutrinos to reappear at lower energies. These both effects can be calculated using the transport equation proposed by :
$$\frac{dI(E,\tau )}{d\tau }=\sigma _{\mathrm{TOT}}(E)I(E,\tau )+\frac{dy}{1y}\frac{d\sigma _{\mathrm{NC}}(E^{},y)}{dy}I(E^{},\tau )$$
(10)
where $`\sigma _{\mathrm{TOT}}=\sigma _{\mathrm{CC}}+\sigma _{\mathrm{NC}}`$ and where $`y`$ is, as usual, the fractional energy loss such that
$$E^{}=\frac{E}{1y}.$$
(11)
The variable $`\tau `$ is the number density of nucleons $`n`$ integrated along a path of length $`z`$ through the Earth
$$\tau (z)=_0^z𝑑z^{}n(z^{}).$$
(12)
The number density $`n(z)`$ is defined as $`n(z)=N_A\rho (z)`$ where $`\rho (z)`$ is the density of Earth along the neutrino path length $`z`$ and $`N_A`$ is the Avogadro number. The number of nucleons $`\tau `$ encountered along the path $`z`$ depends upon the nadir angle $`\theta `$ between the normal to the Earth’s surface (passing through the detector) and the direction of the neutrino beam incident on the detector.
In order to calculate the change of the intensity of the flux with the incident angle one needs to know the density profile of the Earth. We have used the model by A. Dziewoński . Using this parametrisation of the Earth density and the cross sections calculated from the unifed BFKL/DGLAP equation. In Fig. 4 we show the shadowing factor,
$$S(E,\tau )=\frac{I(E,\tau )}{I^0(E)}$$
(13)
where $`I^0(E)=I(E,\tau )`$ is the initial flux at the surface of the Earth. We present $`S`$ as a function of the energy and for different incident angles. The curves exhibit strong suppression for large paths in matter and the energies above $`10^6\mathrm{GeV}`$. Also, the curves corresponding to the flux from AGN , show that the regeneration is important for flat fluxes and large paths. The same effect is nearly negligible in the case of steeply falling atmospheric spectrum .
In Figs. 5 and 6 we show complete simulation for different incident fluxes corresponding to the active galactic nuclei , gamma ray bursts and a sample top-down model . Similar effects of strong attenuation are observed for large energies $`10^6\mathrm{GeV}^2`$ and large paths. This suggests that one will have to choose suitable angle of the observation in order to avoid large muon background from the atmosphere and yet be able to detect the highly energetic neutrinos which are likely to be absorbed by the matter in Earth.
Summary
In this paper we have examined the interactions of ultrahigh energy neutrinos with matter. As the cross section rises with increasing energy the absorption by matter becomes interestingly large. We have seen that the Earth becomes essentially opaque to ultrahigh energy neutrinos. We have calculated neutrino-nucleon cross section using the unifed DGLAP/BFKL evolution equations which treat leading $`\mathrm{ln}(Q^2)`$ and $`\mathrm{ln}(1/x)`$ on equal footing. These equations also resum important subleading effects in $`\mathrm{ln}(1/x)`$ via imposition of consistency constraint. We believe that this form of the evolution is most apropriate in the regime where the parton density is large. The results for the cross section has been compared with the other based on the standard global parton analysis. In that way the uncertainty due to parton density extrapolation has been diminished to $`40\%`$. We have then used the resulting cross sections and solved the transport equation for the neutrinos travelling through the Earth. We have found that the attenuation is large for high energies , above $`10^6\mathrm{GeV}`$ and large paths in matter. We have also found that the regeneration due to neutral current interactions becomes important for flat spectra (like these originating from active galactic nuclei) and large paths. We have simulated the penetration through the Earth for different neutrino fluxes: active galactic nuclei, gamma ray bursts and top-down models. The large attenuation effect reduces substantially the flux for small nadir angles. There is however a window for observation of AGN fluxes by the $`\mathrm{km}^3`$ detectors.
Acknowledgements
This research has been supported in part by the Polish State Committee for Scientific Research (KBN) grant N0 2 P03B 89 13 and by the EU Fourth Framework Programme “Training and Mobility of Researchers”, Network “Quantum Chromodynamics and the Deep Structure of Elementary Particles”, contract FMRX - CT98 - 0194.
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# 1 Introduction
## 1 Introduction
We have recently shown that the modified quantization scheme
$$𝒟A_\mu \{\mathrm{exp}[S_{YM}(A_\mu )f(A_\mu )𝑑x]/𝒟g\mathrm{exp}[f(A_\mu ^g)𝑑x]\}$$
(1.1)
with a gauge non-invariant term $`f(A_\mu )`$, for example,
$$f(A_\mu )=(m^2/2)(A_\mu )^2$$
(1.2)
is identical at least to the conventional local Faddeev-Popov formula
$`{\displaystyle 𝒟A_\mu \{\delta (D^\mu \frac{\delta f(A_\nu )}{\delta A_\mu })/𝒟g\delta (D^\mu \frac{\delta f(A_\nu ^g)}{\delta A_\mu ^g})\}\mathrm{exp}[S_{YM}(A_\mu )]}`$
$`={\displaystyle 𝒟A_\mu \delta (D^\mu \frac{\delta f(A_\nu )}{\delta A_\mu })det\{\delta [D^\mu \frac{\delta f(A_\nu ^g)}{\delta A_\mu ^g}]/\delta g\}\mathrm{exp}[S_{YM}(A_\mu )]}`$ (1.3)
without taking the large mass limit, if one takes into account the variation of the gauge field along the entire gauge orbit parametrized by the gauge parameter $`g`$. Here the operator $`D^\mu \frac{\delta f(A_\nu )}{\delta A_\mu }(x)`$ is defined by an infinitesimal gauge transformation as
$$𝑑xf(A_\nu +D_\nu \omega )=𝑑xf(A_\nu )𝑑x\omega (x)D^\mu \frac{\delta f(A_\nu )}{\delta A_\mu }(x).$$
(1.4)
The above equivalence was discussed in in connection with the analysis of the so-called Gribov problem, and the above formula is valid if the Gribov-type complications are ignored such as in perturbative approach. We, however, note that the choice of the gauge fixing function $`f(A_\mu )`$ is rather general but not completely arbitrary since certain technical conditions need to be satisfied in the above proof of the equivalence.
One can confirm that the starting expression (1.1) is formally identical to the extraction of the gauge volume from the naive path integral measure,
$$\frac{𝒟A_\mu }{vol(g)}\mathrm{exp}[S_{YM}(A_\mu )]$$
(1.5)
and thus the above equivalence is not unexpected. A remarkable aspect is that the apparently non-local expression (1.1) in fact defines a local and thus unitary theory (1.3).
We emphasize that the above equivalence is valid for any value of the mass parameter $`m^2`$, for example. The large mass limit, where the equivalence to the conventional local formula was analyzed in the past, is formally related to the non-linear gauge
$$A_\mu ^2=\lambda =const.,$$
(1.6)
of Dirac and Nambu in the limit $`\lambda =0`$. Nambu used the above gauge to analyze the possible spontaneous breakdown of Lorentz symmetry. In his treatment, the limit $`\lambda =0`$ is singular, and thus the present formulation is not quite convenient for an analysis of the possible breakdown of Lorentz symmetry. See, for example, for the past analyses of the above non-linear gauge.
In this paper, we comment on the possible implications of the above equivalence between (1.1) and the local BRST invariant theory (1.3) in a general context of quantum gauge symmetry, namely, BRST symmetry, which controls the analyses of renormalization and unitarity. The modified quantization scheme is flexible in choosing gauge fixing functions $`f(A_\mu )`$, and we argue that from a view point of quantum gauge symmetry there is no intrinsic difference between the classical theory with some extra terms such as a mass term, which break gauge symmetry, and the classical gauge theory whose gauge symmetry is broken by a gauge fixing term. In particular, the classical massive Yang-Mills theory is more flexibly interpreted either as a gauge fixed version of pure Yang-Mills theory in the modified quantization scheme, or as the conventional massive non-gauge theory.
Our basic argument is a possible re-interpretation of various classical Lagrangians, but it could also be understood as a broader choice of path integral measure than in the conventional analysis starting with the Dirac bracket. The spirit of our approach is probably illustrated by considering the two-dimensional field theory coupled to gravity
$$=\frac{1}{2}\sqrt{g}g^{\mu \nu }_\mu X^a(x)_\nu X^a(x).$$
(1.7)
In the conformal gauge defined by
$$g_{\mu \nu }=\delta _{\mu \nu }\rho (x)$$
(1.8)
the metric $`\rho (x)`$ decouples from the above action, and one may define the path integral
$$𝒟X^a\mathrm{exp}[\frac{1}{2}_\mu X^a(x)^\mu X^a(x)dx].$$
(1.9)
On the other hand, it is well known that a reparametrization invariant path integral measure leads to
$$𝒟X^a𝒟\sqrt{\rho }\mathrm{exp}\{\frac{1}{2}_\mu X^a(x)^\mu X^a(x)dx\frac{26d}{24\pi }[_\mu \mathrm{ln}\rho ^\mu \mathrm{ln}\rho +\mu ^2\rho ]𝑑x\}$$
(1.10)
where $`d`$ stands for the number of variables $`X^a`$. The last term in the action, which is regarded as a Wess-Zumino term, is known as the Liouville action. In both of these expressions, we neglected the Faddeev-Popov ghosts, which should be included in a consistent quantization . The above example, which appears in the first quantization of string theory, illustrates that the choice of path integral measure has more freedom than in the naive prescription starting with Dirac brackets. It is also known that many interesting physical examples can be quantized consistently only when one adds a suitable Wess-Zumino term, which may be regarded as a modified choice of measure. See, for example, the quantization of anomalous gauge theory in 2-dimensions and the quantization of supermembrane in.
## 2 Abelian example
We first briefly illustrate the proof of the above equivalence of (1.1) and (1.3) by using an example of Abelian gauge theory,
$$S_0=\frac{1}{4}𝑑x(_\mu A_\nu _\nu A_\mu )^2$$
(2.1)
for which we can work out everything explicitly. In this note we exclusively work on Euclidean theory with metric convention $`g_{\mu \nu }=(1,1,1,1)`$. Note that there is no Gribov complications in the Abelian theory at least in a continuum formulation. As a simple and useful example, we choose the gauge fixing function
$$f(A)\frac{1}{2}A_\mu A_\mu $$
(2.2)
and
$$D_\mu (\frac{\delta f}{\delta A_\mu })=_\mu A_\mu .$$
(2.3)
Our claim above suggests the relation
$`Z`$ $`=`$ $`{\displaystyle 𝒟A_\mu ^\omega \{e^{S_0(A_\mu ^\omega ){\scriptscriptstyle 𝑑x{\scriptscriptstyle \frac{1}{2}}(A_\mu ^\omega )^2}}/𝒟he^{{\scriptscriptstyle 𝑑x{\scriptscriptstyle \frac{1}{2}}(A_\mu ^{h\omega })^2}}\}}`$ (2.4)
$`=`$ $`{\displaystyle 𝒟A_\mu ^\omega 𝒟B𝒟\overline{c}𝒟ce^{S_0(A_\mu ^\omega )+{\scriptscriptstyle [iB_\mu A_\mu ^\omega +\overline{c}(_\mu _\mu )c]𝑑x}}}`$
where the variable $`A_\mu ^\omega `$ stands for the field variable obtained from $`A_\mu `$ by a gauge transformation parametrized by the gauge orbit parameter $`\omega `$. To establish this result, we first evaluate
$`{\displaystyle 𝒟he^{{\scriptscriptstyle 𝑑x{\scriptscriptstyle \frac{1}{2}}(A_\mu ^{h\omega })^2}}}`$
$`={\displaystyle 𝒟he^{{\scriptscriptstyle 𝑑x{\scriptscriptstyle \frac{1}{2}}(A_\mu ^\omega +_\mu h)^2}}}`$
$`={\displaystyle 𝒟he^{{\scriptscriptstyle 𝑑x{\scriptscriptstyle \frac{1}{2}}[(A_\mu ^\omega )^22(_\mu A_\mu ^\omega )h+h(_\mu _\mu )h]}}}`$
$`={\displaystyle 𝒟B\frac{1}{det\sqrt{_\mu _\mu }}e^{{\scriptscriptstyle 𝑑x{\scriptscriptstyle \frac{1}{2}}[(A_\mu ^\omega )^22(_\mu A_\mu ^\omega ){\scriptscriptstyle \frac{1}{\sqrt{_\mu _\mu }}}B+B^2]}}}`$
$`={\displaystyle \frac{1}{det\sqrt{_\mu _\mu }}}e^{{\scriptscriptstyle 𝑑x{\scriptscriptstyle \frac{1}{2}}(A_\mu ^\omega )^2}+\frac{1}{2}{\scriptscriptstyle _\mu A_\mu ^\omega {\scriptscriptstyle \frac{1}{_\mu _\mu }}_\nu A_\nu ^\omega dx}}`$ (2.5)
where we defined $`\sqrt{_\mu _\mu }h=B`$. Thus
$`Z`$ $`=`$ $`{\displaystyle 𝒟A_\mu ^\omega \{det\sqrt{_\mu _\mu }\}e^{S_0(A_\mu ^\omega )\frac{1}{2}{\scriptscriptstyle _\mu A_\mu ^\omega {\scriptscriptstyle \frac{1}{_\mu _\mu }}_\nu A_\nu ^\omega dx}}}`$ (2.6)
$`=`$ $`{\displaystyle 𝒟A_\mu ^\omega 𝒟B𝒟\overline{c}𝒟ce^{S_0(A_\mu ^\omega )\frac{1}{2}{\scriptscriptstyle B^2𝑑x}+{\scriptscriptstyle [iB{\scriptscriptstyle \frac{1}{\sqrt{_\mu _\mu }}}_\mu A_\mu ^\omega +\overline{c}\sqrt{_\mu _\mu }c]𝑑x}}}`$
which is invariant under the BRST transformation
$`\delta A_\mu ^\omega =i\lambda _\mu c`$
$`\delta c=0`$
$`\delta \overline{c}=\lambda B`$
$`\delta B=0`$ (2.7)
with a Grassmann parameter $`\lambda `$. Note the appearance of the imaginary factor $`i`$ in the term $`iB\frac{1}{\sqrt{_\mu _\mu }}_\mu A_\mu ^\omega `$ in (2.6). The expression (2.6) has been derived in . To show that we in fact obtain a local theory, we need to go one more step further.
We next rewrite the expression (2.6) as
$`{\displaystyle 𝒟A_\mu ^\omega 𝒟B𝒟\mathrm{\Lambda }𝒟\overline{c}𝒟c\delta (\frac{1}{\sqrt{_\mu _\mu }}_\mu A_\mu ^\omega \mathrm{\Lambda })e^{S_0(A_\mu ^\omega )\frac{1}{2}{\scriptscriptstyle (B^2+2i\mathrm{\Lambda }B)𝑑x}+{\scriptscriptstyle \overline{c}\sqrt{_\mu _\mu }c𝑑x}}}`$
$`={\displaystyle 𝒟A_\mu ^\omega 𝒟\mathrm{\Lambda }𝒟\overline{c}𝒟c\delta (\frac{1}{\sqrt{_\mu _\mu }}_\mu A_\mu ^\omega \mathrm{\Lambda })e^{S_0(A_\mu ^\omega )\frac{1}{2}{\scriptscriptstyle \mathrm{\Lambda }^2𝑑x}+{\scriptscriptstyle \overline{c}\sqrt{_\mu _\mu }c𝑑x}}}.`$ (2.8)
We note that we can compensate any variation of $`\delta \mathrm{\Lambda }`$ by a suitable change of gauge parameter $`\delta \omega `$ inside the $`\delta `$-function as
$$\frac{1}{\sqrt{_\mu _\mu }}_\mu _\mu \delta \omega =\delta \mathrm{\Lambda }.$$
(2.9)
By a repeated application of infinitesimal gauge transformations combined with the invariance of the path integral measure under these gauge transformations, we can re-write the formula (2.8) as
$`{\displaystyle 𝒟A_\mu ^\omega 𝒟\mathrm{\Lambda }𝒟\overline{c}𝒟c\delta (\frac{1}{\sqrt{_\mu _\mu }}_\mu A_\mu ^\omega )e^{S_0(A_\mu ^\omega )\frac{1}{2}{\scriptscriptstyle \mathrm{\Lambda }^2𝑑x}+{\scriptscriptstyle \overline{c}\sqrt{_\mu _\mu }c𝑑x}}}`$
$`={\displaystyle 𝒟A_\mu ^\omega 𝒟\overline{c}𝒟c\delta (\frac{1}{\sqrt{_\mu _\mu }}_\mu A_\mu ^\omega )e^{S_0(A_\mu ^\omega )+{\scriptscriptstyle \overline{c}\sqrt{_\mu _\mu }c𝑑x}}}`$
$`={\displaystyle 𝒟A_\mu ^\omega 𝒟B𝒟\overline{c}𝒟ce^{S_0(A_\mu ^\omega )+{\scriptscriptstyle [iB{\scriptscriptstyle \frac{1}{\sqrt{_\mu _\mu }}}_\mu A_\mu ^\omega +\overline{c}\sqrt{_\mu _\mu }c]𝑑x}}}`$
$`={\displaystyle 𝒟A_\mu ^\omega 𝒟B𝒟\overline{c}𝒟ce^{S_0(A_\mu ^\omega )+{\scriptscriptstyle [iB_\mu A_\mu ^\omega +\overline{c}(_\mu _\mu )c]𝑑x}}}.`$ (2.10)
In the last stage of this equation, we re-defined the auxiliary variables $`B`$ and $`\overline{c}`$ as
$`BB\sqrt{_\mu _\mu }`$
$`\overline{c}\overline{c}\sqrt{_\mu _\mu }`$ (2.11)
which is consistent with BRST symmetry and leaves the path integral measure invariant. We have thus established the desired result (2.4). We emphasize that the integral over the entire gauge orbit, as is illustrated in (2.9), is crucial to obtain a local expression (2.10).
It is shown that this procedure works for the non-Abelian case also, though the actual procedure is much more involved and implicit, if the (ill-understood) Gribov-type complications can be ignored such as in perturbative calculations.
## 3 No massive gauge fields?
In the classical level, we traditionally consider
$$=\frac{1}{4}(_\mu A_\nu _\nu A_\mu )^2\frac{1}{2}m^2A_\mu A^\mu $$
(3.1)
as a Lagrangian for a massive vector theory, and
$$=\frac{1}{4}(_\mu A_\nu _\nu A_\mu )^2\frac{1}{2}(_\mu A^\mu )^2$$
(3.2)
as an effective Lagrangian for Maxwell theory with a Feynman-type gauge fixing term added. The physical meanings of these two Lagrangians are thus completely different.
However, the analysis in Section 2 shows that the Lagrangian (3.1) could in fact be regarded as a gauge fixed Lagrangian of massless Maxwell field in quantized theory. To be explicit, by using (2.4), the Lagrangian (3.1) may be regarded as an effective Lagrangian in
$`Z`$ $`=`$ $`{\displaystyle 𝒟A_\mu ^\omega \{e^{{\scriptscriptstyle 𝑑x[{\scriptscriptstyle \frac{1}{4}}(_\mu A_\nu _\nu A_\mu )^2{\scriptscriptstyle \frac{1}{2}}m^2A_\mu ^\omega A^{\omega \mu }]}}/𝒟he^{{\scriptscriptstyle 𝑑x{\scriptscriptstyle \frac{m^2}{2}}(A_\mu ^{h\omega })^2}}\}}`$ (3.3)
$`=`$ $`{\displaystyle 𝒟A_\mu ^\omega 𝒟B𝒟\overline{c}𝒟ce^{{\scriptscriptstyle 𝑑x[{\scriptscriptstyle \frac{1}{4}}(_\mu A_\nu _\nu A_\mu )^2iB_\mu A_\mu ^\omega +\overline{c}(_\mu _\mu )c]𝑑x}}}`$
where we absorbed the factor $`m^2`$ into the definition of $`B`$ and $`\overline{c}`$.
One can also analyze (3.2) by defining
$$f(A_\mu )\frac{1}{2}(_\mu A^\mu )^2$$
(3.4)
in the modified quantization scheme (1.1). The equality of (1.1) and (1.3) then gives
$`{\displaystyle 𝒟A_\mu \delta (D^\mu \frac{\delta f(A_\nu )}{\delta A_\mu })det\{\delta [D^\mu \frac{\delta f(A_\nu ^g)}{\delta A_\mu ^g}]/\delta g\}\mathrm{exp}[S_0(A_\mu )]}`$
$`={\displaystyle 𝒟A_\mu \delta (_\nu ^\nu (^\mu A_\mu ))det[_\nu ^\nu _\mu ^\mu ]\mathrm{exp}[S_0(A_\mu )]}`$ (3.5)
$`={\displaystyle 𝒟A_\mu 𝒟B𝒟\overline{c}𝒟c\mathrm{exp}\{S_0(A_\mu )+𝑑x[iB_\nu ^\nu (^\mu A_\mu )\overline{c}(_\nu ^\nu _\mu ^\mu )c]\}}`$
After the re-definition of auxiliary variables,
$$B_\nu ^\nu B,\overline{c}_\nu ^\nu \overline{c}$$
(3.6)
which preserves BRST symmetry, (3.5) becomes
$$𝒟A_\mu 𝒟B𝒟\overline{c}𝒟c\mathrm{exp}\{S_0(A_\mu )+𝑑x[iB(^\mu A_\mu )+\overline{c}(_\mu ^\mu )c]\}$$
(3.7)
which agrees with (2.10) and (3.3).
We have thus shown that an identical physical meaning can be assigned to two Lagrangians (3.1) and (3.2) in suitably quantized theory.
In passing, a conventional way to relate (3.2) and (3.3) is to note
$`{\displaystyle 𝒟A_\mu ^\omega 𝒟B𝒟\overline{c}𝒟ce^{{\scriptscriptstyle 𝑑x[{\scriptscriptstyle \frac{1}{4}}(_\mu A_\nu _\nu A_\mu )^2iB_\mu A_\mu ^\omega +\overline{c}(_\mu _\mu )c]𝑑x}}}`$
$`={\displaystyle 𝒟A_\mu ^\omega 𝒟B𝒟\overline{c}𝒟ce^{{\scriptscriptstyle 𝑑x[{\scriptscriptstyle \frac{1}{4}}(_\mu A_\nu _\nu A_\mu )^2{\scriptscriptstyle \frac{1}{2}}\xi B^2iB_\mu A_\mu ^\omega +\overline{c}(_\mu _\mu )c]𝑑x}}}`$
$`={\displaystyle 𝒟A_\mu ^\omega 𝒟\overline{c}𝒟ce^{{\scriptscriptstyle 𝑑x[{\scriptscriptstyle \frac{1}{4}}(_\mu A_\nu _\nu A_\mu )^2{\scriptscriptstyle \frac{1}{2\xi }}(_\mu A_\mu ^\omega )^2+\overline{c}(_\mu _\mu )c]𝑑x}}}`$ (3.8)
and by setting the gauge parameter $`\xi =1`$. The first equality of (3.8), namely, the $`\xi `$ independence of the partition function is established by the BRST identity as follows: When one defines
$$Z(\xi )=𝒟A_\mu ^\omega 𝒟B𝒟\overline{c}𝒟ce^{S_0(A^\omega )\frac{\xi }{2}{\scriptscriptstyle B^2𝑑x}+{\scriptscriptstyle [iB_\mu A_\mu ^\omega +\overline{c}(_\mu _\mu )c]𝑑x}}$$
(3.9)
with $`S_0(A^\omega )`$ defined in (2.1), one can show that
$`Z(\xi \delta \xi )`$
$`={\displaystyle 𝒟A_\mu ^\omega 𝒟B𝒟\overline{c}𝒟ce^{S_0(A_\mu ^\omega )\frac{\xi \delta \xi }{2}{\scriptscriptstyle B^2𝑑x}+{\scriptscriptstyle [iB_\mu A_\mu ^\omega +\overline{c}(_\mu _\mu )c]𝑑x}}}`$
$`=Z(\xi )`$
$`+`$ $`{\displaystyle 𝒟A_\mu ^\omega 𝒟B𝒟\overline{c}𝒟c(\frac{\delta \xi }{2}B^2𝑑x)e^{S_0(A_\mu ^\omega )\frac{\xi }{2}{\scriptscriptstyle B^2𝑑x}+{\scriptscriptstyle [iB_\mu A_\mu ^\omega +\overline{c}(_\mu _\mu )c]𝑑x}}}.`$
On the other hand,the BRST invariance of the path integral measure and the effective action in the exponential factor gives rise to
$`{\displaystyle 𝒟A_\mu ^\omega 𝒟B𝒟\overline{c}𝒟c[(\overline{c}B)𝑑x]e^{S_0(A_\mu ^\omega )\frac{\xi }{2}{\scriptscriptstyle B^2𝑑x}+{\scriptscriptstyle [iB_\mu A_\mu ^\omega +\overline{c}(_\mu _\mu )c]𝑑x}}}`$
$`={\displaystyle 𝒟A_\mu ^\omega 𝒟B^{}𝒟\overline{c}^{}𝒟c^{}[(\overline{c}^{}B^{})𝑑x]}`$
$`\times e^{S_0(A_\mu ^\omega )\frac{\xi }{2}{\scriptscriptstyle B_{}^{}{}_{}{}^{2}𝑑x}+{\scriptscriptstyle [iB^{}_\mu A_\mu ^\omega +\overline{c}^{}(_\mu _\mu )c^{}]𝑑x}}`$
$`={\displaystyle 𝒟A_\mu ^\omega 𝒟B𝒟\overline{c}𝒟c[(\overline{c}B+\delta _{BRST}(\overline{c}B))𝑑x]}`$
$`\times e^{S_0(A_\mu ^\omega )\frac{\xi }{2}{\scriptscriptstyle B^2𝑑x}+{\scriptscriptstyle [iB_\mu A_\mu ^\omega +\overline{c}(_\mu _\mu )c]𝑑x}}`$ (3.11)
where the BRST transformed variables are defined by $`A_\mu ^\omega =A_\mu ^\omega +i\lambda _\mu c,B^{}=B,\overline{c}^{}=\overline{c}+\lambda B,c^{}=c`$. The first equality of the above relation means that the path integral itself is independent of the naming of integration variables, and the second equality follows from the BRST invariance of the path integral measure and the action. Namely, the BRST exact quantity vanishes as
$`{\displaystyle 𝒟A_\mu ^\omega 𝒟B𝒟\overline{c}𝒟c(\delta _{BRST}(\overline{c}B)𝑑x)e^{S_0(A_\mu ^\omega )\frac{\xi }{2}{\scriptscriptstyle B^2𝑑x}+{\scriptscriptstyle [iB_\mu A_\mu ^\omega +\overline{c}(_\mu _\mu )c]𝑑x}}}`$
$`={\displaystyle 𝒟A_\mu ^\omega 𝒟B𝒟\overline{c}𝒟c(\lambda B^2𝑑x)e^{S_0(A^\omega )\frac{\xi }{2}{\scriptscriptstyle B^2𝑑x}+{\scriptscriptstyle [iB_\mu A_\mu ^\omega +\overline{c}(_\mu _\mu )c]𝑑x}}}`$
$`=0.`$ (3.12)
Thus from (3.10) we have the relation
$$Z(\xi \delta \xi )=Z(\xi ).$$
(3.13)
Namely, $`Z(\xi )`$ is independent of $`\xi `$, and $`Z(1)=Z(0)`$. This justifies the equality used in (3.8), and we have established the conventional interpretation of (3.2) as an effective Lagrangian with a gauge fixing term added.
Similarly, the two classical Lagrangians related to Yang-Mills fields
$$=\frac{1}{4}(_\mu A_\nu ^a_\nu A_\mu ^a+gf^{abc}A_\mu ^bA_\nu ^c)^2\frac{m^2}{2}A_\mu ^aA^{a\mu }$$
(3.14)
and
$$=\frac{1}{4}(_\mu A_\nu ^a_\nu A_\mu ^a+gf^{abc}A_\mu ^bA_\nu ^c)^2\frac{1}{2}(_\mu A^{a\mu })^2$$
(3.15)
could be assigned an identical physical meaning as an effective gauge fixed Lagrangian associated with the quantum theory defined by
$$𝒟A_\mu ^a𝒟B^a𝒟\overline{c}^a𝒟c^a\mathrm{exp}\{S_{YM}(A_\mu ^a)+dx[iB^a(^\mu A_\mu ^a)+\overline{c}^a(_\mu (D^\mu c)^a]\}$$
(3.16)
which is invariant under the quantum gauge symmetry ( BRST transformation) with a Grassmann parameter $`\lambda `$
$`\delta A_\mu ^a=i\lambda (D_\mu c)^a`$
$`\delta c^a={\displaystyle \frac{i\lambda }{2}}f^{abc}c^bc^c`$
$`\delta \overline{c}^a=\lambda B^a`$
$`\delta B^a=0.`$ (3.17)
In this analysis, we ignore the (ill-understood) Gribov-type complications. This connection with the possible Gribov-type complications becomes transparent if one considers
$$f(A^{a\mu })=\frac{1}{2}(_\mu A^{a\mu })^2$$
(3.18)
in (3.15). We then obtain in the modified scheme (1.1) and its equivalent formula (1.3) (by suppressing the Yang-Mills indices)
$$D^\mu \frac{\delta f(A_\nu )}{\delta A_\mu }=D_\mu ^\mu (_\nu A^\nu )$$
(3.19)
and the associated determinant factor contains the operator
$$D_\mu ^\mu (_\nu D^\nu )+\frac{\delta D_\mu (A_\mu ^\omega )}{\delta \omega }|_{\omega =0}^\mu (_\nu A^\nu ).$$
(3.20)
The gauge fixing and Faddeev-Popov terms then become in the modified scheme
$`_g`$ $`=`$ $`iBD_\mu ^\mu (_\nu A^\nu )\overline{c}[D_\mu ^\mu (_\nu D^\nu )+{\displaystyle \frac{\delta D_\mu (A_\mu ^\omega )}{\delta \omega }}|_{\omega =0}^\mu (_\nu A^\nu )]c`$ (3.21)
$`=`$ $`iBD_\mu ^\mu (_\nu A^\nu )\overline{c}D_\mu ^\mu (_\nu D^\nu )c.`$
In this last step, we used the fact that the gauge fixing condition
$$D_\mu ^\mu (_\nu A^\nu )=0$$
(3.22)
is equivalent to
$$_\nu A^\nu =0$$
(3.23)
in Euclidean theory if the Gribov-type complications are absent and thus the inverse of the operator $`D_\mu ^\mu `$ is well-defined. We thus have the path integral formula in the modified scheme
$`{\displaystyle 𝒟A_\mu ^a𝒟B^a𝒟\overline{c}^a𝒟c^a}`$
$`\times \mathrm{exp}\{S_{YM}(A_\mu ^a)+{\displaystyle }dx[iBD_\mu ^\mu (^\nu A_\nu )\overline{c}D_\mu ^\mu (_\nu D^\nu )c\}`$
$`={\displaystyle 𝒟A_\mu ^a𝒟B^a𝒟\overline{c}^a𝒟c^a}`$
$`\times \mathrm{exp}\{S_{YM}(A_\mu ^a)+{\displaystyle }dx[iB(^\mu A_\mu )+\overline{c}(_\mu (D^\mu c)]\}`$ (3.24)
after the re-definition of auxiliary variables
$$BD_\mu ^\mu B,\overline{c}D_\mu ^\mu \overline{c}$$
(3.25)
which leaves the path integral measure invariant. This last re-definition is allowed only when the operator $`D_\mu ^\mu `$ is well-defined, namely, in the absence of Gribov-type complications in Euclidean theory.
We have illustrated that the apparent “massive gauge field” in the classical level (3.14) has no intrinsic physical meaning. It can be interpreted either as a massive (non-gauge) vector theory, which is a conventional interpretation, or as a gauge-fixed effective Lagrangian for a massless gauge field, which is allowed in the modified quantization scheme. In the framework of path integral quantization, these different interpretations of a given classical Lagrangian could also be understood that we have a certain freedom in the choice of the path integral measure (although the conventional interpretation of the measure is perfectly possible as is shown in (1.5)): One choice of the measure
$`{\displaystyle 𝑑\mu \mathrm{exp}\{𝑑x[\frac{1}{4}(_\mu A_\nu ^a_\nu A_\mu ^a+gf^{abc}A_\mu ^bA_\nu ^c)^2\frac{m^2}{2}A_\mu ^aA^{a\mu }]\}}`$
$`{\displaystyle 𝒟A_\mu \frac{1}{𝒟g\mathrm{exp}[\frac{m^2}{2}(A_\mu ^{ag})^2𝑑x]}}`$ (3.26)
$`\times \mathrm{exp}\{{\displaystyle 𝑑x[\frac{1}{4}(_\mu A_\nu ^a_\nu A_\mu ^a+gf^{abc}A_\mu ^bA_\nu ^c)^2\frac{m^2}{2}A_\mu ^aA^{a\mu }]}\}`$
gives rise to a renormalizable gauge theory (3.16), and the other choice,which is suggested by the naive classical analysis,
$`{\displaystyle 𝑑\mu \mathrm{exp}\{𝑑x[\frac{1}{4}(_\mu A_\nu ^a_\nu A_\mu ^a+gf^{abc}A_\mu ^bA_\nu ^c)^2\frac{m^2}{2}A_\mu ^aA^{a\mu }]\}}`$ (3.27)
$`{\displaystyle 𝒟A_\mu \mathrm{exp}\{𝑑x[\frac{1}{4}(_\mu A_\nu ^a_\nu A_\mu ^a+gf^{abc}A_\mu ^bA_\nu ^c)^2\frac{m^2}{2}A_\mu ^aA^{a\mu }]\}}`$
gives rise to a massive non-gauge vector theory, which is not renormalizable. As we emphasized in Section 1, we here allow a more general definition of path integral measure than the one suggested by the naive classical analysis on the basis of Dirac brackets, if one should take (3.14) as a fundamental Lagrangian. It is interesting that the choice of the measure in (3.26) defines a renormalizable unitary theory starting with the apparently inconsistent theory (3.14). It is an advantage of the modified quantization scheme that we can readily assign a physical meaning to a wider class of theories. A somewhat analogous situation to (3.26) arises when one attempts to quantize the so-called anomalous gauge theory: A suitable choice of the measure with a Wess-Zumino term gives rise to a consistent quantum theory in 2-dimensions, for example .
From a view point of classical-quantum correspondence, one can define a classical theory uniquely starting from quantum theory by considering the limit $`\mathrm{}0`$, but not the other way around in general. For example, the ambiguities related to the operator ordering are well known in any quantization, though the present choices of path integral measure are not directly related to operator ordering ambiguities.
In the context of the present broader interpretation of classical massive gauge fields, the massive gauge fields generated by the Higgs mechanism are exceptional and quite different. The Higgs mechanism for Abelian theory, for example, is defined by (in this part, we use the Minkowski metric with $`g_{\mu \nu }=(1,1,1,1)`$)
$$=(D^\mu \varphi )^{}D_\mu \varphi \mu ^2|\varphi |^2\lambda |\varphi |^4\frac{1}{4}(_\mu A_\nu _\nu A_\mu )^2$$
(3.28)
which is manifestly gauge invariant with $`D_\mu =_\mu igA_\mu `$. The mass $`m=gv`$ for the gauge field is generated after the spontaneous symmetry breaking of gauge symmetry defined by
$$\varphi (x)\varphi ^{}(x)+v/\sqrt{2}$$
(3.29)
with
$$v^2=\mu ^2/\lambda $$
(3.30)
for $`\mu ^2<0`$. In this procedure, all the terms in the Lagrangian including the mass term generated by the Higgs mechanism are gauge invariant. Consequently, our argument discussed so far (i.e,. a possible re-interpretation of the mass term as a gauge fixing term ) does not apply to the present massive vector particle whose mass is generated by the Higgs mechanism. It is quite satisfactory that the Higgs mechanism has an intrinsic physical meaning even in our extended interpretation of mass terms.
In view of the well known fact that the massive non-Abelian gauge theory is inconsistent as a quantum theory (3.27), it may be sensible to treat all the classical massive non-Abelian Lagrangians as a gauge fixed version of pure non-Abelian gauge theory and to restrict the massive non-Abelian gauge fields to those generated by the Higgs mechanism. In this connection, we note that our discussion is quite different from the Stueckelberg formalism of the classical massive vector theory, which was also an attempt to make sense out of a massive vector theory by introducing extra scalar freedom. In the Stueckelberg formalism, one attempts to understand the classical massive vector theory as a quantum massive vector theory, whereas our consideration here proposes to assign a physical meaning to a classical massive vector theory as a gauge fixed version of a massless gauge theory.
## 4 Dynamical generation of gauge fields
It is a long standing question if one can generate gauge fields from some more fundamental mechanism. In fact, there have been numerous attempts in the past to this effect. To our knowledge, however, there exists no definite convincing scheme so far. On the contrary, there is a no-go theorem or several arguments against such an attempt. We here briefly comment on this issue from a view point of our extended scheme of quantum gauge symmetry; our comment in this section is inevitably quite speculative.
Apart from technical details, the basic argument against the “dynamical” generation of gauge fields is that the Lorentz invariant positive definite theory cannot simply generate the negative metric states associated with the time components of massless gauge fields. In contrast, the dynamical generation of the Lagrangian of the structure
$$=\frac{1}{4}(_\mu A_\nu ^a_\nu A_\mu ^a+gf^{abc}A_\mu ^bA_\nu ^c)^2f(A_\mu ^a)$$
(4.1)
does not appear to be prohibited by general arguments so far. Here the term $`f(A_\mu ^a)`$ is Lorentz invariant but not invariant under the local gauge symmetry and thus breaks the gauge symmetry explicitly; $`f(A_\mu ^a)=\frac{m^2}{2}(A_\mu ^a)^2`$ is the simplest and most attractive example, since it carries the lowest scaling dimension. The appearance of $`f(A_\mu ^a)`$, which is not gauge invariant, is also natural if one starts with a fundamental theory without gauge symmetry.
Here comes the issue of interpretation of the induced Lagrangian (4.1). If one regards (4.1) as a quantum theory from the beginning, what one generates is simply a non-gauge theory: This is also the case if one evaluates a general S-matrix, which effectively represents the Lagrangian (4.1), and looks for the possible poles corresponding to massless gauge particles.
However, one might consider that the induced Lagrangian such as (4.1) is a classical object which should be quantized anew: In terms of path integral language, the Lagrangian is induced when one integrates over the “fundamental” degrees of freedom, and one need to perform further path integral over the induced Lagrangian anew. If one takes this latter view point, one might be allowed to regard the part of $`f(A_\mu ^a)`$, which breaks classical gauge symmetry, as a gauge fixing term in the modified quantization scheme. In this latter interpretation, one might be allowed to say that massless gauge fields are generated dynamically. Although a dynamical generation of pure gauge fields is prohibited, a gauge fixed Lagrangian might be allowed to be generated. (In this respect, one may recall that much of the arguments for the no-go theorem would be refuted if one could generate a gauge fixed Lagrangian with the Faddeev-Popov term added.) The mass for the gauge field which has an intrinsic unambiguous physical meaning is then further induced by the spontaneous symmetry breaking of the gauge symmetry thus defined (the Higgs mechanism).
We next comment on a mechanism for generating gauge fields by the violent random fluctuation of gauge degrees of freedom at the beginning of the universe; this scheme is based on the renormalization group flow starting from an initial chaotic theory. In such a scheme, it is natural to think that one is always dealing with quantum theory, and thus no room for our way of re-interpretation of the induced theory. Nevertheless, we find a possible connection in the following sense: To be precise, an example of massive Abelian gauge field in compact lattice gauge theory
$$𝒟U\frac{𝒟\mathrm{\Omega }}{vol(\mathrm{\Omega })}\mathrm{exp}[S_{inv}(U)S_{mass}(U^\mathrm{\Omega })]$$
(4.2)
is analyzed in Ref.. Here $`S_{inv}(U)`$ stands for the gauge invariant part of the lattice Abelian gauge field $`U`$, and $`S_{mass}(U^\mathrm{\Omega })`$ stands for the gauge non-invariant mass term with the gauge freedom $`\mathrm{\Omega }`$. In compact theory, one need not fix the gauge and instead one may take an average over the entire gauge volume of $`\mathrm{\Omega }`$. They argued that the mass term, which breaks gauge symmetry softly, disappears in the long distance limit when one integrates over the entire gauge freedom $`\mathrm{\Omega }`$. Their scheme is apparently dynamical one, in contrast to the kinematical nature of our re-interpretation. Nevertheless, the massive Abelian theory is a free theory in continuum formulation, and the disappearance of the mass term by a mere smearing over the gauge volume may suggest that the mass term in their scheme is also treated as a kind of gauge artifact, just as in our kinematical re-interpretation.
In passing, we note that the lattice gauge theory by using a mass term as a gauge fixing term has been discussed in the context of a numerical simulation. The main interest there is to evaluate gauge dependent quantities such as the propagator.
## 5 Discussion
On the basis of the equivalence between (1.1) and the local expression (1.3), we have discussed a possible more flexible interpretation of classical Lagrangians. In this broader interpretation of the classical action, the quantum gauge symmetry (BRST symmetry) is defined for a much wider class of theories than pure classical gauge theory, such as Maxwell field and Yang-Mills fields, in suitable quantization. In this framework, classical gauge symmetry is sufficient to generate quantum gauge symmetry (up to quantum anomalies), but it is not necessary in general; a theory, whose gauge symmetry is broken by some terms in the Lagrangian, could be re-interpreted as being gauge fixed by those terms in suitably quantized theory.
In this broader interpretation, the BRST symmetry is quite universal. This universality presumably arises from the fact that the essence of BRST symmetry is quite simple; geometrically, it is defined as the translation and scale transformations of a superspace coordinate specified by the real element of the Grassmann algebra
$`Q:\theta \theta +\lambda ,(BRSTcharge)`$
$`D:\theta e^\alpha \theta ,(ghostnumbercharge)`$ (5.1)
where $`\theta `$ and $`\lambda `$ are the real elements of the Grassmann algebra and $`\alpha `$ is a real number. Namely, the abstract BRST symmetry by itself carries no information of classical gauge symmetry.
We hope that the observation in the present note will stimulate further thinking on the meaning of classical and quantized gauge fields and also on the possible origin of gauge fields, the most profound notion of modern field theory.
The present paper is a revised and extended version of our contribution to ICHEP2000 in Osaka in July, 2000. One of us (KF) thanks H. Sugawara for a comment on a classical limit of quantized theory.
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# First-principles study of lattice instabilities in BaxSr1-xTiO3
## Introduction
Ferroelectric solid solutions with perovskite structure can exhibit properties which are much more interesting for technological applications than those of the related pure compounds. This is the case for the new class of single crystal relaxor ferroelectrics, exhibiting exceptionally high piezoelectric constants Park97 . It is also true for Ba<sub>x</sub>Sr<sub>1-x</sub>TiO<sub>3</sub> (BST), which combines a high dielectric constant with other properties making it a good candidate for the construction of dynamic random access memories (DRAM) Bilodeau97 .
In contrast to barium titanate (BaTiO<sub>3</sub>) which is a prototype ferroelectric (FE), strontium titanate (SrTiO<sub>3</sub>) is an incipient ferroelectric which undergoes an antiferrodistortive (AFD) phase transition at 105K Lines77 . Due to the different behavior of the related pure compounds, the structural properties of BST evolve with composition Lemanov96 . It exhibits a sequence of ferroelectric transitions similar to that of BaTiO<sub>3</sub> for Ba concentrations ranging from 20 to 100%. Close to pure SrTiO<sub>3</sub> the behavior is complicated by the competition between FE and AFD instabilities yielding the phase diagram reported in Ref. Lemanov96 .
In the past few years, first-principles effective Hamiltonians sucessfully described the temperature-dependent behavior of various pure perovskite oxides Vanderbilt97 . The generalization of this theoretical approach to solid solutions is a new challenge which includes some additional complications to treat the alloy disorder accurately.
A straightforward procedure is to try to extract the relevant information on the alloy from computation for various ordered supercells. This idea was applied to the study of Pb<sub>x</sub>Ge<sub>1-x</sub>Te solid solutions Cockayne98 but has the inconvenience of being computationally intensive. An alternative and promising approach is based on the Virtual Crystal Approximation (VCA) which allows the retention of a crystal with the primitive periodicity but composed of virtual atomic potentials averaging those of the atoms in the parent compounds. A careful construction of the VCA for PbZr<sub>x</sub>Ti<sub>(1-x)</sub>O<sub>3</sub> yielded a stress induced phase transition Ramer00 and a compositional phase boundary Ramer01 . Going further, the construction of a model Hamiltonian based on VCA calculations for the same compound recently was used to describe most of its temperature-dependent properties Bellaiche00 . In the latter case, the VCA calculation provided a reference structure to which corrections had to be added to treat the behavior of the real ions correctly.
If the purpose of the VCA is to define an average crystal from which the real solid solution can be described by including a restricted number of small corrections, one can ask if a similarly good estimate of the mixed compound could be obtained directly from the information available for the pure materials by averaging some key quantities. This idea is reinforced by the observation that the interatomic force constants seem very similar in the different ABO<sub>3</sub> compounds Ghosez99 .
In what follows we define an Average Crystal Approximation (ACA) by averaging the lattice constant ($`a_{cell}`$), the Born effective charges (Z), the optical dielectric tensor ($`ϵ_{\mathrm{}}`$), and the interatomic force constants in real space (IFC) of the pure compounds. We compare our results to those obtained (i) within the VCA and (ii) for a cubic ordered supercell with alternating planes of Ba and Sr along the (111) direction. We show that our ACA yields an excellent approximation to the dielectric and dynamical properties of the mixed compound. It allows the prediction of the lattice dynamics in the full range of composition without requiring any first-principles calculations other than those for the related pure compounds.
## Methodology
### General framework
Our calculations have been performed within the density functional theory Jones89 and the local density approximation, using the abinit package abinit . The ionic potentials screened by the core electrons have been replaced by highly transferable extended norm conserving pseudopotentials as proposed by M. Teter Teter93 . The Ba 5s, Ba 5p, Sr 4s, Sr 4p, Ti 3s, Ti 3p, Ti 3d, Ti 4s, O 2s and O 2p electrons have been treated as valence states. The wavefunction has been expanded in plane waves up to a kinetic energy cutoff of 45 Ha. The electronic exchange-correlation energy has been computed using a polynomial parametrizationGoedeker96 of Ceperley Alder data Ceperley80 . Integrals over the simple cubic Brillouin zone have been replaced by sums on a $`6\times 6\times 6`$ mesh of special $`k`$-points. A grid providing an equivalent sampling has been used for the supercell calculations.
The Born effective charges, the dielectric constants and the dynamical matrices have been computed within a variational formulation of the density functional perturbation theory Gonze92 . The interatomic force constants in real space (IFC) and the phonon dispersion curves have been extracted following the scheme described in Ref. Gonze94 , which deals correctly with the long range character of the dipolar interaction. The IFC were obtained from the knowledge of the dynamical matrix on a $`2\times 2\times 2`$ BCC mesh of $`q`$-points.
### Virtual atom pseudopotential
The virtual crystal approximation is based on the construction of a fictitious “virtual atom” potential by averaging the ionic potentials of the atoms alternating at the same site of the structure. As illustrated recently by Ramer and Rappe Ramer00 , there are different ways of combining non-local atomic pseudopotentials, each of them leading to different physical results. A careful choice seems necessary to predict accurately the alloy properties at the VCA level. However, this choice could be less stringent if we think about the virtual crystal as a reference structure from which the real crystal can be obtained by including a limited number of relevant corrections.
For BST, we choose therefore to stay with the most straightforward combination of the reference atomic pseudopotentials. The potential at the A-site can be formally defined for a concentration of $`x`$% of barium as:
$$V_{VCA}^{ps}[x]=xV_{Ba}^{ps}+(1x)V_{Sr}^{ps}$$
(1)
where $`V_{Ba}^{ps}`$ and $`V_{Sr}^{ps}`$ are the pseudopotentials of Ba and Sr atoms.
Non-local pseudopotentials, such as those we are using, can be expressed as a sum of a local contribution and short-range non-local corrections:
$$V^{ps}(r,r^{})=V^{loc}(r)\delta (rr^{})+\underset{l,m}{}\underset{q}{}\frac{V_l^{ps}|\varphi _{lm,q}><\varphi _{lm,q}|V_l^{ps}}{<\varphi _{lm,q}|V_l^{ps}|\varphi _{lm,q}>}.$$
(2)
Our virtual atom pseudopotential is then defined as follows Desquesnes :
$`V_{VCA}^{ps}(r,r^{})`$ $`=`$ $`[xV_{Ba}^{loc}(r)+(1x)V_{Sr}^{loc}(r)]\delta (rr^{})`$ (3)
$`+`$ $`{\displaystyle \underset{l,m}{}}{\displaystyle \underset{q}{}}{\displaystyle \frac{xV_{Ba,l}^{ps}|\varphi _{lm,q}^{Ba}><\varphi _{lm,q}^{Ba}|xV_{Ba,l}^{ps}}{<\varphi _{lm,q}^{Ba}|xV_{Ba,l}^{ps}|\varphi _{lm,q}^{Ba}>}}`$
$`+`$ $`{\displaystyle \underset{l,m}{}}{\displaystyle \underset{q}{}}{\displaystyle \frac{(1x)V_{Sr,l}^{ps}|\varphi _{lm,q}^{Sr}><\varphi _{lm,q}^{Sr}|(1x)V_{Sr,l}^{ps}}{<\varphi _{lm,q}^{Sr}|(1x)V_{Sr,l}^{ps}|\varphi _{lm,q}^{Sr}>}}`$
For $`x`$ ranging from 0 to 1, our virtual atom evolves smoothly from Sr to Ba.
## Ground-state properties
First we focus on the lattice constant of the cubic perovskite structure and look at its evolution with composition. In Table 1, we report the LDA lattice constants for BaTiO<sub>3</sub> and SrTiO<sub>3</sub>, for Ba<sub>1-x</sub>Sr<sub>x</sub>TiO<sub>3</sub> within the VCA (for x = 0.25, 0.50 and 0.75), and for the (111) ordered supercell. We compare the computed values with experimental results, which at intermediate compositions are accurately given by Vegard’s law Lemanov96 . While our results show the lattice constant underestimate typical of LDA, the Vegard’s law behavior is well reproduced.
We note that the theoretical LDA underestimate of the experimental lattice constant fortuituously corresponds to the effect of a similar external pressure of 8.7 GPa in both BaTiO<sub>3</sub> and SrTiO<sub>3</sub>. In what follows, we will work at the experimental lattice constants which is therefore equivalent to imposing a fixed fictitious negative pressure of -8.7 GPa to the LDA results. For the mixed compounds we will work at the values obtained by Vegard’s law as reported on the third line of Table 1.
## Dielectric and dynamical properties
### Born effective charges and dielectric tensor
In Table 2, we compare the Born effective charges and the macroscopic optical dielectric constants obtained in our various calculations. The values of O and O refer respectively to the change of polarization for a displacement of the oxygen atom along the Ti–O direction and perpendicular to it.
We observe that the values of Z and $`ϵ_{\mathrm{}}`$ in the supercell are very similar to both those averaged from the pure compounds (ACA) and those obtained in the VCA. In particular, the ACA and the VCA similarly reproduce the large anomalous values of $`Z_{Ti}^{}`$ and $`Z_O^{}`$ which play an important role in generating the FE instability. For Z$`{}_{A}{}^{}{}_{}{}^{}`$, both kinds of approximation yield the average of the values computed for Ba and Sr in the supercell. A similar observation was reported by Bellaiche et al. for PZT Bellaiche00 .
### Phonon frequencies and interatomic force constants
Finally we compare the phonon dispersion curves, yielding information at the harmonic level on the full lattice dynamics of the crystal. Phonon frequencies are obtained by diagonalizing the dynamical matrix. The latter is usually decomposed into a dipolar part computed from the values of Z and $`ϵ_{\mathrm{}}`$, and a short-range part computed using the IFC. Our ACA dynamical matrix is constructed by simply averaging all the previous quantities from their values in the pure compounds and using an average mass at the A-site, while the VCA requires a full set of additional first-principles calculations with the virtual atom potential.
In Fig. 1, we show the phonon dispersion curves of a mixed crystal of BST with x=50%, obtained both within the ACA and the VCA. We compare the results to those obtained for the related pure compounds. The dispersion curves of BaTiO<sub>3</sub> were discussed in Ref. Ghosez98 . Those of SrTiO<sub>3</sub> have been recomputed and agree with the calculation reported previously by LaSota et al. LaSota97 . Compared to BaTiO<sub>3</sub> which exhibits only a FE instability with a flat dispersion in the $`\mathrm{\Gamma }`$–X–M plane, the FE instability in SrTiO<sub>3</sub> is more localized in $`q`$-space (around the $`\mathrm{\Gamma }`$ point) and coexists with an AFD instability extending along the Brillouin zone edges (R–M lines).
We observe that averaging the different quantities (Z, $`ϵ_{\mathrm{}}`$, IFC, atomic masses) does not yield an average of the phonon frequencies but rather is closer to the results for pure BaTiO<sub>3</sub>. This is also what is found in the VCA, the two phonon dispersions being very similar. The major discepancy appears for the low frequency Ti dominated branch along the $`\mathrm{\Gamma }`$–M–X lines where the value of $`\omega ^2`$, being near zero, is especially sensitive to small changes in the dynamical matrix.
In Table 3, we compare the ACA and VCA phonon frequencies at $`\mathrm{\Gamma }`$ and R in the reference perovskite structure to the related frequencies at the $`\mathrm{\Gamma }`$ point in the (111) ordered supercell. The different symmetry of the supercell does not allow a one to one correspondence but the comparison is made between phonons having the highest eigenvector overlap. We note that the supercell eigenmode associated with the R<sub>15</sub> modes of the perovskite structure, and involving displacement of the A atom, exhibits a small LO-TO splitting.
The overall agreement between ACA, VCA and supercell calculations is very good. This is particularly true for the $`\mathrm{\Gamma }_{15}`$ and R$`_{25^{}}`$ modes related to the FE instability. There is a larger discrepancy for the R<sub>25</sub> AFD mode. We note however that projecting the supercell dynamical matrix within the subspace of the VCA eigenvector yields a value of 125 cm<sup>-1</sup> (instead of 86 cm<sup>-1</sup>) closer to VCA and ACA results. This shows that the discrepancy originates in a change of the eigenvector rather than in a modification of the interatomic force constants.
### Evolution of the phonon instabilities with composition
As the ACA seems to give accurate results for the lattice dynamics of BST solid solutions, it offers an easy way to make a guess for the evolution of the FE and AFD instabilities with composition. We observe in Figure 2 that the square of the frequency of the R<sub>25</sub> AFD mode evolves linearly with composition and becomes unstable for Ba concentrations smaller than 20 $`\%`$. This corresponds precisely to the composition for which a change of behavior is observed experimentally in the phase diagram of BST Lemanov96 . In contrast, the FE instability exhibits a non-linear behavior and remains unstable for the full range of composition.
## Conclusions
In this paper, we have compared different approaches for the study of the lattice dynamics of BST solid solutions and have introduced a new Average Crystal Approximation. We have shown that the latter yields results equivalent to the VCA but with the advantage that it does not require any additional first-principles calculations if the information for the related pure compounds is available. Looking at the actual values of the IFC in the supercell could provide indications of the corrections to be included in the harmonic part of an effective Hamiltonian based on the ACA for the correct description of BST solid solutions.
## Acknowledgments
This work was supported by ONR grant No N00014-97-0047. Supercomputing support was provided by MHPCC. X.G. acknowledges financial support from the F.N.R.S. (Belgium), the F.R.F.C. (project 2.4556.99), the P.A.I. P4/10.
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# Pearling Instabilities of Membrane Tubes with Anchored Polymers
\[
## Abstract
We have studied the pearling instability induced on hollow tubular lipid vesicles by hydrophilic polymers with hydrophobic side groups along the backbone. The results show that the polymer concentration is coupled to local membrane curvature. The relaxation of a pearled tube is characterized by two different well-separated time scales, indicating two physical mechanisms. We present a model, which explains the observed phenomena and predicts polymer segregation according to local membrane curvature at late stages.
PACS numbers: 87.16.Dg, 68.10.-m \]
Single-component phospholipid membranes have been the focus of intense interest in recent years, as the simplest model system of biological membranes . The latter are highly complex systems, comprised of a bilayer consisting of many types of lipids as well as a mesh of macromolecules such as proteins and polysaccharides, which participate in a large variety of cellular processes. A natural step to mimic this complexity in a simple model system is to study the association of polymeric molecules with self-assembled single-component phospholipid membranes.
Experiments with polymers, which associate with membranes by anchoring, have revealed changes in the bending moduli of bilayers of single-tailed surfactants , and striking morphological changes in vesicles. Anchoring occurs by the penetration of a number of hydrophobic side-groups grafted along a hydrophilic backbone into a bilayer. Hollow tubular vesicles incubated in a solution of anchoring polymers having a polysaccharide backbone develop a pearling instability, above a threshold polymer concentration. This instability was not observed with purely hydrophilic polymers, and was effected by hydrophobic groups alone (without the backbone), but only at concentrations five orders of magnitude higher than with anchoring polymers.
In this experimental and theoretical study we investigate the mechanisms responsible for pearling in our system. It has been suggested that the induction of curvature by the anchors, which sink into the membrane to a depth of half a bilayer, may constitute a mechanism which drives the pearling instability . We present a novel nonequilibrium experiment, which shows that two independent mechanisms – spontaneous curvature and area difference – contribute to the pearling phenomenon. We also present new experimental and theoretical findings showing inhomogeneous shapes at late stages.
In our experiments, vesicles were made of stearoyl-oleoyl-phosphatidyl-choline (SOPC) with $`C_{18}`$ alkyl chains. The polymer used was hydrophilic dextran with a molecular weight of $`162,000g/mol`$, functionalized both with palmitoyl alkyl chains $`C_{16}`$ and dodecanoic NBD chains as fluorescent markers. The anchors are distributed statistically along the backbone, spaced four persistence lengths apart on average (1 alkyl chain per 25 glucose units). $`1\mu l`$ droplets of SOPC in a 7:1 chloroform-methanol solution ($`10mg/ml`$) were placed on a glass coverslip forming small lipid patches. The sample was prehydrated for 20 minutes under water-saturated nitrogen, then hydrated with $`0.1mM`$ potassium buffer at pH 6.5. A large number of hollow tubes with one or more lamellae formed after hydration, connected to the lipid patch. The latter constitutes a reservoir with which the tubes can exchange lipid molecules. Experiments were conducted at room temperature, insuring that the membranes were in a fluid-like state. This allowed free diffusion of both lipids and anchored polymers along the bilayers. A drop of polymer solution of a given concentration was introduced through one of the cell sides. The polymer was added after the tubes were formed, and we therefore assume that it anchored mostly on the outermost leaflet of the membrane. Its concentration on the membrane grew from zero as more and more chains anchored from solution. Events were observed by phase contrast and fluorescence microscopy and recorded on video. The NBD markers were excited with an argon laser ($`488nm`$), and observed with a CCD camera.
Fig. 1 shows a tube undergoing pearling after addition of polymer. The instability typically starts near the end cap of the tube (Fig. 1a), and gradually propagates along the axis (Fig. 1b). This is presumably since the polymer, diffusing from the side of the chamber, first reaches the tip of the tube. Existing pearls become gradually more spherical (Fig. 1c). The rate at which pearls form depends on the number of lamellae in the walls of the tube, as well as on the concentration of polymer in the water ($`2100mg/l`$), varying from several seconds to many minutes. A salient feature of Fig. 1c is the increase in pearl size towards the end of the tube. This size gradient becomes much more pronounced at very long times, as Fig. 1d illustrates for another vesicle. The string of pearls separates into a set of small, nearly-uniform spheres connected to a group of much larger spheres.
We have studied the onset of pearling by measuring the wavelength of the fastest growing mode, $`P`$, just above threshold, as a function of the radius of the unpearled tube, $`R_0`$. Within experimental accuracy $`P`$ is linear in $`R_0`$, with a slope $`k2\pi R_0/P=1.07\pm 0.05`$. This result agrees well both with the value $`k=1`$ predicted by theoretical models based on induced curvature as the driving mechanism , and with experimental results on pearling instabilities in single-tailed surfactant systems. In tension-induced pearling $`k0.7`$ , and thus our measurement excludes pearling due to polymer or flow-induced tension in our system.
Polymers anchoring to one side of a membrane can induce curvature by two mechanisms. The first is an increase in the area of the outermost monolayer into which anchors sink . The second mechanism is a local deformation of the membrane which can be induced either by the anchors regarded as inclusions , or by an entropic pressure exerted by the polymer backbone . These mechanisms form the basis of two models that describe the tendency of a membrane to display curvature : the area difference elasticity model (ADE) and the spontaneous curvature model (SC), respectively. Calculations of equilibrium shapes of vesicles with cylindrical symmetry based on both the SC and ADE models indeed yield pearled shapes of constant mean curvature called Delaunay surfaces . The equilibrium shape is not sensitive to the pearling mechanism.
Nevertheless, the dynamics of pearling may allow us to distinguish between the two mechanisms and reveal their presence, since they are characterized by well-separated relaxation time scales. Inhomogeneities in both area difference and spontaneous curvature decay diffusively . The relaxation of spontaneous curvature is associated with polymer diffusion in the membrane. The relevant diffusion constant was measured for various macromolecules and falls in the range of $`1\mu m^2/sec<D_{sc}<5\mu m^2/sec`$ . Area difference relaxes via the sliding of one monolayer with respect to the other, and does not involve diffusion of molecules over large distances. The diffusion constant associated with ADE can be estimated from dimensional analysis to be $`D_{ade}K_0/b`$, where $`K_0`$ is the compression modulus of the membrane, and $`b`$ is the friction coefficient between the two leaflets of a bilayer. Experimental estimates of these parameters lead to $`50\mu m^2/sec<D_{ade}<500\mu m^2/sec`$ . Changes in the shape of a vesicle are coupled to motion of the surrounding water, where the energy dissipated in the water is of the order of the curvature energy of the bilayer. From dimensional analysis we find that for both mechanisms, ADE and SC, the diffusion constant associated with the displacement of water is $`D_{hyd}\kappa /\eta R_050\mu m^2/sec`$, where $`\kappa `$ is the bending modulus of the membrane and $`\eta `$ is the viscosity of water. Thus, it is difficult to differentiate between hydrodynamical and ADE effects by measuring the diffusion constant, but it should be easy to distinguish between the SC and ADE mechanisms.
To study the relaxational dynamics, a micropipette was used to deliver locally a small volume ($`10^4\mu l`$) of polymer solution of concentration $`20100mg/l`$ close to a tube (Fig. 2). Fluorescence images show that the polymer was indeed concentrated in a local region shortly after the injection. Pearling occurred, after an induction time in which enhanced undulations were observed (Figs. 2a-2d). The injection of the polymer was stopped after the pearled region increased to a length of order $`100\mu m`$. The pearled region then shrank gradually (Figs. 2e-2h) as pearls opened up one at a time starting with those farthest from the region of polymer injection. This decay process took several minutes. Fluorescence images show that the amount of polymer in the surrounding water is negligible during the decay.
We measured the length of the pearled region, $`L`$, as a function of time during the decay process. For diffusive decay, we expect $`L^22D(\overline{t}t)`$, where $`D`$ is the relevant diffusion constant and $`\overline{t}`$ is a constant. Figure 3 is a typical example of the dependence of $`L^2`$ on time (measurements on several tubes yielded similar results). The system exhibits diffusive behavior at early times with a diffusion constant of $`D200\mu m^2/sec`$. We attribute this decay to a combination of ADE and hydrodynamical effects. At later times there is a sharp crossover to a much slower diffusive behavior with $`D5\mu m^2/sec`$. According to our estimates, this corresponds to polymer diffusion, and is associated with the SC mechanism. Our results thus provide clear evidence that both ADE and SC mechanisms influence the pearling instability.
We now carry out a theoretical analysis of the pearling phenomena in the case of global application of the polymer (Fig. 1). We consider closed vesicles with polymer molecules only on the outer side of the vesicle, and assume that pearling is a result of the SC mechanism (All our predictions, except for the inhomogeneities in polymer concentration, apply equally well to the ADE model). In contrast with standard curvature models , the spontaneous curvature in our system is a local quantity. We assume that the spontaneous curvature is proportional to the polymer concentration on the membrane, $`\rho H_0`$. $`\rho (\stackrel{}{r})`$ is the fraction of the membrane area covered by polymer molecules at the position $`\stackrel{}{r}`$, and takes values in the interval $`[0,1]`$. $`H_0`$ is the spontaneous curvature induced by full coverage of the polymer. We further assume that the observed vesicle shapes are close to the equilibrium shapes under the constraints of constant vesicle volume and membrane area. These shapes can be obtained by minimizing the free energy of the membrane in the presence of the polymer. The simplest free energy for our system is a sum of the curvature energy and the entropy of mixing of the polymer:
$`F={\displaystyle }`$ $`dA`$ $`\{2\kappa (H\rho H_0)^2+`$ (1)
$`{\displaystyle \frac{k_BT}{a^2}}`$ $`[\rho \mathrm{ln}\rho +(1\rho )\mathrm{ln}(1\rho )]\},`$ (2)
where $`H`$ is the local mean curvature, $`a`$ is the characteristic linear size of an anchored polymer molecule and the integration is over the area of the membrane. In principle, the effects of gradients of polymer concentration should also be included in the free energy. However, the inhomogeneities at equilibrium do not induce an extensive free energy increase (see below), and we therefore ignore such terms.
We have measured $`\kappa (20\pm 5)k_BT`$ using the pipette aspiration technique . Structures with radii smaller than optical resolution ($`0.2\mu m`$) were observed in our experiments, giving a lower bound of $`H_010\mu m^1`$ on the spontaneous curvature induced by the polymer. Finally, $`a`$ can be estimated as the radius of gyration of a polymer performing a two dimensional random walk on the membrane. Since the hydrophilic backbone is in water, we assume a good solvent in semi-dilute conditions, giving a radius between 40 and $`80nm`$.
We consider very long, cylindrically symmetric vesicles and ignore the existence of the end caps, since the length of most of the experimental tubes is larger than their radii by two orders of magnitude. To find the equilibrium configuration of the system, the free energy (2) was minimized with respect to the vesicle shape as well as the local polymer density. This was done under the constraints of constant vesicle volume, membrane area and total number of polymer molecules.
Ideally, the equilibrium configuration of the system would have a homogeneous polymer distribution and a curvature $`H=\rho H_0`$ everywhere. Such a configuration minimizes the energy and maximizes the entropy simultaneously. This is indeed possible for a range of values of $`\rho H_0`$. For example, a vesicle with volume-to-area ratio $`\lambda `$ can have the shape of a cylinder of radius $`2\lambda `$ and curvature $`H=1/(4\lambda )`$. Another shape, having the same value of $`\lambda `$, is a chain of identical spheres of radius $`3\lambda `$ connected by infinitesimally narrow necks. In this case $`H=1/(3\lambda )`$ everywhere. In fact, the vesicle’s curvature may take any intermediate value, $`1/(4\lambda )<H<1/(3\lambda )`$, because for each of these curvatures there corresponds a Delaunay shape , with the same value of $`\lambda `$. On the other hand, it is not possible to construct a shape of constant curvature for $`H>1/(3\lambda )`$ or for $`H<1/(4\lambda )`$.
In our experiment the polymer adsorbs onto the membrane gradually. Therefore, at the very early stages of the experiment $`\rho <1/(4\lambda H_0)`$, and the shape of constant curvature which minimizes the energy is a cylinder of radius $`2\lambda `$ and curvature $`H=1/(4\lambda )`$. This is consistent with our experimental observations. In principle, we also have to consider possible inhomogeneities in membrane curvature. However, our calculations show that for small $`\rho `$ such inhomogeneities only increase the free energy.
At intermediate stages of the experiment the polymer concentration is in the range $`1/(4\lambda H_0)\rho 1/(3\lambda H_0)`$, and long vesicles are expected to have a Delaunay shape. Indeed, the shapes of vesicles we observe at intermediate stages of the experiment are similar to Delaunay shapes. Such pearled shapes have already been observed , and the importance of Delaunay shapes for this system has been discussed in the literature .
The situation becomes more interesting at late stages of the experiment when $`\rho `$ exceeds the value $`1/(3\lambda H_0)`$. In this case, the best Delaunay shape is a chain of identical spheres of radius $`3\lambda `$ and curvature $`H=1/(3\lambda )`$. However, a detailed calculation (to be presented elsewhere) shows that a configuration consisting of a chain of small spheres connected by infinitesimal necks to each other and to one large sphere has the lowest free energy.
The small spheres have polymer density $`\rho _1`$ and radius $`r_1=1/(\rho _1H_0)`$. The energy of this subsystem vanishes, since its curvature is $`H_1=\rho _1H_0`$. The large sphere has radius $`r_2`$ and polymer density $`\rho _2`$, and plays the role of a reservoir for the excess volume and polymer molecules. The ratios $`r_2/r_1`$ and $`\rho _2/\rho _1`$ can be calculated numerically as functions of the average polymer density, the area of the system and its volume. For the values of the model parameters discussed above we obtained $`r_2/r_1>10`$ and $`\rho _2/\rho _1<0.3`$. This ratio of the radii is consistent with the experimental configurations seen at long times, where chains of very small spheres coexist with few very large spheres, all connected by narrow necks (see Fig. 1d).
The strong inhomogeneity in polymer concentration predicted by the theory, is particularly interesting because it may show a qualitative difference between predictions of the local SC model and those of ADE models. The latter does not differentiate between polymer molecules and lipids; i.e., exchanging polymers with a larger number of lipids of the same total area does not change the energy. One can always use such exchanges to turn an inhomogeneous polymer distribution into a homogeneous one of higher entropy without changing the energy of the vesicle. Hence, ADE models predict a homogeneous polymer distribution. In the local SC model, inhomogeneities in polymer distribution induce the same undesirable decrease in the entropy of mixing. However, inhomogeneities lower the energy, and our calculations show that for reasonable parameter values, the reduction in energy does lead to sizable inhomogeneities. We intend to measure the polymer concentration on the membrane to test this interesting theoretical prediction.
We thank L. Jullien, for his help and encouragement, and R. Lipowsky, E. Moses and S. Safran for useful exchanges. This research was supported by The Israel Science Foundation - Recanati and IDB Group Foundation and The Minerva Foundation. V.F. acknowledges support from The Research Council of Norway (NFR).
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# Quantum superconductor-metal transition
## I Introduction
Quantum superconductor-insulator (or superconductor-normal metal) phase transitions can take place at zero temperature due to variation of parameters of the system. For example in experiments the transition takes place as a function of the degree of disorder in a superconducting film. In experiments of the transition was mediated by a magnetic field.
It has been suggested that a disordered superconducting film can be described as a network of Josephson-coupled superconducting grains shunted by resistors . In this case the Coulomb interaction between electrons in grains suppresses fluctuations of the number of electrons in a grain and, due to uncertanty principle, increases the amplitude of fluctuations of the phase of the superconducting order parameter . The competition between the charging energy and the Josephson inter-grain coupling energy leads to the phase transition. Shunting resistors play a double role in the model: a) dissipation in the resistors tends to suppress fluctuations of the phase of the order parameter b)tunneling between superconducting grains and the resistors renormalizes the capacitance of the grains and thereby the charging energy. As a result, in the framework of this model and in two dimensional case the onset of superconductivity corresponds to the normal state film’s conductance $`G`$ of the film of order of $`\frac{e^2}{h}`$. Since at $`G<\frac{e^2}{h}`$ the system should be in the insulating state, it has been conjectured that the transition is of the superconductor-insulator type and that the behaviour of the system near the transition is universal. On the other hand, a renormalization group analysis , which starts from a perturbation theory for slightly disordered uniform superconductors, showed in 2D case a zero temperature superconductor-metal transition.
Recently Feigelman and Larkin reconsidered the problem in the framework of a model of superconducting grains embedded in a normal metal film. They showed that a) the superconductor-normal metal transition takes place and b) deep in the metallic phase parameters of the system can be calculated with the help of perturbation theory. At $`G>\frac{e^2}{h}`$, however, the critical concentration of grains turned out to be exponentially small.
On the other hand, the theories concerning granular superconductors were done in the limit when the modulus of the order parameter on a grain does not fluctuate. In the case of no reflection on the superconductor-normal metal boundary this corresponds to $`R\xi `$. Here $`R`$ is the grain’s radius, $`\xi =min[\frac{v_F}{\mathrm{\Delta }_0};\sqrt{\frac{D}{\mathrm{\Delta }_0}}]`$ is the zero temperature coherence length of the bulk superconductor, $`D=\frac{v_Fl}{3}`$ is the electron diffusion coefficient in the metal, $`v_F`$ is the Fermi velocity, $`l`$ is the electron elastic mean free path and $`\mathrm{\Delta }_0`$ is the zero temperature value of the gap in the bulk superconductor.
In this paper we consider the opposite case when $`R<\xi `$ and show that the zero temperature quantum superconductor-metal transition as a function of grains’ concentration can take place in samples with arbitrary large conductance and can exist even in the absence of any disorder in the sample and in the approximation when electrons in the normal metal do not interact. The critical concentration of the grains turns out to be relatively large.
We will consider a quasi-two-dimensional film of thickness $`a\xi `$, which consists of superconducting grains of radius $`R`$ embedded in a nonsuperconducting metal. The results can be easily generalized for d=3 case. We will assume that densities of states $`\nu `$ of the superconductor and the metal are the same and that the spatial dependence of the electron-electron interaction constant in the the sample has the form
$$\lambda (𝒓)=\{\begin{array}{cc}\lambda _s>0\hfill & \text{if }|𝒓𝒓_i|<R\hfill \\ \lambda _N<0\hfill & \text{if }|𝒓𝒓_i|>R\hfill \end{array}$$
(1)
Here $`\lambda _s`$ and $`\lambda _N`$ are electron interaction constants in the superconductor and in the metal respectively; $`𝒓_i`$ are coordinates of centers of the superconducting grains and index $`i`$ labels the grains. At zero temperature the linearized mean field equation for the order parameter $`\mathrm{\Delta }(𝒓)`$
$$\mathrm{\Delta }(𝒓)=\lambda (𝒓)𝑑𝒓_1K(𝒓,𝒓_1)\mathrm{\Delta }(𝒓_1)$$
(2)
has a solution $`\mathrm{\Delta }(𝒓)=f(|𝒓𝒓_i|)`$ at $`R=R_c^{(mf)}`$. Assuming no reflection at the superconductor-metal boundary we get $`R_c^{(mf)}\xi `$ . Here $`K(𝒓,𝒓_1)=𝑑ϵG(ϵ,𝒓,𝒓_1)G(ϵ,𝒓,𝒓_1)`$, $`G(𝒓,𝒓_1)`$ is the electron Green’s function in the normal metal. It is convenient to normalize $`𝑑𝒓|f(𝒓)|^2=1`$. At $`R<R_c^{(mf)}`$ the mean field value of the order parameter is zero ($`\mathrm{\Delta }(𝒓)=0`$). To find the value of $`\mathrm{\Delta }(𝒓)`$ in the case $`R>R_c^{(mf)}`$ one has to add to Eq.2 terms nonlinear in $`\mathrm{\Delta }`$ .
## II Quantum fluctuations of the order parameter in an individual grain
We will consider the case $`R<R_c^{(mf)}`$ when for an individual grain Eq.2 does not have nonzero mean field solution. To calculate the correlation function of quantum fluctuations of the order parameter it is convenient to use a parametrization $`\mathrm{\Delta }_i(𝒓,t)=\alpha _i(t)f(|𝒓𝒓_i|)`$, which reflects the fact that, at $`(R_c^{(mf)}R)R_c^{(mf)}`$, the amplitude of quantum fluctuations of $`\alpha _i(t)`$ is large while the amplitude of fluctuations of the shape of $`f(r)`$ is small. To describe dynamics of the order parameter in a grain we use the effective action
$$S_i=\nu \tau _0\frac{d\omega }{2\pi }(i|\omega |+\frac{1}{\tau })|\alpha _i(\omega )|^2$$
(3)
whose derivation we outline in the Appendix. Here $`\tau _0=min[\frac{r}{v_F};\frac{R^2}{D}]`$ is the the time of electron flight through the grain and
$$\tau =\frac{\tau _0R}{R_c^{(mf)}R}$$
(4)
Using Eq.3 we get
$$<\alpha _i(\omega ),\alpha _i^{}(\omega )>_{(0)}=\frac{1}{\nu \tau _0(i|\omega |+\frac{1}{\tau })}$$
(5)
which in $`t`$ representation corresponds to
$$<\alpha _i(t),\alpha _i^{}(0)>_{(0)}=\{\begin{array}{cc}\frac{1}{\nu \tau _0}(\frac{\tau }{t})^2\hfill & \text{if }t\tau \hfill \\ \frac{1}{\nu \tau _0}i[i\pi +2ln(\frac{t}{\tau })]\hfill & \text{if }t\tau \hfill \end{array}$$
(6)
Here the subscript $`\mathrm{"}(0)\mathrm{"}`$ indicates that the correlation function is calculated for a single existing grain $`\mathrm{"}i\mathrm{"}`$. Eqs.5,6 hold as long as terms nonlinear in $`|\alpha _i|^2`$ in the effective action can be neglected, i.e. if $`\tau \delta ^1`$. Here $`\delta =(\nu R^2a)^1`$ is the average level spacing in the grain. They correspond to the casual Green’s function $`G_c=<T(\mathrm{\Psi }_\sigma (𝒓,t)\mathrm{\Psi }_\sigma (𝒓,t)\mathrm{\Psi }_\sigma ^+(0,0)\mathrm{\Psi }_\sigma ^+(0,0))>`$, where $`\sigma `$ is the spin index. To get the retarded Green’s function $`G_R=<\theta (t)[\mathrm{\Psi }_\sigma ^+(𝒓,t)\mathrm{\Psi }_\sigma ^+(𝒓,t),\mathrm{\Psi }_\sigma (0,0)\mathrm{\Psi }_\sigma (0,0)]_+>`$ one has to make analytical continuation of Eq.5 with respect to $`\omega `$ and then to make a Fourier transform. As a result, $`G_R(t)\mathrm{exp}(\frac{t}{\tau })`$. It is interesting that the asymptotic time dependence of Eq.6 is the same as the one obtained in the case $`RR_c^{(mf)}`$ with the help of a complicated renormalization group analysis of the Caldeira-Leggett effective action . In the latter case there is a non-zero mean field order parameter on a grain and the correlation function decays with time due to phase fluctuations mediated by the interaction with quantum electromagnetic fluctuations in conducting environment. At $`G>\frac{e^2}{h}`$, however, the corresponding correlation time turns out to be exponentially large, which is different from the case $`R<R_c^{(mf)}`$ Eq.4.
## III The superconductor-metal transition in a system of superconducting grains
Let us consider now the case when the concentration of superconducting grains $`N`$ embedded into the quasi-two-dimensional metallic film is finite. To describe this system in the case when $`\alpha _i`$ are small we will use the effective action $`S=S_0+S_{int}`$ where $`S_0=_iS_i`$, while $`S_{int}`$ describes inter-grain interaction via the normal metal and also the nonlinear in $`|\alpha _i|^2`$ contributions to the action.
$$S_{int}=\underset{i,j}{}𝑑\omega J_{ij}\alpha _i(\omega )\alpha _j^{}(\omega )+b\underset{i}{}𝑑\omega _1𝑑\omega _2𝑑\omega _3\alpha _i(\omega _1)\alpha _i^{}(\omega _2)\alpha _i(\omega _3)\alpha _i^{}(\omega _1\omega _2\omega _3)$$
(7)
where $`J_{ij}`$ have the meaning of Josephson coupling between grains $`i`$ and $`j`$, and $`b\frac{R_c^2\nu }{aD^2}`$.
The mean field approximation corresponds to the minimum of $`S`$ at $`\omega =0`$. If $`|𝒓_i𝒓_j|`$ is small enough the solution of the Uzadel equation Eq.25 yields Josephson coupling between two grains of the form ($`|𝒓_i𝒓_j|R`$)
$$J_{ij}=\frac{\nu R^2}{|𝒓_i𝒓_j|^2}$$
(8)
However, in the case of finite two-dimensional concentration $`N`$ such an expression for $`J_{ij}`$ would lead to a logarithmic divergence of the ground state energy density.
On the other hand, at large $`|𝒓_i𝒓_j|`$ electron-hole pairs diffusing through the metal between grains $`i`$ and $`j`$ will experience Andreev scattering from the superconducting grains situated between grains $`i`$ and $`j`$. An example of such a grain $`\mathrm{"}k\mathrm{"}`$ is shown in Fig.1. Due to the Andreev nature of reflection electrons scattered by the grain $`\mathrm{"}k\mathrm{"}`$ will be reflected into a hole moving in the direction opposite to the initial direction of electron’s motion. Thus this electron will never reach the grain $`j`$. As a result we have
$$J_{ij}=\frac{\nu R^2}{|𝒓_i𝒓_j|^2}\mathrm{exp}(\frac{|𝒓_i𝒓_j|}{L_0^{(mf)}})$$
(9)
$$L_0^{(mf)}=(\frac{l}{NR}\frac{Ra^{1/2}\mathrm{\Delta }_0}{\alpha _{(mf)}})^{1/2}$$
(10)
Therefore the mean field equation for the ground state mean field value $`\alpha _{(mf)}`$ of the order parameter on a grain has a form
$$\frac{\nu \tau _0}{\tau }2\pi \nu R^2N\mathrm{ln}(N^{1/2}L_0^{(mf)})+b\alpha _{(mf)}^2=0$$
(11)
Thus at $`T=0`$ an exponentially small mean field solution for the order parameter
$$\alpha _{(mf)}\mathrm{\Delta }_0la^{1/2}\mathrm{exp}(\frac{\tau _0}{\tau R^2N})$$
(12)
exists at arbitrary small $`N`$.
However, at small enough $`N`$ the amplitude of quantum fluctuations of the order parameter Eqs.5,6 becomes larger than its mean field value. In this case the mean field theory is not applicable. To show that at small $`N`$ quantum fluctuations destroy superconductivity completely and that the normal metal state is stable we use a perturbation theory procedure similar to the one used in . A requirement for the perturbation theory in therms of $`J_{ij}`$ to be valid is the convergence of the integral
$$<\alpha _i(t)\alpha _i^{}(0)>𝑑t<\mathrm{}$$
(13)
which in our case follows from Eq.6. The integral in this case equals $`\frac{\tau }{\nu \tau _0}`$.
Again, the Josephson couplings of the form Eq.8 would lead to divergency of the perturbation theory. To cut off the divergence in the absence of the magnetic field one can consider the case when there is a repulsion between electrons in the metal and $`\lambda _N0`$ . Then in the two-dimensional case we have $`J_{ij}\nu R^2r^2[1+2\nu |\lambda _N|ln(\frac{r}{R})]^2`$ and the perturbation theory will converge on the lengthscale $`L_\lambda R\mathrm{exp}(\frac{1}{\nu |\lambda _N|})`$. In the presence of a weak magnetic field $`H`$ the Josephson intergrain coupling decays exponentially on distances larger than the magnetic length $`L_H=\sqrt{\frac{m}{eH}}`$. Thus the cut off length relevant for the convergence of the perturbation theory with respect to the term in Eq.7 proportional to $`J_{ij}`$ is $`L_0=\mathrm{min}[L_\lambda ;L_H]`$. It gives small corrections to the Eqs.5, 6 as long as $`NR^2ln(L_0N^{1/2})\frac{\tau _0}{\tau }`$. In the opposite case the ground state of the system is superconducting. Thus we can estimate a relation between the critical concentration of the grains $`N_c`$ and their critical radius $`R_c`$ from the equation
$$\frac{|R_c^{(mf)}R_c|}{R_c^{(mf)}}R_c^2N_cln(L_0N_c^{1/2})$$
(14)
For example, in the case $`R<R_c^{(mf)}`$ and at $`H=0`$ we have the estimate
$$N_c\frac{1}{R_c^2}\nu |\lambda _N|$$
(15)
This is different from the case $`RR_c^{(mf)}`$ where at $`G>\frac{e^2}{h}`$ the critical concentration is exponentially small. The difference originates from the difference in correlation times.
Neglecting the second term in Eq.7 and assuming for simplicity that superconducting grains form a square lattice we get an expression for the correlation functions
$`<\alpha _k(\omega )\alpha _l^{}(\omega )>={\displaystyle \frac{1}{\nu \tau _0}}(i|\omega |+{\displaystyle \frac{1}{\tau }}+{\displaystyle \frac{1}{\nu \tau _0}}J_{ij})_{k,l}^1`$ (16)
$`={\displaystyle \frac{1}{\nu \tau _0}}{\displaystyle \underset{𝒒}{}}\mathrm{exp}(i𝒒𝒓){\displaystyle \frac{1}{i|\omega |+\frac{1}{\tau }K(𝒒)}}`$ (17)
where components of $`𝒒`$ are whole multiples of $`\frac{2\pi }{L}`$ and $`L`$ is the samle size.
$$K(q)=\frac{1}{\nu \tau _0}\underset{i}{}exp(i𝒒𝒓)K(𝒓_i)$$
(18)
At $`|𝒒|N^{1/2}1`$ we have
$$K(𝒒)=\pi DNln((1+(qL_0)^2)L_0^2N)$$
(19)
The mean field Eqs.16-18 are valid because the characteristic radius of interaction between grains $`N^{1/2}ln(L_0N^{\frac{1}{2}})N^{1/2}`$ is much larger than the average intergrain distance.
Qualitatively the picture of quantum superconducting fluctuations in the normal metal is similar to superconducting fluctuations in uniform metals at temperatures which are close to the critical one :
a) Due to quantum fluctuations conductivity of the system is enhanced as compared to the normal metal value. It exhibits a big positive magnetoresistance.
b) At small magnetic field the zero temperature Hall coefficient is suppressed as compared to its normal metal value . It also exhibits strong magnetic field dependence.
c) Diamagnetic susceptibility is enhanced and exhibits a strong nonlinearity as a function of magnetic field.
d) Energy dependence of the density of states at the Fermi surface has a dip, whose amplitude is magnetic field dependent.
## IV Conclusion
We have considered a model in which superconducting grains with the radius $`R<R_c`$ are embedded in the normal metal and have shown that in this case there is a zero temperature quantum superconductor-normal metal phase transition as a function of $`N`$ and $`R`$. This transition is driven, primarily, by fluctuations of the modulus of the superconducting order parameter. Though the parameters of the transition, in principle, depend on $`D`$, it exists even in the case when there is no disorder in the sample. The critical concentration of grains turns out to be relatively large.
Calculations presented above did not take into account localization effects . In d=3 case they are small. The question of whether in two-dimensional case the metallic state is localized requires additional investigation. On the other hand, as we have discussed, the transition can take place at very large sample conductances when the localization length is also very large.
We would like to mention that in the presence of the electron-electron repulsion in the metal the metal-insulator transition can exist even in the mean field approximation described by the Eq.2. For example, in the case $`\lambda _s|\lambda _N|`$ and $`R<R_c^{mf}`$ we have $`N_c\frac{1}{R^2}`$.
Finally, we would like to mention, that in our opinion experimental data on 2-D films do not contradict the possibility of a zero temperature superconductor-normal metal quantum phase transition.
## V Appendix
We start with a standard expression for the partition function in a superconductors (see for example ):
$$Z=D\mathrm{\Delta }D\mathrm{\Delta }^{}\mathrm{exp}(iS_{eff})$$
(20)
$$S_{eff}=𝑑𝒓𝑑t[\frac{|\mathrm{\Delta }|^2}{\lambda (𝒓)}+_0^1𝑑\eta Tr[\widehat{\mathrm{\Delta }}(𝒓,t)\widehat{G}_\eta (𝒓t,𝒓t)]]$$
(21)
where $`\widehat{G}_\eta (𝒓t,𝒓^{}t^{})`$ is a matrix Green function which is a solution of the Gorkov equation
$$[i\frac{d}{dt}\xi (i)\sigma _z+u(𝒓)\eta \widehat{\mathrm{\Delta }}(𝒓,t)]\widehat{G}_\eta (𝒓,t;𝒓^{}t^{})=\delta (𝒓𝒓^{})\delta (tt^{})\widehat{I}$$
(22)
Here $`u(𝒓)`$ is the external potential.
$$\widehat{\mathrm{\Delta }}=\left(\begin{array}{cc}0& \mathrm{\Delta }\\ \mathrm{\Delta }^{}& 0\end{array}\right)$$
(23)
and $`\widehat{I}`$ is a $`22`$ unit matrix in Nambu space. We assume that $`u(𝒓)`$ has a white noise statistics with correlation functions $`<<u(𝒓)>>`$ and $`<<u(𝒓)u(𝒓^{})>>=\frac{v_F}{l\nu }\delta (𝒓𝒓^{})`$. Here brackets $`<<>>`$ stand for averaging over realizations of the scattering potential.
Averaging Eq.21 over realizations of $`u(𝒓)`$, neglecting all weak localization and mesoscopic corrections, and making the diffusion approximation we get Uzadel equations for the normal and the anomalous Green’s functions (See, for example, )
$$\omega F(\omega ,𝒓)+\frac{1}{2}D(F(\omega ,𝒓)^2G(\omega ,𝒓)G(\omega ,𝒓)^2F(\omega ,𝒓))=\eta \mathrm{\Delta }^{}(𝒓)G(\omega ,𝒓)$$
(24)
$$G^2+|F|^2=1$$
(25)
Expanding Eqs.20-24 with respect to $`\mathrm{\Delta }`$ we get an expression for the effective action Eq.3 for an individual grain.
To get Eqs.9,10 we have to solve the Uzadel equation Eq.24 in the normal metal between the superconducting grains. This solution corresponds to an electron diffusion in the metal and Andreev reflections from the superconducting grains. In the case $`RR_c^{mf}`$ we have $`|\mathrm{\Delta }(t)|=\mathrm{\Delta }_0`$ and the Andreev scattering cross section equals to $`R^2`$. Since in our case $`\mathrm{\Delta }(t)\mathrm{\Delta }_0`$ the cross section is of order of $`R^2\frac{\mathrm{\Delta }^2(t)}{\mathrm{\Delta }_0^2}`$ and we get Eq.9,10.
This work was supported by Division of Material Sciences, U.S.National Science Foundation under Contract No. DMR-9205144. We would like to thank A.Andreev, S.Chakravarty, M.Feigelman, A.I.Larkin, D.Khmelnitskii, S.Kivelson, and F.Zhou for useful discussions.
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# Far-ultraviolet Spectroscopy of Venus and Mars at 4 Å Resolution with the Hopkins Ultraviolet Telescope on Astro-2
## 1 INTRODUCTION
Ultraviolet observations of the atmospheres of Venus and Mars, primarily from fly-by or orbiting spacecraft, but also from platforms above the terrestrial atmosphere, have played an important role in the study of the composition and structure of the CO<sub>2</sub> atmospheres of these planets. Early spacecraft experiments were broad-band photometers that provided information about the spatial distribution of the emissions in the passband but required remote spectroscopic measurements to assure the identity of the emitting species. A thorough review of this subject is given by Paxton & Anderson (1992). Once the emitting species were known, narrowband polychromators were developed and flown on missions such as Mariner 10 (Broadfoot et al., 1974) and Venera 11 and 12 (Bertaux et al., 1981), although the interpretation of the data from these instruments was far from unambiguous. With Mariners 6, 7 and 9, flown to Mars in the early 1970s, and Pioneer Venus Orbiter launched to Venus in 1978, ultraviolet spectrometers were included in the payloads and provided modest spectral resolution. Very little additional spectroscopy, particularly from newer generations of Earth-orbiting observatories such as the International Ultraviolet Explorer (IUE) and the Hubble Space Telescope (HST), has been done and is also summarized by Paxton & Anderson. In February 1990, Galileo flew by Venus on its way to Jupiter and the ultraviolet spectrometers on board made observations that were reported by Hord et al. (1991).
The flight of the Hopkins Ultraviolet Telescope (HUT) on the Astro-2 mission on the Space Shuttle Endeavour in March 1995 provided an opportunity to measure the ultraviolet disk spectra of both Venus and Mars at a spectral resolution ($`4`$ Å) significantly higher than any of the prior spacecraft observations. Moreover, since the first order spectral range of HUT extended to wavelengths as short as 830 Å, the observations of Venus enabled the resolution of the identity of the emissions recorded in the narrow-band photometric channels of the Venera 11 and 12 EUV instruments that had been interpreted in terms of analogue terrestrial spectra (Bertaux et al., 1981). This paper presents the HUT disk spectra of Venus and Mars together with the spectral identifications and the disk-averaged brightnesses. CO fluorescence in both the $`BX`$ (0,0) and $`CX`$ (0,0) Birge-Hopfield bands is identified in both Venus and Mars and shown to be consistent with current models of the CO/CO<sub>2</sub> mixing ratios in the atmospheres of these planets. Several other emissions are definitively identified for the first time.
## 2 INSTRUMENT and OBSERVATIONS
The HUT instrument consists of a 0.9-m SiC coated mirror that feeds a prime focus spectrograph with a photon-counting microchannel plate detector and photodiode array readout. Details of the instrument and its performance and calibration are given by Davidsen et al. (1992) and Kruk et al. (1995, 1999). The spectrograph covered 820 – 1840 Å in first order with a dispersion of 0.51 Å per pixel. The absolute calibration, monitored several times during the mission by observations of pure hydrogen white dwarfs, is considered accurate to better than 5% at all wavelengths longward of 912 Å (Kruk et al., 1999). For the observations reported here, the $`20^{\prime \prime }`$ diameter spectrograph aperture was used, and was underfilled by both Venus and Mars. The mean spectral resolution was determined by the diameter of the planet, 15$`\stackrel{}{\mathrm{.}}`$1 for Venus and 12$`\stackrel{}{\mathrm{.}}`$2 for Mars, and the illuminated fraction, giving a resolution of 4–4.5 Å over the entire spectral band.
At the time of the Astro-2 mission, Venus was at a solar elongation slightly inside the 45 solar avoidance constraint of HUT. Nevertheless, since the performance of HUT and the supporting Instrument Pointing System was nominal after ten days of the mission, it was decided to attempt to observe Venus at a solar elongation of 40$`\stackrel{}{\mathrm{.}}`$2. This observation occurred at UT 05:25 on 13 March 1995 and produced no adverse heating of the telescope assembly nor appreciably elevated scattered light. A second observation was made about 9 hours later. Both observations were made during orbit day.
A spectrum of Venus from the first observation is shown in Figure Far-ultraviolet Spectroscopy of Venus and Mars at 4 Å Resolution with the Hopkins Ultraviolet Telescope on Astro-2. It includes not only the desired disk spectrum of Venus but also contributions from the terrestrial day airglow. The Earth’s dayglow includes H i Lyman-$`\alpha `$, $`\beta `$ and $`\gamma `$, as well as emissions of atomic and ionic oxygen and atomic nitrogen that are also present in the dayglow of Venus. Fortuitously, it took longer to acquire and lock onto Venus during the second observation so there are $``$500 seconds of background data taken during this observation that correspond to approximately the same viewing geometry (with respect to the Earth’s atmosphere) as when Venus was in the aperture during the first observation. This is illustrated in Figure Far-ultraviolet Spectroscopy of Venus and Mars at 4 Å Resolution with the Hopkins Ultraviolet Telescope on Astro-2, which shows the count rates for two of the brighter emissions, O i $`\lambda `$1304 and N i $`\lambda `$1200, as a function of time for both observations. Figure Far-ultraviolet Spectroscopy of Venus and Mars at 4 Å Resolution with the Hopkins Ultraviolet Telescope on Astro-2 shows two 516 second integrations, Venus + Earth and Earth alone. The difference spectrum is then the true spectrum of Venus, and this spectrum is used to derive the disk brightnesses for the O i, N i, and H i emissions listed in Table 1. For the other emissions, the longer exposure spectrum of Figure Far-ultraviolet Spectroscopy of Venus and Mars at 4 Å Resolution with the Hopkins Ultraviolet Telescope on Astro-2 is used. The error bars given in Table 1 are 1-$`\sigma `$ statistical uncertainties in the observed counts within a given emission feature.
Mars was observed during orbit night beginning at UT 22:33 on 12 March 1995. The HUT spectrum of Mars is shown in Figure Far-ultraviolet Spectroscopy of Venus and Mars at 4 Å Resolution with the Hopkins Ultraviolet Telescope on Astro-2, and the derived disk brightnesses are also listed in Table 1.
## 3 DISCUSSION
The HUT observations of Venus and Mars represent the highest quality spectra, in terms of both spectral resolution and instrument sensitivity, obtained to date for the ultraviolet below 1800 Å. The present discussion will focus on the identification of spectral features and comparison with previous observations. Future work will concern the modelling of these spectra in terms of current atmospheric models based on in situ measurements made over the past two decades (Paxton & Anderson, 1992).
### 3.1 Carbon Monoxide and Atomic Carbon
Below 2000 Å, the ultraviolet dayglow of Venus and Mars is dominated by emissions of carbon monoxide and carbon (Durrance, 1981; Fox, 1992). From both Fig. Far-ultraviolet Spectroscopy of Venus and Mars at 4 Å Resolution with the Hopkins Ultraviolet Telescope on Astro-2 and Fig. Far-ultraviolet Spectroscopy of Venus and Mars at 4 Å Resolution with the Hopkins Ultraviolet Telescope on Astro-2 we see that a large number of individual bands of the CO Fourth Positive ($`A^1\mathrm{\Pi }`$ – $`X^1\mathrm{\Sigma }^+`$) system are clearly identified. When compared with calculated optically thin fluorescence efficiencies (“g-factors”) for this system (Tozzi et al., 1998), the effect of the twenty-fold increase in the mean CO<sub>2</sub> absorption cross-section from 1700 to 1500 Å (Yoshino et al., 1996) is clearly seen in the enhancement of the weaker longer wavelength bands. In addition, the CO is optically thick in that the expected strong ($`v^{},0`$) bands are particularly weak implying that photons are being pumped out of these bands into other members of the ($`v^{},v^{\prime \prime }`$) progression by repeated absorption and re-emission. Bands of the ($`14,v^{\prime \prime }`$) and ($`9,v^{\prime \prime }`$) progressions, pumped by solar H i Lyman-$`\alpha `$ and O i $`\lambda `$1304, respectively (Kassal, 1976; Durrance, 1981; Wolven & Feldman, 1998), are clearly seen, particularly the (14,3) band at 1316 Å and the (9,2) band at 1378 Å, both of which are free of blending by other CO bands or atomic emissions. The strong emission feature at $`1355`$ Å is a blend of O i $`\lambda `$1356 and the CO (14,4) band at 1352 Å, the strongest in the solar Lyman-$`\alpha `$ pumped progression, with most of the observed emission due to CO.
Other observed features include the CO Hopfield-Birge bands, $`B^1\mathrm{\Sigma }^+`$ – $`X^1\mathrm{\Sigma }^+`$ (0,0) at 1151 Å and $`C^1\mathrm{\Sigma }^+`$ – $`X^1\mathrm{\Sigma }^+`$ (0,0) at 1088 Å, which are identified for the first time in the spectra of Venus and Mars. They appear to be present in the Galileo EUV spectrum of Venus (Hord et al., 1991), but cannot be positively identified at the 30 Å resolution of that instrument. The $`BX`$ band is blended with emission from O i $`\lambda `$1152. At first sight, the relative intensities of these bands and the similarity with low energy electron impact laboratory spectra (Kanik et al., 1995), suggest that excitation of CO by photoelectrons is a major source of the emission. However, as is shown in Section 3.2, it is possible to account for the observed brightness of the $`CX`$ (0,0) band solely on the basis of resonance fluorescence of solar ultraviolet radiation. No other bands of either of these systems is detected, as expected from the strongly diagonal nature of both band systems and the high pre-dissociation fractions for $`v^{}1`$ (Eidelsberg et al., 1991). There is also a weak indication of the $`E^1\mathrm{\Pi }`$ – $`X^1\mathrm{\Sigma }^+`$ (0,0) band at 1076 Å in the spectrum of Venus.
In addition to the principal C i multiplets at 1561 and 1657 Å, a number of other atomic carbon lines are identified in Fig. Far-ultraviolet Spectroscopy of Venus and Mars at 4 Å Resolution with the Hopkins Ultraviolet Telescope on Astro-2. Some of these, the multiplets at 1329, 1280, 1277 and 1261 Å, were identified in the Mariner 6 and 7 spectra of the upper atmosphere of Mars (Barth et al., 1971). In addition, we identify emission from C ii $`\lambda `$1335, blended with the (9,1) CO Fourth Positive band at 1339 Å, indicating the presence of C<sup>+</sup> in the ionospheres of Venus and Mars.
It is interesting to compare the ratio of these emissions between Venus and Mars. From Table 1 we find that for all of the CO bands and the atomic carbon emissions this ratio is $`15`$, whereas the ratio of solar flux is 5.3. This had been previously noted by Moos (1974), and reflects the lower CO to CO<sub>2</sub> mixing ratio on Mars compared to that on Venus.
### 3.2 The CO Hopfield-Birge Bands
Emissions from CO shortward of Lyman-$`\alpha `$ appear in the form of the Hopfield-Birge bands $`C^1\mathrm{\Sigma }^+`$ – $`X^1\mathrm{\Sigma }^+`$ (0,0) at 1088 Å and $`B^1\mathrm{\Sigma }^+`$ – $`X^1\mathrm{\Sigma }^+`$ (0,0) at 1151 Å. The $`BX`$ band is blended with emission from O I $`\lambda `$1152. The intensities of these bands are not consistent with optically thin emission. Since no other bands of these systems were detected a scenario of resonant scattering was assumed. The band is scattered from a column of CO defined by the optical thickness due to pure absorption by CO<sub>2</sub>:
$$𝒩_{CO}=_0^{\mathrm{}}N_{CO}(z)e^{\tau _{CO_2}(z)}𝑑z$$
and computed using the number densities of the Venus International Reference Atmosphere (VIRA) model (Keating et al., 1985) and the model of Fox & Dalgarno (1979) for Mars, with CO<sub>2</sub> absorption coefficients taken from Nakata et al. (1965).
To test the applicability of these atmospheric profiles to modeling the CO emission, we compared the column density derived from the (14,3) band at 1316 Å of the Fourth Positive system ($`A^1\mathrm{\Pi }`$ – $`X^1\mathrm{\Sigma }^+`$) and that derived using the above relationship. The brightness of the (14,3) band (in rayleighs) in the optically thin limit is
$$B=g(14,3)𝒩_{CO}\times 10^6$$
and measured to be $`146\pm 25`$ R in the Venus spectrum. The g-factor for Lyman-$`\alpha `$ absorption in the (14,0) band was evaluated using the oscillator strength recommended by Eidelsberg et al. (1999), $`1.80\times 10^5`$, based on their recent work and the work of Jolly et al. (1997) and Stark et al. (1998), together with the solar Lyman-$`\alpha `$ line profile from Lemaire et al. (1998) normalized to UARS/SOLSTICE measurements of the solar Lyman-$`\alpha `$ flux (Woods et al., 1996) made at the time of the Astro-2 mission. Branching ratios were taken from Kurucz (1976). We thus derive a CO column density of $`3.4\times 10^{16}`$ cm<sup>-2</sup>, whereas the VIRA model predicts a column of $`5.2\times 10^{16}`$ cm<sup>-2</sup>. The close agreement, considering the uncertainty in all of the parameters involved in the g-factor calculation, supports our use of these model atmospheres.
Following the treatment of Liu & Dalgarno (1996), a theoretical curve-of-growth was calculated using the rates of emission from the excited electronic state of each rovibrational transition to determine the expected omnidirectional brightness, $`4\pi `$, in rayleighs. The equivalent width ($`EQW`$) of the band is then defined as $`EQW=4\pi /\pi _{}`$, where $`\pi _{}`$ is the incident solar flux determined from HUT observations of the moon and the lunar albedo as given by Henry et al. (1995). To compute the curve-of-growth, we used a density weighted temperature average to characterize both the kinetic and rotational temperatures. Including transitions through rotational quantum number $`J=30`$, we produced curves-of-growth for both Venusian and Martian $`CX`$ (0,0) emissions as shown in Figure Far-ultraviolet Spectroscopy of Venus and Mars at 4 Å Resolution with the Hopkins Ultraviolet Telescope on Astro-2. The curves appear as bands with the upper edge corresponding to an oscillator strength of $`1.177\times 10^1`$ (Chan et al., 1993) as recommended by Morton & Noreau (1994), and the lower to $`6.19\times 10^2`$ (Eidelsberg et al., 1991).
The emissions appear to be consistent with resonant scattering from the upper atmospheres of both planets. As with the $`AX`$ bands, the observed ratio of the emissions between Venus and Mars is $`15`$, whereas the ratio of the solar flux is 5.3. This difference reflects the lower CO mixing ratio on Mars as well as a reduced thermospheric temperature at $`\tau _{CO_2}=1`$. Curves of growth for $`BX`$ (0,0) ($`f=6.5\times 10^3`$, Stark et al. 1999) were determined in the same manner as for $`CX`$ (0,0) and the equivalent width corresponding to the appropriate column density was used to calculate an expected $`28\pm 10`$ and $`1.8\pm 0.8`$ R of emission from Venus and Mars respectively. Subtracting these from the total observed emission results in O I $`\lambda `$1152 emission of about $`100\pm 14`$ R for Venus and $`4.5\pm 2.0`$ R for Mars, most likely due to direct electron impact excitation of atomic oxygen.
### 3.3 Argon and Helium
Argon emission at 1048 and 1066 Å is detected only in the spectrum of Mars (Fig. Far-ultraviolet Spectroscopy of Venus and Mars at 4 Å Resolution with the Hopkins Ultraviolet Telescope on Astro-2). The upper limits for Venus are consistent with the known argon mixing ratio of $`<100`$ ppm. Note that Ar i $`\lambda `$1066 appears brighter than Ar i $`\lambda `$1048, even though the fluorescence efficiency of the 1048 Å line is 3.5 times higher than that of the 1066 Å line.<sup>1</sup><sup>1</sup>1The Ar i g-factors are evaluated using solar fluxes inferred from the HUT observations of sunlight reflected from the Moon (Henry et al., 1995) This is the result of the higher CO<sub>2</sub> absorption cross-section at 1048 Å. Using the approach of Bertaux et al. (1981), we find an argon mixing ratio of 0.016 and 0.021, from the brightness of the 1048 and 1066 Å lines, respectively. Considering the uncertainties in both the data and the fluorescence efficiencies, this result is in good agreement with the known value of 0.016 for Mars (Barth, 1985). The Ar i multiplets near 867 Å are not detected on either Mars or Venus.
Figure Far-ultraviolet Spectroscopy of Venus and Mars at 4 Å Resolution with the Hopkins Ultraviolet Telescope on Astro-2 also shows the presence of He i $`\lambda `$584 in second order at 1168 Å in the spectrum of Venus. The second order effective area of HUT was only $`1`$ cm<sup>2</sup>, compared to the peak first order effective area of 25 cm<sup>2</sup> near 1150 Å (Kruk et al., 1999), so that the derived brightness has a large statistical uncertainty, $`300\pm 200`$ R. Nevertheless, this result is consistent with prior measurements of He i emission from Venera and Galileo.
### 3.4 Other Emissions
Shortward of Lyman-$`\alpha `$, other emission features detected include O ii $`\lambda `$834, O i $`\lambda `$989, H i Lyman-$`\beta `$, and N i $`\lambda `$$`\lambda `$1134 and 1200. The atomic nitrogen features are identified for the first time, although, as above, N i $`\lambda `$1134 appears to be present in the Galileo EUV spectrum of Venus. The disk brightnesses for both Venus and Mars are listed in Table 1. While there are many weak features present in the spectrum of Venus between 840 and 1000 Å, most of these are of terrestrial origin, primarily of atomic and singly ionized nitrogen, and the contributions from Venus are at least an order of magnitude lower than those reported by Stern et al. (1996) from a sounding rocket observation.
### 3.5 Comparison with Venera 11/12 and Galileo
Table 2 gives a comparison of the HUT measurements of the disk of Venus with those reported from Venera 11 and 12 (Bertaux et al., 1981) and, more recently, Galileo (Hord et al., 1991). The Galileo flyby of Venus was in February 1990, during solar maximum, while the HUT observations in March 1995 were at a time approaching solar minimum. This is reflected in the roughly factor of two difference between the two sets of data, which must be considered excellent agreement. The Venera 11/12 measurements were made in December 1978, also approaching solar maximum. Thus, a similar argument accounts for the difference in O i and O ii brightnesses between HUT and Venera. However, it appears that the Venera channels at wavelengths longward of 1500 Å were severely contaminated by scattered light.
## 4 CONCLUSION
We have obtained far-ultraviolet spectra of Venus and Mars in the range 820 – 1840 Å at $``$4 Å resolution with the Hopkins Ultraviolet Telescope (HUT) during the Astro-2 mission in March 1995. The spectra of both planets are rich in CO band emission, some of the systems being identified for the first time, together with strong O i and C i multiplets. Resonance fluorescence is identified as the dominant source of the CO emission. Atomic nitrogen emissions are also identified in the spectrum of Venus, while the Ar i doublet is seen only in the spectrum of Mars. These spectra, obtained at higher spectral resolution than was possible from earlier fly-by and orbiting missions to these planets, elucidates and extends the earlier spectroscopic measurements and should provide guidance in the interpretation of the far-ultraviolet spectra obtained during the recent Cassini fly-by of Venus (Stewart et al., 1999).
It is a pleasure to thank the Spacelab Operations Support Group at Marshall Space Flight Center for their support during the Astro-2 mission. We also thank our many colleagues at the Johns Hopkins University and the Applied Physics Laboratory for their contributions to the success of the Hopkins Ultraviolet Telescope. This work was supported by NASA grant NAG5-5122 and contract NAS5-27000 to the Johns Hopkins University.
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# 1 Introduction
## 1 Introduction
Observational constraints are of fundamental importance to build a realistic chemical evolution model. With respect to these constraints the last years have been of crucial importance (Pagel 1997) and, in the case of the Milky Way, the new observational data required a revision of the previous chemical evolution models (see Pagel and Tautvaisiene 1995 and Chiappini et al. 1997, hereinafter CMG, for a discussion on this point).
In particular, Gratton et al. (1996, 1999) showed that the distribution of the abundances of $`\alpha `$-elements to Fe for a large homogeneous sample of stars in the solar neighbourhood seems to indicate a short timescale for the evolution of the halo and thick disk phases and a sudden decrease in the star formation in the epoch preceding the formation of the thin disk. An analogous result was found by Fuhrmann (1998) for the \[Mg/Fe\] ratio. Moreover, Beers and Sommer-Larsen (1995) have shown that the thick disk population extends to very low metallicities. Those are very important new information which stimulated us (CMG) to consider a different picture of Galaxy formation than those previously adopted.
Previous models, in fact, (e.g. Matteucci and François 1989) were based on a pure colapse picture where the disk formed from gas shed from the halo. These models, however, are difficult to conciliate with the new results discussed above indicating a discontinuity between halo and thin disk. We than suggested the so-called Two-Infall Model, a model that assumes that the Galactic thin disk was not only formed from gas shed from the halo but was formed mainly from extra-galactic gas. We assume two main infall episodes, the first one is responsible for the formation of the population made of that fraction of the halo and thick disk stars which originated from a fast dissipative collapse, such as suggested by Eggen et al. (1962). The second infall episode forms the thin disk component with a timescale much longer than that of the halo formation.
In this new picture the disk was formed slowly (with a timescale of 7-8 Gyrs at the solar vicinity) and from inside-out (Matteucci and François 1989). A direct consequence of that is that at high redshift we should expect to see smaller disks in size (Roche et al. 1998). This long timescale for the formation of the thin disk at the solar vicinity, required to produce a good fit of the obseved G-dwarf metallicity distribution (Rocha-Pinto and Maciel 1996) was also confirmed by recent chemical evolution models (eg. Portinari et al. 1998, Prantzos and Silk 1998, Chang et al. 1999) as well as by chemo-dynamical models (eg. Hensler 1998) and is also in agreement with the results showed in this conference by Carraro (2000).
The two-infall model adopted a constant IMF. On the empirical grounds there is at present no clear direct evidence that the IMF in the Galaxy has varied with time. A detailed discussion about possible observed variations in the IMF in different environments is given by Scalo (1998), but such variations are comparable with the uncertainties still involved in the IMF determinations. The present uncertainties in the observational results prevent any conclusion against the idea of an universal IMF.
However, a variable IMF, which formed relatively more massive stars during the earlier phases of the evolution of the Galaxy compared to the one observed today in the solar vicinity, has often been suggested as being one of the possible solutions for the G-dwarf problem (namely the deficiency of metal-poor stars in the solar neighborhood when compared with the number of such stars predicted by the simple model). Such an IMF would also be physically plausible from the theoretical point of view if the IMF depends on a mass scale such as the Jeans mass. Given the uncertainties in both theoretical and observational grounds, the proposed IMFs can in principle be tested by means of a detailed chemical evolution model.
An example of a theoretical approach to the IMF problem is the one proposed by Padoan et al. (1997 - hereinafter PNJ). Since random motions are probably ubiquitous in sites of star formation, PNJ suggested to describe the formation of protostars as the gravitational collapse of Jeans masses in a density distribution shaped by random supersonic motions (but see Scalo et al. 1998). This IMF was already tested in models of elliptical galaxies by Chiosi et al. (1998) and they concluded that a strongly varying IMF could be suitable for such galaxies.
Our goal was to address the question of what time-dependent IMF properties are allowed in order to still match the observational constraints when adopting the two-infall model (CMG). To do this we tested different IMFs in our chemical evolution code, going from those where the variability is contained in the slope of the power-law, assumed to be a function of the metallicity (parametrizations adopted by Scully et al. 1996, Matteucci and Tornambé 1987), to the one by PNJ which predicts also a change in the stellar mass which contributes most to the IMF as a function of time.
A detailed description of the Two-Infall model as well as of the PNJ IMF and the hypothesis needed to introduce it in our chemical evolution code are discussed in detail in Chiappini et al. 2000a (hereinafter CMP; see also CMG and PNJ).
## 2 Results
### 2.1 The Solar Vicinity
The two-infall chemical evolution model, published by CMG (hereinafter Model A) has been modified to test the IMF proposed by PNJ (hereinafter Model B and C) and the other two IMFs (Models D - Matteucci and Tornambé 1987 and E - Scully et al. 1996). We address the reader to the CMP paper where a detailed description of each one of the IMFs adopted here can be found.
The predictions of models A and B turn out to be very similar concerning the observed properties of the solar vicinity (tables 2 and 3 of CMP). This is due to the fact that the PNJ IMF, when applied to our Galaxy, does not vary much over most of the Galactic lifetime which is a consequence of our assumptions of constant molecular cloud temperature during the Galactic evolution. In fact, if the cloud temperature and density have varied strongly over the history of the Galaxy, the predicted IMF will also vary more and this would certainly worsen the agreement with the observational constraints considered here.
On the other hand, the other two IMFs (models D and E) vary substantially over the entire Galactic evolution. As a consequence of that, these IMFs do not give a good agreement with the solar vicinity observational constraints. Figure 1 shows the predicted G-dwarf metallicity distribution for models D and E compared with models A and B. Those models predict too few metal-poor stars and a higher solar metallicity peak than the observed one. The predicted solar abundances are also not in agreement with the observed ones (Table 3 of CMP).
It is worth noting that here we adopted a long timescale for the formation of the solar vicinity (8 Gyrs) and that a better agreement with the G-dwarf observed metallicity distribution could in principle be achieved by adopting these variable IMFs in a closed box model scenario. However, as recently showed by Martinelli and Matteucci (1999), models which can fit the G-dwarf metallicity distribution with a more variable IMF and which adopt a shorter infall timescale for the solar vicinity formation do not give good agreement with the other observed properties. Moreover, in a close box model the predicted gradients along the disk would be essentially flat.
From the comparison with the solar vicinity observational constraints we can confirm the result by Matteucci & Tornambé (1985) that only a constant IMF or an IMF that varied only at early times can be in agreement with the solar vicinity properties.
### 2.2 The Galactic Disk
The IMF in model B combined with an inside-out cenario for the disk formation predicts a higher number of low mass-stars towards the galactic center, where the metallicity is higher (see Figure 8 of CMP). A direct consequence of this fact is the predicted flatter gradient with respect to the one of Model A. Figure 2a shows the oxygen abundance gradient as predicted by model A (CMG), model B (adopting the PNJ IMF) and model C (same as B but with star formation efficiency increasing with decreasing galactocentric distance). Model B predicts a flat gradient between 6 and 10 kpc, and a small negative one in the inner part of the Galaxy ($`R<6kpc`$). Model C, which adopts an increasing star formation efficiency $`\nu `$ towards the galactic center, predicts a steeper gradient inside the solar circle. In this case a the bimodal abundance gradient, steeper in the inner part and flatter outside, similar to the one found by CMG, is recovered. However, as it can be seen in Figure 2b, a model with a higher star formation efficiency in the central parts (Model C) consumes the gas very fast thus reaching the threshold density value very soon. As a consequence, model C predicts a flat gas density radial profile at variance with observations, whereas Model B predicts a gas density distribution very similar to the one predicted by model A and in better agreement with observations.
Figure 2b also shows models D and E predictions for the radial gas profile. From this figure it can be seen that the only two models in good agreement with the observed radial profile are models A and B. However, as shown before, model B predicts even flatter abundance gradients than model A. Models D and E predict a flatter radial gas profile than the observed one but in better agreement with the observations than the one predicted by model C. Moreover, the oxygen abundance gradient predicted by models D and E (see figure 11 of CMP) has a slope which is similar to the one obtained with model A. In any case, those two models do not give predictions in agreeement with the solar vicinity constraints as already mentioned in the previous section.
This is a very important result showing that only model A, with a constant IMF, is in agreement simultaneously with the solar vicinity and the disk observational constraints. The abundance gradients predicted by model A, although slightly flatter than the observed ones, are steeper in the inner parts of the Galactic disk and also steepen with time, in agreement with the recent results by Maciel & Quireza (1999) (but see Chiappini et al. 2000b for a detailed discussion on abundance gradients).
## 3 Conclusions
The main results are summarized below:
a) We tested the IMF proposed by PNJ in a model for the chemical evolution of the Milky Way (CMG) and we showed that this IMF gives good agreement with the observed properties of the solar vicinity. However, such an agreement is due to the fact that this IMF when applied to the two-infall model shows a time variation that is important only in the early phases of Galactic evolution. This in turn is due to the simplifying assumptions adopted here, like neglecting the dependence of the IMF on the gas temperature which would produce more sharply varying IMF. In these early phases the IMF is biased towards massive stars.
c) the PNJ IMF combined with the inside-out picture for the thin disk formation predicts a gradient flatter than the one predicted by a model which adopts a constant IMF. This situation cannot be reversed by changing the SFR because in this case the abundance gradient is recovered but the gas density profile is destroyed.
d) Models which adopt IMFs strongly dependent on metallicity, thus simulating a dependence also on the gas temperature, are not in agreement with the most important observational constraints of the solar vicinity and predict radial gas profiles at variance with observations, therefore they should be rejected.
e) We conclude that a constant IMF and the assumption of a continuous infall onto the Galactic disk is still the best way to explain the observational constraints in the Milky Way including the G-dwarf metallicity distribution. A probable source of the required infall could be the HVCs (see Burton and Braun 1999 for a discussion on this point).
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# 1. Introduction
## 1. Introduction
In spite of the fact that the top quark has been discovered already several years ago its interactions are still very weakly constrained. It remains an open question if top-quark couplings obey the Standard Model (SM) scheme of the electroweak forces or there exists a contribution from physics beyond the SM. We could interpret the great success of the 1-loop precision tests of the SM as a strong indication that the third generation also obeys the SM scheme. However, an independent and direct measurement of the top-quark couplings is definitely necessary before drawing any definite conclusion concerning non-standard physics.
Over the past several years there was a substantial effort devoted to a possibility of determining top-quark couplings through measurements performed at the open top region<sup>♯1</sup><sup>♯1</sup>♯1Recently an interesting and complementary analysis by Jezabek, Nagano and Sumino has been published where the authors discussed possibility of determining $`CP`$-violating production form factors at the $`t\overline{t}`$ threshold region. of future $`e^+e^{}`$ linear colliders . The existing studies focused mainly on tests of $`CP`$ violation in top-quark interactions. In this article we will construct some new tools which could help to measure both $`CP`$ violating and $`CP`$ conserving top-quark couplings at linear colliders and therefore reveal the structure of fundamental interactions beyond the SM.
The top quark decays immediately after being produced and its huge mass $`m_t174`$ GeV leads to a decay width $`\mathrm{\Gamma }_t`$ much larger than $`\mathrm{\Lambda }_{\mathrm{QCD}}`$. Therefore the decay process is not influenced by any fragmentation effects and decay products will provide useful information on top-quark properties. Here we will consider distributions of either $`\mathrm{}^\pm `$ in the inclusive process $`e^+e^{}t\overline{t}\mathrm{}^\pm \mathrm{}`$ or bottom quarks from $`e^+e^{}t\overline{t}\stackrel{()}{b}\mathrm{}`$. It turns out that the analysis of the leptonic and $`b`$-quark final states is similar and could be presented simultaneously. Although $`t\overline{t}`$ are also produced via $`WW`$ fusion , we do not consider here this mechanism since $`\sigma (e^+e^{}t\overline{t}\nu \overline{\nu })`$ is expected to be much smaller than $`\sigma (e^+e^{}t\overline{t})`$ for the energy of our interest ($`\sqrt{s}\stackrel{<}{}2\text{TeV}`$) .
This paper is organized as follows. First in sec.2 we describe the basic framework of our analysis, and then show the angular and energy distributions of the lepton and $`b`$-quark in sec.3. In sec.4, after briefly reviewing the optimal-observable procedure , we estimate to what precision all the non-standard parameters can be measured or constrained adjusting the initial beam polarizations. Finally, we summarize our results in sec.5. In the appendix we collect several functions used in the main text for completeness, though some of them could also be found in our previous papers .
## 2. Framework and Formalism
We parameterize $`t\overline{t}`$ couplings to the photon and the $`Z`$ boson in the following way
$$\mathrm{\Gamma }_{vt\overline{t}}^\mu =\frac{g}{2}\overline{u}(p_t)\left[\gamma ^\mu \{A_v+\delta A_v(B_v+\delta B_v)\gamma _5\}+\frac{(p_tp_{\overline{t}})^\mu }{2m_t}(\delta C_v\delta D_v\gamma _5)\right]v(p_{\overline{t}}),$$
(2.1)
where $`g`$ denotes the $`SU(2)`$ gauge coupling constant, $`v=\gamma ,Z`$, and
$$A_\gamma =\frac{4}{3}\mathrm{sin}\theta _W,B_\gamma =0,A_Z=\frac{1}{2\mathrm{cos}\theta _W}\left(1\frac{8}{3}\mathrm{sin}^2\theta _W\right),B_Z=\frac{1}{2\mathrm{cos}\theta _W}$$
denote the SM contributions to the vertices. Among the above non-SM form factors, $`\delta A_v`$, $`\delta B_v`$, $`\delta C_v`$ describe $`CP`$-conserving while $`\delta D_v`$ parameterizes $`CP`$-violating interactions. Similarly, we adopt the following parameterization of the $`Wtb`$ vertex suitable for the $`t`$ and $`\overline{t}`$ decays:
$`\mathrm{\Gamma }_{Wtb}^\mu ={\displaystyle \frac{g}{\sqrt{2}}}V_{tb}\overline{u}(p_b)\left[\gamma ^\mu (f_1^LP_L+f_1^RP_R){\displaystyle \frac{i\sigma ^{\mu \nu }k_\nu }{M_W}}(f_2^LP_L+f_2^RP_R)\right]u(p_t),`$
$`\overline{\mathrm{\Gamma }}_{Wtb}^\mu ={\displaystyle \frac{g}{\sqrt{2}}}V_{tb}^{}\overline{v}(p_{\overline{t}})\left[\gamma ^\mu (\overline{f}_1^LP_L+\overline{f}_1^RP_R){\displaystyle \frac{i\sigma ^{\mu \nu }k_\nu }{M_W}}(\overline{f}_2^LP_L+\overline{f}_2^RP_R)\right]v(p_{\overline{b}}),`$ (2.2)
where $`P_{L/R}=(1\gamma _5)/2`$, $`V_{tb}`$ is the $`(tb)`$ element of the Kobayashi-Maskawa matrix and $`k`$ is the momentum of $`W`$. In the SM $`f_1^L=\overline{f}_1^L=1`$ and all the other form factors vanish. On the other hand, it is assumed here that interactions of leptons with gauge bosons are properly described by the SM. Throughout the calculations all fermions except the top quark are considered as massless. We also neglect terms quadratic in the non-standard form factors.
Using the technique developed by Kawasaki, Shirafuji and Tsai one can derive the following formula for the inclusive distributions of the top-quark decay product $`f`$ in the process $`e^+e^{}t\overline{t}f+\mathrm{}`$ :
$$\frac{d^3\sigma }{d𝒑_f/(2p_f^0)}(e^+e^{}f+\mathrm{})=4𝑑\mathrm{\Omega }_t\frac{d\sigma }{d\mathrm{\Omega }_t}(n,0)\frac{1}{\mathrm{\Gamma }_t}\frac{d^3\mathrm{\Gamma }_f}{d𝒑_f/(2p_f^0)}(tf+\mathrm{}),$$
(2.3)
where $`\mathrm{\Gamma }_t`$ is the total top-quark decay width and $`d^3\mathrm{\Gamma }_f`$ is the differential decay rate for the process considered. $`d\sigma (n,0)/d\mathrm{\Omega }_t`$ is obtained from the angular distribution of $`t\overline{t}`$ with spins $`s_+`$ and $`s_{}`$ in $`e^+e^{}t\overline{t}`$, $`d\sigma (s_+,s_{})/d\mathrm{\Omega }_t`$, by the following replacement:
$$s_{+\mu }n_\mu ^f=\left[g_{\mu \nu }\frac{p_{t}^{}{}_{\mu }{}^{}p_{t}^{}{}_{\nu }{}^{}}{m_t^2}\right]\frac{{\displaystyle 𝑑\mathrm{\Phi }\overline{B}\mathrm{\Lambda }_+\gamma _5\gamma ^\nu B}}{{\displaystyle 𝑑\mathrm{\Phi }\overline{B}\mathrm{\Lambda }_+B}},s_\mu 0,$$
(2.4)
where the matrix element for $`t(s_+)f+\mathrm{}`$ was expressed as $`\overline{B}u_t(p_t,s_+)`$, $`\mathrm{\Lambda }_+p/_t+m_t`$, $`d\mathrm{\Phi }`$ is the relevant final-state phase-space element and $``$ denotes the appropriate spin summation.
## 3. Angular/Energy Distributions
In this section we present $`d^2\sigma /dx_fd\mathrm{cos}\theta _f`$ for the top-quark decay product $`f(=\mathrm{}^\pm /\stackrel{()}{b})`$, where $`x_f`$ denotes the normalized energy of $`f`$ defined in terms of its energy $`E_f`$ and the top-quark velocity $`\beta (\sqrt{14m_t^2/s})`$ as
$$x_f\frac{2E_f}{m_t}\sqrt{\frac{1\beta }{1+\beta }}$$
and $`\theta _f`$ is the angle between the $`e^{}`$ beam direction and the $`f`$ momentum, all in the $`e^+e^{}`$ CM frame.
Direct calculations performed in presence of the general decay vertex (2.2) lead to the following result for the $`n_\mu ^f`$ vector defined in eq.(2.4):
$$n_\mu ^f=\alpha ^f\left(g_{\mu \nu }\frac{p_{t\mu }p_{t\nu }}{m_t^2}\right)\frac{m_t}{p_tp_f}p_f^\nu $$
(3.1)
where for a given final state $`f`$, $`\alpha ^f`$ is a calculable depolarization factor
$$\alpha ^f=\{\begin{array}{cc}1\hfill & \mathrm{for}f=\mathrm{}^+\hfill \\ & \\ \frac{2r1}{2r+1}\left[\mathrm{\hspace{0.25em}1}+\frac{8\sqrt{r}(1r)}{(2r1)(2r+1)}\mathrm{Re}(f_2^R)\right]\hfill & \mathrm{for}f=b\hfill \end{array}$$
(3.2)
with $`r(M_W/m_t)^2`$. Similarly we have $`\alpha ^{\overline{f}}=\alpha ^f`$ with replacement $`f_2^R\overline{f}_2^L`$. It should be emphasized here that the above result means that there are no corrections to the “polarization vector” $`n_\mu ^{\mathrm{}}`$ for the semileptonic top-quark decay. On the other hand, one can see that the corrections to $`\alpha ^b`$ could be substantial as the kinematical suppression factor in the leading term $`2r1(=0.56`$) could be canceled by the appropriate contribution from the non-standard form factor $`f_2^R`$.
Applying the strategy described above and adopting the general formula for the $`t\overline{t}`$ distribution $`d\sigma (s_+,s_{})/d\mathrm{\Omega }_t`$ from refs., one obtains the following result for the double distribution of the angle and the rescaled energy of $`f`$ for longitudinally polarized $`e^+e^{}`$ beams:
$$\frac{d^2\sigma ^{()}}{dx_fd\mathrm{cos}\theta _f}=\frac{3\pi \beta \alpha _{\text{EM}}^2}{2s}B_f\left[\mathrm{\Theta }_0^{f()}(x_f)+\mathrm{cos}\theta _f\mathrm{\Theta }_1^{f()}(x_f)+\mathrm{cos}^2\theta _f\mathrm{\Theta }_2^{f()}(x_f)\right],$$
(3.3)
where $`\alpha _{\text{EM}}`$ is the fine structure constant and $`B_f`$ denotes the appropriate branching fraction. The energy dependence is specified by the functions $`\mathrm{\Theta }_i^{f()}(x_f)`$, explicit forms of which for unpolarized beams were shown in ref. .<sup>♯2</sup><sup>♯2</sup>♯2The functions $`\mathrm{\Theta }_i^{f()}(x_f)`$ for polarized beams could be easily obtained from formulas for unpolarized beams replacing $`D_{V,A,VA}`$, $`E_{V,A,VA}`$, $`F_{14}`$, $`G_{14}`$ defined by eq.(A.18) with $`D_{V,A,VA}^{()}`$, $`E_{V,A,VA}^{()}`$, $`F_{14}^{()}`$, $`G_{14}^{()}`$ as in eq.(A.17) in the appendix. They are parameterized both by the production and the decay form factors.
The angular distribution for $`f`$ could be easy obtained from eq.(3.3) by the integration over the energy of $`f`$:
$$\frac{d\sigma ^{()}}{d\mathrm{cos}\theta _f}_x_{}^{x_+}𝑑x_f\frac{d^2\sigma ^{()}}{dx_fd\mathrm{cos}\theta _f}=\frac{3\pi \beta \alpha _{\text{EM}}^2}{2s}B_f\left(\mathrm{\Omega }_0^{f()}+\mathrm{\Omega }_1^{f()}\mathrm{cos}\theta _f+\mathrm{\Omega }_2^{f()}\mathrm{cos}^2\theta _f\right),$$
(3.4)
where $`\mathrm{\Omega }_i^{f()}=_x_{}^{x_+}𝑑x\mathrm{\Theta }_i^{f()}`$ are shown by eq.(A.1) in the appendix and $`x_\pm `$ define kinematical energy range of $`x`$:
$$r(1\beta )/(1+\beta )x_{\mathrm{}}1\mathrm{and}(1r)(1\beta )/(1+\beta )x_b1r.$$
(3.5)
The decay vertex is entering the double distribution, eq.(3.3), through i) the functions $`F^f(x_f)`$, $`G^f(x_f)`$ and $`H_{1,2}^f(x_f)`$ defined in the appendix, and ii) the depolarization factor $`\alpha ^f`$. All the non-SM parts of $`F^f`$, $`G^f`$ and $`H_{1,2}^f`$ disappear upon integration over the energy $`x_f`$ both for $`\mathrm{}^+`$ and $`b`$, as it could be seen from the explicit forms for $`\mathrm{\Omega }_i^{f()}`$. Since $`\alpha ^f=1`$ for the leptonic distribution, we observe that the total dependence of the lepton distribution on non-standard structure of the top-quark decay vertex drops out through the integration over the energy .<sup>♯3</sup><sup>♯3</sup>♯3The same conclusion has also been reached through a different approach using the helicity formalism in ref.. However, one can expect substantial modifications for the bottom-quark distribution since corrections to $`\alpha ^b`$ could be large.
The fact that the angular leptonic distribution is insensitive to corrections to the $`VA`$ structure of the decay vertex allows for much more clear tests of the production vertices through measurements of the distribution, since that way we can avoid a contamination from a non-standard structure of the decay vertex. As an application of the angular distribution let us consider the following $`CP`$-violating forward-backward charge asymmetry:<sup>♯4</sup><sup>♯4</sup>♯4Which is an integrated version of the asymmetry we have considered in ref..
$$𝒜_{CP}^f(P_e^{},P_{e^+})=\frac{{\displaystyle _{c_m}^0}d\mathrm{cos}\theta _f{\displaystyle \frac{d\sigma ^{+()}(\theta _f)}{d\mathrm{cos}\theta _f}}{\displaystyle _0^{+c_m}}d\mathrm{cos}\theta _f{\displaystyle \frac{d\sigma ^{()}(\theta _f)}{d\mathrm{cos}\theta _f}}}{{\displaystyle _{c_m}^0}d\mathrm{cos}\theta _f{\displaystyle \frac{d\sigma ^{+()}(\theta _f)}{d\mathrm{cos}\theta _f}}+{\displaystyle _0^{+c_m}}d\mathrm{cos}\theta _f{\displaystyle \frac{d\sigma ^{()}(\theta _f)}{d\mathrm{cos}\theta _f}}},$$
(3.6)
where $`P_e^{}`$ and $`P_{e^+}`$ are the polarizations of $`e`$ and $`\overline{e}`$ beams, $`d\sigma ^{+/()}`$ is referring to $`f`$ and $`\overline{f}`$ distributions respectively, and $`c_m`$ expresses the experimental polar-angle cut. As $`\theta _f\pi \theta _{\overline{f}}`$ under $`CP`$, this asymmetry is a true measure of $`CP`$ violation. Since $`d\sigma ^{()}/d\mathrm{cos}\theta _f`$ is obtained from $`d\sigma ^{+()}/d\mathrm{cos}\theta _f`$ by reversing the sign of $`\mathrm{cos}\theta _f`$ and $`F_{1,4}^{()}`$ terms and replacing $`\alpha ^f`$ with $`\alpha ^{\overline{f}}`$ in $`\mathrm{\Omega }_{0,1,2}^{f()}`$, the asymmetry is explicitly given by the following formula
$$𝒜_{CP}^f=N_A^f/D_A^f$$
(3.7)
with (in the leading order)
$`N_A^f=2c_m\alpha _0^f\left[(1c_m^2)\mathrm{Re}(F_1^{()})+c_m\mathrm{Re}(F_4^{()})\right]\left[\mathrm{\hspace{0.25em}1}{\displaystyle \frac{1\beta ^2}{2\beta }}\mathrm{ln}{\displaystyle \frac{1+\beta }{1\beta }}\right]`$
$`c_m(1\beta ^2)\alpha _1^f\mathrm{Re}(f_2^R\overline{f}_2^L)`$
$`\times \{2(1c_m^2)\mathrm{Re}(D_{VA}^{(0,)})+c_mE_A^{(0,)}`$
$`[\mathrm{\hspace{0.25em}2}(1c_m^2)\mathrm{Re}(D_{VA}^{(0,)})+c_m(E_V^{(0,)}+E_A^{(0,)})]{\displaystyle \frac{1}{2\beta }}\mathrm{ln}{\displaystyle \frac{1+\beta }{1\beta }}\}`$
$`D_A^f=2c_m\left[\mathrm{\hspace{0.25em}1}+c_m^2\left(1{\displaystyle \frac{2}{3}}\beta ^2\right)\right]D_V^{(0,)}2c_m\left[(12\beta ^2)c_m^2\left(1{\displaystyle \frac{2}{3}}\beta ^2\right)\right]D_A^{(0,)}`$
$`4c_m(1c_m^2)\alpha _0^f(1\beta ^2)\mathrm{Re}(D_{VA}^{(0,)})`$
$`2c_m^2[\alpha _0^f(1\beta ^2)E_A^{(0,)}+2\mathrm{R}\mathrm{e}(E_{VA}^{(0,)})]`$
$`+c_m\{(1c_m^2)[D_V^{(0,)}+D_A^{(0,)}+2\alpha _0^f\mathrm{Re}(D_{VA}^{(0,)})]`$
$`+c_m[\alpha _0^f(E_V^{(0,)}+E_A^{(0,)})+2\mathrm{R}\mathrm{e}(E_{VA}^{(0,)})]\}{\displaystyle \frac{1\beta ^2}{\beta }}\mathrm{ln}{\displaystyle \frac{1+\beta }{1\beta }},`$ (3.8)
where all the coefficients are specified in the appendix, the superscript $`(0)`$ indicates the SM contribution and we expressed $`\alpha ^f`$ as $`\alpha _0^f+\alpha _1^f\mathrm{Re}(f_2^R)`$ with
$`\alpha _0^f=1,\alpha _1^f=0(\mathrm{for}f=\mathrm{}),`$
$`\alpha _0^f={\displaystyle \frac{2r1}{2r+1}},\alpha _1^f={\displaystyle \frac{8\sqrt{r}(1r)}{(1+2r)^2}}(\mathrm{for}f=b).`$
As one could have anticipated, the asymmetry for $`f=\mathrm{}`$ is sensitive to $`CP`$ violation originating exclusively from the production mechanism: It depends only on $`F_{1,4}^{()}`$ that contains contributions from the $`CP`$-violating form factors $`\delta D_\gamma `$ and $`\delta D_Z`$ while the contributing decay-vertex part consists of SM $`CP`$-conserving couplings only. For bottom quarks the effect of the modification of the decay vertex is contained in the corrections to $`b`$ and $`\overline{b}`$ depolarization factors, $`\alpha ^b+\alpha ^{\overline{b}}=\alpha _1^b\mathrm{Re}(f_2^R\overline{f}_2^L)`$, with SM $`CP`$-conserving contributions from the production process.<sup>♯5</sup><sup>♯5</sup>♯5One can show that $`f_1^{L,R}=\pm \overline{f}_1^{L,R}`$ and $`f_2^{L,R}=\pm \overline{f}_2^{R,L}`$ where upper (lower) signs are those for $`CP`$-conserving (-violating) contributions . Therefore, when only linear terms in non-standard form factors are kept, any $`CP`$-violating observable defined for the top-quark decay must be proportional to $`f_1^{L,R}\overline{f}_1^{L,R}`$ or $`f_2^{L,R}\overline{f}_2^{R,L}`$.
It will be instructive to give the following remark here: The asymmetry is defined for various initial beam polarizations $`P_{e^\pm }`$. For $`P_e^{}P_{e^+}`$, the initial state seems not to be $`CP`$ invariant and therefore one might expect contributions to the asymmetry originating from the $`CP`$-conserving part of the top-quark couplings. However, as it is seen from eq.(3.8), this is not the case. It turns out that even for $`P_e^{}P_{e^+}`$ the asymmetry is still proportional only to the $`CP`$-violating couplings embedded in $`F_{1,2,3,4}`$. The explanation is the following: Whatever the polarizations of the initial beams are, the electron (positron) beam consists of $`e(\pm 1)(\overline{e}(\pm 1))`$ where $`\pm 1`$ indicates the helicity, and only $`e(\pm 1)`$ and $`\overline{e}(1)`$ can interact non-trivially in the limit of $`m_e=0`$ since they couple to vector bosons. Therefore the interacting initial states are always $`CP`$ invariant.
Now, since we have observed in ref. that the differential version of the asymmetry discussed here could be substantial for higher collider energy, in order to illustrate the potential power of the asymmetry we present in tabs.1 and 2 (as a function of $`\sqrt{s}`$) the expected statistical significance ($`N_{SD}`$) for the asymmetry:
$$N_{SD}\frac{|𝒜_{CP}^f|}{\mathrm{\Delta }𝒜_{CP}^f}=|𝒜_{CP}^f|\sqrt{\frac{L_{eff}\sigma _{tot}}{1(𝒜_{CP}^f)^2}},$$
(3.9)
where $`L_{eff}ϵL`$ is an effective integrated luminosity for the tagging efficiency $`ϵ`$. Hereafter we adopt the integrated luminosity $`L=500\text{fb}^1`$ and the efficiency $`ϵ=60\%`$ both for lepton and $`b`$-quark detection.<sup>♯6</sup><sup>♯6</sup>♯6That low efficiency is supposed to take into account cuts necessary to suppress the background. If the $`b`$-tagging is applied then, as shown in the second paper of ref., the irreducible background to top events due $`W^\pm +2b+2j`$ is negligible, provided that a vertex tagging efficiency $`ϵ_b\stackrel{>}{}.5`$ can be achieved. Therefore for the $`b`$-tagging case the efficiency we have employed is definitely conservative. Since $`N_{SD}`$ scales as $`\sqrt{ϵL}`$ it would be easy to estimate the statistical significance for any given luminosity and efficiency. In addition, to fit the typical detector shape we impose a polar-angle cut $`|\mathrm{cos}\theta _f|<0.9`$, i.e. $`c_m=0.9`$ in eq.(3.8), both for leptons and bottom quarks. On the other hand, we will not impose any cut on the lepton/ $`b`$-quark energy since their kinematical lower bounds $`E_{\mathrm{}}^{min}=7.5`$ GeV and $`E_b^{min}=27.5`$ GeV (for $`\sqrt{s}=500`$ GeV) are large enough to be detected. Perfect angular resolution will be assumed both for lepton and $`b`$-quark final states. Also ideal leptonic-energy resolution will be used.
As it is seen from the tables, the asymmetry $`𝒜_{CP}^f`$ turned out to be a very sensitive $`CP`$-violating observable; even for unpolarized beams and $`CP`$-violating couplings of the order of $`0.05`$ one can expect $`2.4\sigma 5.5\sigma `$ effect both for lepton and $`b`$-quark asymmetries once $`L=500\text{fb}^1`$ is achieved.
The $`CP`$-violating form factors discussed here could be also generated within the SM. However, it is easy to notice that the first non-zero contribution to $`\delta D_{\gamma ,Z}`$ would require at least two loops. For the top-quark decay process $`CP`$ violation could appear at the one-loop level, however it is strongly suppressed by the double GIM mechanism . Therefore we can conclude that an experimental detection of $`CP`$-violating form factors considered here would be a clear indication for physics beyond the SM. In particular, non-vanishing $`𝒜_{CP}^l`$ in the lepton distribution will strongly indicate some new-physics in $`t\overline{t}\gamma /Z`$ couplings.
## 4. Optimal-Observable Analysis
4.1. Optimal observables
Let us briefly recall the main points of the optimal-observable (OO) technique . Suppose we have a distribution
$$\frac{d\sigma }{d\varphi }(\mathrm{\Sigma }(\varphi ))=\underset{i}{}c_if_i(\varphi )$$
(4.1)
where $`f_i(\varphi )`$ are known functions of the location in final-state phase space $`\varphi `$ and $`c_i`$’s are model-dependent coefficients. The goal would be to determine $`c_i`$’s. It can be done by using appropriate weighting functions $`w_i(\varphi )`$ such that $`𝑑\varphi w_i(\varphi )\mathrm{\Sigma }(\varphi )=c_i`$. Generally, different choices for $`w_i(\varphi )`$ are possible, but there is a unique choice so that the resultant statistical error is minimized. Such functions are given by
$$w_i(\varphi )=\underset{j}{}X_{ij}f_j(\varphi )/\mathrm{\Sigma }(\varphi ),$$
(4.2)
where $`X_{ij}`$ is the inverse matrix of $`M_{ij}`$ which is defined as
$$M_{ij}𝑑\varphi \frac{f_i(\varphi )f_j(\varphi )}{\mathrm{\Sigma }(\varphi )}.$$
(4.3)
The statistical uncertainty of $`c_i`$-determination through $`d\sigma /d\varphi `$ measurement becomes
$$\mathrm{\Delta }c_i=\sqrt{X_{ii}\sigma _T/N},$$
(4.4)
where $`\sigma _T𝑑\varphi (d\sigma /d\varphi )`$ and $`N`$ is the total number of events.
It is clear from the definition of the matrix $`M_{ij}`$, eq.(4.3), that $`M_{ij}`$ has no inverse if the functions $`f_i(\varphi )`$ are linearly dependent, and then we cannot perform any meaningful analysis. One can see it more intuitively as follows: if $`f_i(\varphi )=f_j(\varphi )`$ the splitting between $`c_i`$ and $`c_j`$ would be totally arbitrary and only $`c_i+c_j`$ could be determined.
4.2. For application
In order to apply the OO procedure to the processes under consideration, we have to reexpress the distributions in the form shown in eq.(4.1). The angular distribution, eq.(3.4), has already an appropriate form for this purpose, where $`f_i(\varphi )=\mathrm{cos}^i\theta _f`$ $`(i=0,1,2)`$ and $`\mathrm{\Omega }_i^{f()}`$ are the coefficients to be determined. On the other hand, the double angular and energy distribution eq.(3.3) must be modified. We reexpress the distribution in the following way, keeping only the SM contribution and terms linear in the non-standard form factors:
$$\frac{d^2\sigma ^{()}}{dx_fd\mathrm{cos}\theta _f}=\frac{3\pi \beta \alpha _{\text{EM}}^2}{2s}B_fS_f^{()}(x_f,\theta _f),$$
(4.5)
where
$`S_f^{()}(x_f,\theta _f)=S_f^{(0,)}(x_f,\theta _f)`$
$`+{\displaystyle \underset{v=\gamma ,Z}{}}[\mathrm{Re}(\delta A_v)_{Av}^{f()}(x_f,\theta _f)+\mathrm{Re}(\delta B_v)_{Bv}^{f()}(x_f,\theta _f)`$
$`+\mathrm{Re}(\delta C_v)_{Cv}^{f()}(x_f,\theta _f)+\mathrm{Re}(\delta D_v)_{Dv}^{f()}(x_f,\theta _f)]`$
$`+\mathrm{Re}(f_2^R)_{2R}^{f()}(x_f,\theta _f).`$
As it is seen from the above formula, the coefficients $`c_i`$ of eq.(4.1) are just the anomalous form factors to be determined. The SM contribution reads:
$$S_f^{(0,)}(x_f,\theta _f)=\mathrm{\Theta }_0^{f(0,)}(x_f)+\mathrm{cos}\theta _f\mathrm{\Theta }_1^{f(0,)}(x_f)+\mathrm{cos}^2\theta _f\mathrm{\Theta }_2^{f(0,)}(x_f)$$
(4.6)
with
$`\mathrm{\Theta }_0^{f(0,)}(x)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[(3\beta ^2)D_V^{(0,)}(13\beta ^2)D_A^{(0,)}2\alpha _0^f(1\beta ^2)\mathrm{Re}(D_{VA}^{(0,)})\right]f^f(x)`$ (4.7)
$`+2\alpha _0^f\mathrm{Re}(D_{VA}^{(0,)})g^f(x)`$
$`+{\displaystyle \frac{1}{2}}\left[D_V^{(0,)}+D_A^{(0,)}+2\alpha _0^f\mathrm{Re}(D_{VA}^{(0,)})\right]\left[\mathrm{\hspace{0.25em}2}h_1^f(x)h_2^f(x)\right],`$
$`\mathrm{\Theta }_1^{f(0,)}(x)`$ $`=`$ $`2\left[\mathrm{\hspace{0.25em}2}\mathrm{Re}(E_{VA}^{(0,)})+\alpha _0^f(1\beta ^2)E_A^{(0,)}\right]f^f(x)+2\alpha _0^f(E_V^{(0,)}+E_A^{(0,)})g^f(x)`$ (4.8)
$`2\left[\mathrm{\hspace{0.25em}2}\mathrm{Re}(E_{VA}^{(0,)})+\alpha _0^f(E_V^{(0,)}+E_A^{(0,)})\right]h_1^f(x),`$
$`\mathrm{\Theta }_2^{f(0,)}(x)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[(3\beta ^2)(D_V^{(0,)}+D_A^{(0,)})+6\alpha _0^f(1\beta ^2)\mathrm{Re}(D_{VA}^{(0,)})\right]f^f(x)`$ (4.9)
$`+2\alpha _0^f\mathrm{Re}(D_{VA}^{(0,)})g^f(x)`$
$`{\displaystyle \frac{3}{2}}\left[D_V^{(0,)}+D_A^{(0,)}+2\alpha _0^f\mathrm{Re}(D_{VA}^{(0,)})\right]\left[\mathrm{\hspace{0.25em}2}h_1^f(x)h_2^f(x)\right].`$
Explicit forms of the functions $`_{\{A,B,C,D\}\{\gamma ,Z\}}^{f()}`$ and $`_{2R}^{f()}`$ are shown in the appendix together with the functions $`f^f(x)`$, $`g^f(x)`$ and $`h_{1,2}^f(x)`$.
There are ten functions entering eq.(4.5): $`S_f^{(0,)}`$, $`_{\{A,B,C,D\}\{\gamma ,Z\}}^{f()}`$ and $`_{2R}^{f()}`$. As explained earlier, one cannot determine their coefficients separately if they are not independent. As could be found from the appendix, for the double lepton distribution, the first nine functions are linear combinations of
$`f^{\mathrm{}}(x),f^{\mathrm{}}(x)\mathrm{cos}\theta ,f^{\mathrm{}}(x)\mathrm{cos}^2\theta ,`$
$`g^{\mathrm{}}(x),g^{\mathrm{}}(x)\mathrm{cos}\theta ,g^{\mathrm{}}(x)\mathrm{cos}^2\theta ,`$
$`h_{1,2}^{\mathrm{}}(x)(13\mathrm{cos}^2\theta ),h_1^{\mathrm{}}(x)\mathrm{cos}\theta ,`$ (4.10)
while the last one, $`_{2R}^{\mathrm{}()}`$, is a combination of $`\delta \{f^{\mathrm{}},g^{\mathrm{}},h_{1,2}^{\mathrm{}}\}(x)`$ and $`\mathrm{cos}^n\theta `$ ($`n=0,1,2`$). Since there are ten coefficients to be measured,<sup>♯7</sup><sup>♯7</sup>♯7Counting the SM coefficient in front of $`S_f^{(0,)}`$ which is normalized to $`1`$. it looks always possible to determine all of them. However, it turns out not to be the case in some special cases. Indeed the possibility for the determination of all the ten form factors depends crucially on the chosen beam polarization.
This can be understood considering the invariant amplitude for $`e\overline{e}t\overline{t}`$, which could be expressed in terms of eight independent parameters as
$`(e\overline{e}t\overline{t})=C_{VV}[\overline{v}_e\gamma _\mu u_e\overline{u}_t\gamma ^\mu v_{\overline{t}}]+C_{VA}[\overline{v}_e\gamma _\mu u_e\overline{u}_t\gamma _5\gamma ^\mu v_t]`$
$`+C_{AV}[\overline{v}_e\gamma _5\gamma _\mu u_e\overline{u}_t\gamma ^\mu v_t]+C_{AA}[\overline{v}_e\gamma _5\gamma _\mu u_e\overline{u}_t\gamma _5\gamma ^\mu v_t]`$
$`+C_{VS}[\overline{v}_eq/u_e\overline{u}_tv_t]+C_{VP}[\overline{v}_eq/u_e\overline{u}_t\gamma _5v_t]`$
$`+C_{AS}[\overline{v}_e\gamma _5q/u_e\overline{u}_tv_t]+C_{AP}[\overline{v}_e\gamma _5q/u_e\overline{u}_t\gamma _5v_t].`$
However, if $`e`$ or $`\overline{e}`$ is perfectly polarized, contributions from $`[\overline{v}_e\gamma _\mu u_e]`$ and $`[\overline{v}_e\gamma _5\gamma _\mu u_e]`$ are identical. For example, when $`e`$ has fully left-handed polarization, $`u_e`$ is replaced with $`u_{eL}(1\gamma _5)u_e/2`$ and in that case they are changed as
$$\overline{v}_e\gamma _\mu u_e\overline{v}_e\gamma _\mu u_{eL},\overline{v}_e\gamma _5\gamma _\mu u_e\overline{v}_e\gamma _\mu u_{eL}.$$
Therefore, the invariant amplitude becomes
$`(e\overline{e}t\overline{t})`$
$`=(C_{VV}+C_{AV})[\overline{v}_e\gamma _\mu u_{eL}\overline{u}_t\gamma ^\mu v_{\overline{t}}]+(C_{VA}+C_{AA})[\overline{v}_e\gamma _\mu u_{eL}\overline{u}_t\gamma _5\gamma ^\mu v_t]`$
$`+(C_{VS}+C_{AS})[\overline{v}_eq/u_{eL}\overline{u}_tv_t]+(C_{VP}+C_{AP})[\overline{v}_eq/u_{eL}\overline{u}_t\gamma _5v_t],`$
and one ends with just four independent functions and therefore only four coefficients could be determined. More details could be found in the appendix below eq.(A.16). Of course, such singular configurations of polarization are not considered in our analyses.
As for $`b`$-quark distributions $`\delta f^b(x)=\delta g^b(x)=\delta h_1^b(x)=\delta h_2^b(x)=0`$, instead of ten functions $`\varphi _i(x)`$ we have in that case only nine of them given by the $`b`$-quark version eq.(4.10). Therefore, at most nine couplings could be determined. Since $`b`$-quark energy resolution is expected to be relatively poor, we will not apply OO procedure to the $`b`$-quark double distribution.
4.3. Numerical analysis
Below, we will adjust beam polarizations to perform the best measurement of the form factors. In order to gain some intuition we show in figs. 1 and 2 the functions $`_{\{A,B,C,D\}\{\gamma ,Z\}}^{\mathrm{}()}`$, $`_{2R}^{\mathrm{}()}`$ plus $`S_{\mathrm{}}^{(0,)}`$ for unpolarized beams (fig.1) and for the beam polarization $`P_e^{}=P_{e^+}=+0.5`$ (fig.2). The figures illustrates how much the polarization could modify the functions and therefore influence the possibility for the determination of the form factors.
Lepton angular distribution
Since we have only three independent functions $`\{1,\mathrm{cos}\theta ,\mathrm{cos}^2\theta \}`$, $`M`$ and its inverse $`X`$ are $`(3,3)`$ matrices. We have considered the following polarization set-ups: $`P_e^{}=P_{e^+}=0,\pm 0.5`$ and $`\pm 1`$. Since $`1>|\mathrm{cos}\theta |>\mathrm{cos}^2\theta `$ we observe that $`X_{11}<X_{22}<X_{33}`$, therefore the statistical uncertainty for $`\mathrm{\Omega }_0^{()}`$ measurement, $`\mathrm{\Delta }\mathrm{\Omega }_0^{()}`$, is always the smallest one.
Once we assume the detection efficiency $`ϵ`$ and the integrated luminosity $`L`$, we can compute the statistical significance of measuring the non-SM part of $`\mathrm{\Omega }_i^{()}`$
$$N_{SD}^{(i)}=|\mathrm{\Omega }_i^{()}\mathrm{\Omega }_i^{(0,)}|/\mathrm{\Delta }\mathrm{\Omega }_i^{()}.$$
For the efficiency and luminosity specified earlier we obtain
$`P_e^{}=P_{e^+}=0`$
$`M_{11}=2.06,M_{22}=0.55,M_{33}=0.27`$
$`X_{11}=1.09,X_{22}=1.83,X_{33}=8.40`$
$`N_{SD}^{(0)}=16.1,N_{SD}^{(1)}=3.5,N_{SD}^{(2)}=1.9`$
$`P_e^{}=P_{e^+}=+0.5`$
$`M_{11}=3.06,M_{22}=0.95,M_{33}=0.49`$
$`X_{11}=0.82,X_{22}=1.66,X_{33}=6.25`$
$`N_{SD}^{(0)}=7.5,N_{SD}^{(1)}=2.7,N_{SD}^{(2)}=1.3`$
$`P_e^{}=P_{e^+}=+1`$
$`M_{11}=3.20,M_{22}=1.25,M_{33}=0.72`$
$`X_{11}=1.00,X_{22}=2.39,X_{33}=7.00`$
$`N_{SD}^{(0)}=5.2,N_{SD}^{(1)}=3.0,N_{SD}^{(2)}=1.3`$
$`P_e^{}=P_{e^+}=0.5`$
$`M_{11}=1.27,M_{22}=0.35,M_{33}=0.17`$
$`X_{11}=1.81,X_{22}=2.91,X_{33}=13.4`$
$`N_{SD}^{(0)}=26.2,N_{SD}^{(1)}=4.9,N_{SD}^{(2)}=3.0`$
$`P_e^{}=P_{e^+}=1`$
$`M_{11}=0.76,M_{22}=0.21,M_{33}=0.11`$
$`X_{11}=3.06,X_{22}=4.90,X_{33}=22.3`$
$`N_{SD}^{(0)}=35.5,N_{SD}^{(1)}=6.4,N_{SD}^{(2)}=4.0,`$ (4.11)
where we put all the non-SM parameters $`\mathrm{Re}(\delta \{A,B,C,D\}_{\gamma ,Z})`$ and $`\mathrm{Re}(f_2^R)`$ to be $`+0.05`$ as an example.
As one can see, the precision is better for negative beam polarization, partly because of larger number of events. However we cannot conclude that using negatively-polarized beams is always more effective for new-physics search, since $`N_{SD}^{(i)}`$ strongly depends on the non-SM parameters used in the computations. In fact, positively-polarized beams give smaller $`X_{ii}`$ and this is independent of the choice of non-SM parameters. Therefore polarization of the initial beams should be carefully adjusted for each tested model in actual experimental analysis.
$`b`$-quark angular distribution
We can compute $`M`$, $`X`$ and $`N_{SD}^{(i)}`$ in the same way as for the lepton distribution:
$`P_e^{}=P_{e^+}=0`$
$`M_{11}=2.23,M_{22}=0.63,M_{33}=0.31`$
$`X_{11}=1.04,X_{22}=1.85,X_{33}=7.98`$
$`N_{SD}^{(0)}=37.3,N_{SD}^{(1)}=17.8,N_{SD}^{(2)}=1.9`$
$`P_e^{}=P_{e^+}=+0.5`$
$`M_{11}=2.38,M_{22}=0.65,M_{33}=0.32`$
$`X_{11}=0.96,X_{22}=1.61,X_{33}=7.29`$
$`N_{SD}^{(0)}=15.5,N_{SD}^{(1)}=7.6,N_{SD}^{(2)}=1.8`$
$`P_e^{}=P_{e^+}=+1`$
$`M_{11}=1.63,M_{22}=0.45,M_{33}=0.22`$
$`X_{11}=1.39,X_{22}=2.29,X_{33}=10.5`$
$`N_{SD}^{(0)}=9.7,N_{SD}^{(1)}=4.8,N_{SD}^{(2)}=2.2`$
$`P_e^{}=P_{e^+}=0.5`$
$`M_{11}=1.45,M_{22}=0.42,M_{33}=0.21`$
$`X_{11}=1.63,X_{22}=3.01,X_{33}=12.5`$
$`N_{SD}^{(0)}=62.6,N_{SD}^{(1)}=29.3,N_{SD}^{(2)}=2.4`$
$`P_e^{}=P_{e^+}=1`$
$`M_{11}=0.87,M_{22}=0.25,M_{33}=0.13`$
$`X_{11}=2.74,X_{22}=5.10,X_{33}=21.0`$
$`N_{SD}^{(0)}=85.1,N_{SD}^{(1)}=39.7,N_{SD}^{(2)}=3.1,`$ (4.12)
for $`\mathrm{Re}(\delta \{A,B,C,D\}_{\gamma ,Z})=\mathrm{Re}(f_2^R)=+0.05`$. Negatively-polarized beams give better precision again, but the same remark as to the lepton angular distribution should be kept in mind also here.
The above results prove that the optimal observables utilizing the angular distributions should be very efficient seeking for the non-SM parts of $`\mathrm{\Omega }_i^{f()}`$. However, since they are combinations of the form factors, we can only constrain them. Of course, it would be exciting if we found any signal of non-standard physics, however our final goal is to determine each form factor separately. That is why we proceed to the next analysis using the double angular and energy distributions.
Lepton angular and energy distribution
Because of high precision of direction and energy determination of leptons we adopted the double energy and angular distributions, eq.(4.5), also for OO analysis. As discussed earlier, in principle all nine form factors could be determined with the expected statistical uncertainties $`\mathrm{\Delta }c_i`$ for $`c_i=\mathrm{Re}(\delta \{A,B,C,D\}_{\gamma ,Z})`$ and $`\mathrm{Re}(f_2^R)`$. The beam polarizations $`P_e^{}`$ and $`P_{e^+}`$ were adjusted to minimize the statistical error for determination of each form factor. We found that positive polarizations lead to a smaller $`\mathrm{\Delta }c_i`$ for eight form factors in the production vertices. Unfortunately, however, the optimal polarizations for each form-factor measurement is different. Below we present the smallest statistical uncertainties and the corresponding beam polarizations for each parameter:
$`\mathrm{\Delta }[\text{Re}(\delta A_\gamma )]=0.16\mathrm{for}P_e^{}=0.7\mathrm{and}P_{e^+}=0.7,`$
$`(_{\delta C_Z:\mathrm{\hspace{0.25em}4.69},\delta D_\gamma :\mathrm{\hspace{0.25em}27.2},\delta D_Z:\mathrm{\hspace{0.25em}53.2},f_2^R:\mathrm{\hspace{0.25em}0.02}}^{\delta A_Z:\mathrm{\hspace{0.25em}0.13},\delta B_\gamma :\mathrm{\hspace{0.25em}0.25},\delta B_Z:\mathrm{\hspace{0.25em}0.49},\delta C_\gamma :\mathrm{\hspace{0.25em}2.47}}),`$
$`\mathrm{\Delta }[\text{Re}(\delta A_Z)]=0.07\mathrm{for}P_e^{}=0.5\mathrm{and}P_{e^+}=0.4,`$
$`(_{\delta C_Z:\mathrm{\hspace{0.25em}1.76},\delta D_\gamma :\mathrm{\hspace{0.25em}7.09},\delta D_Z:\mathrm{\hspace{0.25em}20.6},f_2^R:\mathrm{\hspace{0.25em}0.02}}^{\delta A_\gamma :\mathrm{\hspace{0.25em}0.23},\delta B_\gamma :\mathrm{\hspace{0.25em}0.11},\delta B_Z:\mathrm{\hspace{0.25em}0.27},\delta C_\gamma :\mathrm{\hspace{0.25em}0.70}}),`$
$`\mathrm{\Delta }[\text{Re}(\delta B_\gamma )]=0.09\mathrm{for}P_e^{}=0.2\mathrm{and}P_{e^+}=0.2,`$
$`(_{\delta C_Z:\mathrm{\hspace{0.25em}1.17},\delta D_\gamma :\mathrm{\hspace{0.25em}0.95},\delta D_Z:\mathrm{\hspace{0.25em}14.6},f_2^R:\mathrm{\hspace{0.25em}0.03}}^{\delta A_\gamma :\mathrm{\hspace{0.25em}0.43},\delta A_Z:\mathrm{\hspace{0.25em}0.11},\delta B_Z:\mathrm{\hspace{0.25em}0.36},\delta C_\gamma :\mathrm{\hspace{0.25em}0.21}}),`$
$`\mathrm{\Delta }[\text{Re}(\delta B_Z)]=0.27\mathrm{for}P_e^{}=0.4\mathrm{and}P_{e^+}=0.4,`$
$`(_{\delta C_Z:\mathrm{\hspace{0.25em}1.56},\delta D_\gamma :\mathrm{\hspace{0.25em}5.43},\delta D_Z:\mathrm{\hspace{0.25em}18.5},f_2^R:\mathrm{\hspace{0.25em}0.02}}^{\delta A_\gamma :\mathrm{\hspace{0.25em}0.25},\delta A_Z:\mathrm{\hspace{0.25em}0.07},\delta B_\gamma :\mathrm{\hspace{0.25em}0.10},\delta C_\gamma :\mathrm{\hspace{0.25em}0.56}}),`$
$`\mathrm{\Delta }[\text{Re}(\delta C_\gamma )]=0.11\mathrm{for}P_e^{}=0.1\mathrm{and}P_{e^+}=0.0,`$
$`(_{\delta C_Z:\mathrm{\hspace{0.25em}1.11},\delta D_\gamma :\mathrm{\hspace{0.25em}1.76},\delta D_Z:\mathrm{\hspace{0.25em}14.6},f_2^R:\mathrm{\hspace{0.25em}0.03}}^{\delta A_\gamma :\mathrm{\hspace{0.25em}0.82},\delta A_Z:\mathrm{\hspace{0.25em}0.22},\delta B_\gamma :\mathrm{\hspace{0.25em}0.10},\delta B_Z:\mathrm{\hspace{0.25em}0.65}}),`$
$`\mathrm{\Delta }[\text{Re}(\delta C_Z)]=1.11\mathrm{for}P_e^{}=0.1\mathrm{and}P_{e^+}=0.0,`$
$`(_{\delta C_\gamma :\mathrm{\hspace{0.25em}0.11},\delta D_\gamma :\mathrm{\hspace{0.25em}1.76},\delta D_Z:\mathrm{\hspace{0.25em}14.6},f_2^R:\mathrm{\hspace{0.25em}0.03}}^{\delta A_\gamma :\mathrm{\hspace{0.25em}0.82},\delta A_Z:\mathrm{\hspace{0.25em}0.22},\delta B_\gamma :\mathrm{\hspace{0.25em}0.10},\delta B_Z:\mathrm{\hspace{0.25em}0.65}}),`$
$`\mathrm{\Delta }[\text{Re}(\delta D_\gamma )]=0.08\mathrm{for}P_e^{}=0.2\mathrm{and}P_{e^+}=0.1,`$
$`(_{\delta C_\gamma :\mathrm{\hspace{0.25em}0.15},\delta C_Z:\mathrm{\hspace{0.25em}1.13},\delta D_Z:\mathrm{\hspace{0.25em}14.4},f_2^R:\mathrm{\hspace{0.25em}0.03}}^{\delta A_\gamma :\mathrm{\hspace{0.25em}0.52},\delta A_Z:\mathrm{\hspace{0.25em}0.13},\delta B_\gamma :\mathrm{\hspace{0.25em}0.09},\delta B_Z:\mathrm{\hspace{0.25em}0.42}}),`$
$`\mathrm{\Delta }[\text{Re}(\delta D_Z)]=14.4\mathrm{for}P_e^{}=0.2\mathrm{and}P_{e^+}=0.1,`$
$`(_{\delta C_\gamma :\mathrm{\hspace{0.25em}0.15},\delta C_Z:\mathrm{\hspace{0.25em}1.13},\delta D_\gamma :\mathrm{\hspace{0.25em}0.08},f_2^R:\mathrm{\hspace{0.25em}0.03}}^{\delta A_\gamma :\mathrm{\hspace{0.25em}0.52},\delta A_Z:\mathrm{\hspace{0.25em}0.13},\delta B_\gamma :\mathrm{\hspace{0.25em}0.09},\delta B_Z:\mathrm{\hspace{0.25em}0.42}}),`$ (4.13)
where we also showed the expected precision of the other parameter measurements for the same beam polarizations. For instance, we can expect $`\mathrm{\Delta }[\text{Re}(\delta A_\gamma )]=0.16`$ for $`P_e^{}=P_{e^+}=0.7`$ while the expected precision of $`\text{Re}(\delta A_Z)`$, $`\text{Re}(\delta B_\gamma )`$, $`\mathrm{}`$ for the same polarizations are 0.13, 0.25, $`\mathrm{}`$, respectively. This result is independent of the choice of the non-SM parameters in contrast to the preceding results.
As it is seen the precision of $`\delta \{C,D\}_Z`$ measurement would be very poor even for the optimal polarization. This is mainly a consequence of the size of $`_{\{C,D\}Z}^{\mathrm{}()}`$, which is illustrated in figs.1 and 2: These two functions are very small in a large area. More quantitatively, the size of the elements of $`M`$ matrix, $`M_{ij}`$, is $`O(1)`$ for $`i,j7,9`$, while the size of $`M_{i7}(=M_{7i})`$ and $`M_{i9}(=M_{9i})`$ is at most $`O(10^2)`$. In addition, determination of $`\delta D_\gamma `$ would be practically difficult, as well, since its error varies rapidly with the polarization. For example, $`\mathrm{\Delta }[\text{Re}(\delta D_\gamma )]`$ becomes 0.86 for $`P_e^{}=0.1/P_{e^+}=0.1`$ and 0.99 for $`P_e^{}=0.3/P_{e^+}=0.1`$. The source of that sensitivity is hidden in the neutral-current structure with $`\mathrm{sin}^2\theta _W0.23`$. Indeed, the optimal polarization becomes $`P_e^{}=0.1`$ instead of 0.2 $`(\mathrm{\Delta }[\text{Re}(\delta D_\gamma )]=0.09)`$ for $`\mathrm{sin}^2\theta _W=0.25`$. On the other hand, a good determination (almost independently of the polarization) could be expected for $`f_2^R`$. Indeed, the best precision is
$$\mathrm{\Delta }[\text{Re}(f_2^R)]=0.01\mathrm{for}P_e^{}=0.8\mathrm{and}P_{e^+}=0.8$$
(4.14)
whereas even for the unpolarized beams we obtain $`\mathrm{\Delta }[\text{Re}(f_2^R)]=0.03`$.
At the time a linear collider will be operating, data from Tevatron Run II and LHC will also provide independent constraints on top-quark couplings. Below we provide an example of a combined analysis assuming $`\delta A_v`$, $`\delta B_v`$ and $`f_2^R`$ are known and OO are used to determine $`\delta C_v`$ and $`\delta D_v`$ only (here we put $`\delta A_v=\delta B_v=f_2^R=0`$ for simplicity). The results are as follows:
$$\begin{array}{cc}\mathrm{\Delta }[\text{Re}(\delta C_\gamma )]=0.04\hfill & \mathrm{for}P_e^{}=0.2\mathrm{and}P_{e^+}=0.2\hfill \\ \mathrm{\Delta }[\text{Re}(\delta C_Z)]=0.23\hfill & \mathrm{for}P_e^{}=0.2\mathrm{and}P_{e^+}=0.1\hfill \\ \mathrm{\Delta }[\text{Re}(\delta D_\gamma )]=0.03\hfill & \mathrm{for}P_e^{}=0.2\mathrm{and}P_{e^+}=0.1\hfill \\ \mathrm{\Delta }[\text{Re}(\delta D_Z)]=2.97\hfill & \mathrm{for}P_e^{}=0.2\mathrm{and}P_{e^+}=0.1\hfill \end{array}$$
(4.15)
The error for $`\delta D_Z`$ became much smaller but still too large for practical use. However, as we have seen in sec.3, the $`CP`$-sensitive asymmetry $`𝒜_{CP}^f`$ would provide much stronger constraints on $`\delta D_{\gamma ,Z}`$.
## 5. Summary and Conclusions
We have presented here the angular and energy distributions for $`\stackrel{()}{f}`$ in the process $`e^+e^{}t\overline{t}\stackrel{()}{f}\mathrm{}`$, where $`f=\mathrm{}`$ or $`b`$ quark in the form suitable for an application of the optimal observables (OO). The most general ($`CP`$-violating and $`CP`$-conserving) couplings for $`\gamma t\overline{t}`$, $`Zt\overline{t}`$ and $`Wtb`$ have been assumed. All fermion masses except $`m_t`$ have been neglected and we have kept only terms linear in anomalous couplings. We have assumed the tagging efficiency at the level of $`60\%`$ both for lepton and $`b`$ quark detection, the range of the polar angle restricted by $`|\mathrm{cos}\theta _f|<0.9`$ and the integrated luminosity $`L=500\text{fb}^1`$.
$`CP`$-violating charge forward-backward asymmetry $`𝒜_{CP}^f`$ has been introduced as an efficient way for testing $`CP`$-violation in top-quark couplings. Since the angular distribution for leptons is insensitive to variations of the standard V$``$A structure of the $`Wtb`$ coupling, the asymmetry could be utilized for a pure test of $`CP`$-violation in the top-quark production process. The expected statistical significance $`N_{SD}`$ for the measurement of the asymmetry has been calculated. We have found that it should be possible to detect $`𝒜_{CP}^f`$ at $`5.5\sigma `$ ($`4.3\sigma `$) level for bottom quarks (leptons) for unpolarized beams, assuming $`CP`$-violating couplings of the order of $`0.05`$. Having both beams polarized at $`80\%`$ the signal for bottom quarks (leptons) could reach even $`16\sigma `$ ($`6.6\sigma `$).
Next, the OO procedure has been applied to the angular distributions. In the case of the lepton angular distribution, the expected statistical significance for signals of non-standard physics varies between $`1.3\sigma `$ and 35.5$`\sigma `$ assuming non-standard form factors of the order of 0.05. It turned out that in the case of the bottom-quark angular distribution the statistical significance of the signal is in general higher than for leptons because of larger event rate and varies between $`1.8\sigma `$ and 85.1$`\sigma `$ for the same non-standard form factors.
When deriving the above results we have fixed all the non-SM parameters to be $`+0.05`$ as a reasonable example for the strength of beyond-the-SM physics. However, final results for statistical significances considered here depend on the size of the non-standard parameters. The most convenient beam polarizations for a measurement of the asymmetry $`𝒜_{CP}^f`$ and for testing the angular distributions varies with the non-standard parameters, as well. Therefore one should stress that the beam polarizations should be carefully adjusted for each model to be tested in actual experimental analysis. However, in any case, the above results show that a measurement of $`𝒜_{CP}^f`$ and OO analysis of the angular distributions are both very efficient for new-physics search.
Then we have analyzed the angular and energy distribution of the lepton toward separate determinations of the anomalous form factors. In order to reach the highest precision we have been adjusting beam polarizations to minimize errors for each form factor. We have found that at $`\sqrt{s}=500`$ GeV with the integrated luminosity $`L=500\text{fb}^1`$ the best determined coupling would be the axial coupling of the $`Z`$ boson with the error $`\mathrm{\Delta }[\text{Re}(\delta A_Z)]=0.07`$ while the lowest precision is expected for $`\text{Re}(\delta D_Z)`$ with $`\mathrm{\Delta }[\text{Re}(\delta D_Z)]=14.4`$. This result is independent of the choice of the non-SM parameters in contrast to the above two types of analyses.
Concluding, we have observed that the angular distributions and the angular and energy distributions of top-quark decay products both provide very efficient tools for studying top-quark couplings to gauge bosons at linear colliders.
ACKNOWLEDGMENTS
We would like to thank K. Fujii for useful discussions concerning details of lepton and bottom-quark detection and J. Pliszka for his remarks on the statistical analysis. One of us (Z.H.) is grateful to Y. Sumino, T. Nagano, T. Takahashi and K. Ikematsu for stimulating discussions. This work is supported in part by the State Committee for Scientific Research (Poland) under grant 2 P03B 014 14 and by Maria Skłodowska-Curie Joint Fund II (Poland-USA) under grant MEN/NSF-96-252.
Appendix
Integrals of $`\mathrm{\Theta }_i^{f()}(x)`$ denoted in the main text by $`\mathrm{\Omega }_i^{f()}`$ in the angular distribution eq.(3.4) are the following:
$`\mathrm{\Omega }_0^{f()}=D_V^{()}(12\beta ^2)D_A^{()}2\mathrm{Re}(G_1^{()})`$
$`\alpha ^f[\mathrm{\hspace{0.25em}2}(1\beta ^2)\mathrm{Re}(D_{VA}^{()})\mathrm{Re}(F_1^{()})+(32\beta ^2)\mathrm{Re}(G_3^{()})]`$
$`+[D_V^{()}+D_A^{()}+2\mathrm{Re}(G_1^{()})`$
$`+\alpha ^f\mathrm{Re}(2D_{VA}^{()}F_1^{()}+3G_3^{()})]{\displaystyle \frac{1\beta ^2}{2\beta }}\mathrm{ln}{\displaystyle \frac{1+\beta }{1\beta }},`$
$`\mathrm{\Omega }_1^{f()}=4\mathrm{Re}(E_{VA}^{()})+2\alpha ^f[(1\beta ^2)E_A^{()}\mathrm{Re}(F_4^{()}G_2^{()})]`$
$`\left\{2\mathrm{R}\mathrm{e}(E_{VA}^{()})+\alpha ^f[E_V^{()}+E_A^{()}\mathrm{Re}(F_4^{()}G_2^{()})]\right\}{\displaystyle \frac{1\beta ^2}{\beta }}\mathrm{ln}{\displaystyle \frac{1+\beta }{1\beta }},`$
$`\mathrm{\Omega }_2^{f()}=(32\beta ^2)[D_V^{()}+D_A^{()}+2\mathrm{R}\mathrm{e}(G_1^{()})]`$
$`+3\alpha ^f[\mathrm{\hspace{0.25em}2}(1\beta ^2)\mathrm{Re}(D_{VA}^{()})\mathrm{Re}(F_1^{()})+(32\beta ^2)\mathrm{Re}(G_3^{()})]`$
$`3[D_V^{()}+D_A^{()}+2\mathrm{Re}(G_1^{()})`$
$`+\alpha ^f\mathrm{Re}(2D_{VA}^{()}F_1^{()}+3G_3^{()})]{\displaystyle \frac{1\beta ^2}{2\beta }}\mathrm{ln}{\displaystyle \frac{1+\beta }{1\beta }}.`$ (A.1)
Next we present explicit formulas of the coefficient functions for the nine anomalous form factors in eq.(4.5) $`_{\{A,B,C,D\}\{\gamma ,Z\}}^{f()}(x,\theta )`$ and $`_{2R}^{f()}(x,\theta )`$ ($`f=\mathrm{}/b`$):
$`_{Av}^{f()}(x,\theta )`$
$`=[{\displaystyle \frac{1}{2}}(3\beta ^2)C(D_V:A_v)f^f(x)+2\alpha _0^fC(D_{VA}:A_v)g^f(x)](1+\mathrm{cos}^2\theta )`$
$`[\alpha _0^f(1\beta ^2)C(D_{VA}:A_v)f^f(x)`$
$`{\displaystyle \frac{1}{2}}\{C(D_V:A_v)+2\alpha _0^fC(D_{VA}:A_v)\}\{2h_1^f(x)h_2^f(x)\}](13\mathrm{cos}^2\theta )`$
$`+2[\alpha _0^fC(E_V:A_v)\{g^f(x)h_1^f(x)\}+2C(E_{VA}:A_v)\{f^f(x)h_1^f(x)\}]\mathrm{cos}\theta ,`$
(A.2)
$`_{Bv}^{f()}(x,\theta )`$
$`={\displaystyle \frac{1}{2}}\beta ^2C(D_A:B_v)f^f(x)(3\mathrm{cos}^2\theta )+2\alpha _0^fC(D_{VA}:B_v)g^f(x)(1+\mathrm{cos}^2\theta )`$
$`{\displaystyle \frac{1}{2}}[\{C(D_A:B_v)+2\alpha _0^f(1\beta ^2)C(D_{VA}:A_v)\}f^f(x)`$
$`\{C(D_A:B_v)+2\alpha _0^fC(D_{VA}:B_v)\}\{2h_1^f(x)h_2^f(x)\}](13\mathrm{cos}^2\theta )`$
$`+2[\{\alpha _0^f(1\beta ^2)C(E_A:B_v)+2C(E_{VA}:B_v)\}f^f(x)+\alpha _0^fC(E_A:B_v)g^f(x)`$
$`\{\alpha _0^fC(E_A:B_v)+2C(E_{VA}:B_v)\}h_1^f(x)]\mathrm{cos}\theta ,`$ (A.3)
$`_{Cv}^{f()}(x,\theta )`$
$`=\beta ^2C(G_1:C_v)f^f(x)(1+\mathrm{cos}^2\theta )`$
$`+2\alpha _0^fC(G_2:C_v)[f^f(x)+g^f(x)h_1^f(x)]\mathrm{cos}\theta `$
$`[\{C(G_1:C_v)+\alpha _0^f(2\beta ^2)C(G_3:C_v)\}f^f(x)+\alpha _0^fC(G_3:C_v)g^f(x)`$
$`\{2C(G_1:C_v)+3\alpha _0^fC(G_3:C_v)\}h_1^f(x)`$
$`+\{C(G_1:C_v)+\alpha _0^fC(G_3:C_v)\}h_2^f(x)](13\mathrm{cos}^2\theta )`$ (A.4)
$`_{Dv}^{f()}(x,\theta )`$
$`=\alpha _0^fC(F_1:D_v)[f^f(x)h_1^f(x)](13\mathrm{cos}^2\theta )\alpha _0^fC(F_1:D_v)g^f(x)(1+\mathrm{cos}^2\theta )`$
$`2\alpha _0^fC(F_4:D_v)[f^f(x)+g^f(x)h_1^f(x)]\mathrm{cos}\theta ,`$ (A.5)
while $`_{2R}^{\mathrm{}()}(x,\theta )`$ takes different forms for $`f=\mathrm{}`$ and $`f=b`$ as
$`_{2R}^{\mathrm{}()}(x,\theta )`$
$`={\displaystyle \frac{1}{2}}\left[(3\beta ^2)D_V^{(0,)}(13\beta ^2)D_A^{(0,)}2(1\beta ^2)\mathrm{Re}(D_{VA}^{(0,)})\right]\delta f^{\mathrm{}}(x)`$
$`+2\mathrm{R}\mathrm{e}(D_{VA}^{(0,)})\delta g^{\mathrm{}}(x)(1+\mathrm{cos}^2\theta )`$
$`+{\displaystyle \frac{1}{2}}\left[D_V^{(0,)}+D_A^{(0,)}+2\mathrm{R}\mathrm{e}(D_{VA}^{(0,)})\right]\left[\mathrm{\hspace{0.25em}2}\delta h_1^{\mathrm{}}(x)\delta h_2^{\mathrm{}}(x)\right](13\mathrm{cos}^2\theta )`$
$`+2\left[(1\beta ^2)E_A^{(0,)}+2\mathrm{R}\mathrm{e}(E_{VA}^{(0,)})\right]\delta f^{\mathrm{}}(x)\mathrm{cos}\theta `$
$`+2(E_V^{(0,)}+E_A^{(0,)})\delta g^{\mathrm{}}(x)\mathrm{cos}\theta `$
$`2\left[E_V^{(0,)}+E_A^{(0,)}+2\mathrm{R}\mathrm{e}(E_{VA}^{(0,)})\right]\delta h_1^{\mathrm{}}(x)\mathrm{cos}\theta `$
$`+{\displaystyle \frac{1}{2}}\left[(3\beta ^2)(D_V^{(0,)}+D_A^{(0,)})+6(1\beta ^2)\mathrm{Re}(D_{VA}^{(0,)})\right]\delta f^{\mathrm{}}(x)\mathrm{cos}^2\theta ,`$ (A.6)
and
$`_{2R}^{b()}(x,\theta )`$
$`=\alpha _1^b\{\mathrm{Re}(D_{VA}^{(0,)})[\{(1\beta ^2)f^b(x)2h_1^b(x)+h_2^b(x)\}(13\mathrm{cos}^2\theta )`$
$`+2g^b(x)(1+\mathrm{cos}^2\theta )]`$
$`+2[(1\beta ^2)E_A^{(0,)}f^b(x)+(E_V^{(0,)}+E_A^{(0,)})\{g^b(x)h_1^b(x)\}]\mathrm{cos}\theta \},`$ (A.7)
where the functions $`f^f(x)`$, $`g^f(x)`$, $`h_{1,2}^f(x)`$, $`\delta f^f(x)`$, $`\delta g^f(x)`$ and $`\delta h_{1,2}^f(x)`$ are defined as
$`F^f(x)=f^f(x)+\mathrm{Re}(f_2^R)\delta f^f(x),`$
$`G^f(x)=g^f(x)+\mathrm{Re}(f_2^R)\delta g^f(x),`$
$`H_{1,2}^f(x)=h_{1,2}^f(x)+\mathrm{Re}(f_2^R)\delta h_{1,2}^f(x),`$ (A.8)
with $`F^f(x)`$, $`G^f(x)`$ and $`H_{1,2}^f(x)`$ being given as follows
$`F^f(x){\displaystyle \frac{1}{B_f}}{\displaystyle 𝑑\omega \frac{1}{\mathrm{\Gamma }_t}\frac{d^2\mathrm{\Gamma }_f}{dxd\omega }},G^f(x){\displaystyle \frac{1}{B_f}}{\displaystyle 𝑑\omega \left[\mathrm{\hspace{0.25em}1}x\frac{1+\beta }{1\omega }\right]\frac{1}{\mathrm{\Gamma }_t}\frac{d^2\mathrm{\Gamma }_f}{dxd\omega }},`$
$`H_1^f(x){\displaystyle \frac{1}{B_f}}{\displaystyle \frac{1\beta }{x}}{\displaystyle 𝑑\omega (1\omega )\frac{1}{\mathrm{\Gamma }_t}\frac{d^2\mathrm{\Gamma }_f}{dxd\omega }},`$
$`H_2^f(x){\displaystyle \frac{1}{B_f}}\left({\displaystyle \frac{1\beta }{x}}\right)^2{\displaystyle 𝑑\omega (1\omega )^2\frac{1}{\mathrm{\Gamma }_t}\frac{d^2\mathrm{\Gamma }_f}{dxd\omega }},`$ (A.9)
and $`\omega `$ is defined as $`\omega (p_tp_f)^2/m_t^2`$.
After performing the above integrations using
$$\frac{1}{\mathrm{\Gamma }_t}\frac{d^2\mathrm{\Gamma }_f}{dxd\omega }=\{\begin{array}{cc}\frac{1+\beta }{\beta }\frac{3B_{\mathrm{}}}{W}\omega \left[\mathrm{\hspace{0.25em}1}+2\mathrm{R}\mathrm{e}(f_2^R)\sqrt{r}\left(\frac{1}{1\omega }\frac{3}{1+2r}\right)\right]\hfill & \mathrm{for}f=\mathrm{}^+,\hfill \\ & \\ \frac{1+\beta }{2\beta (1r)}\delta (\omega r)\hfill & \mathrm{for}f=b.\hfill \end{array}$$
one obtains the following explicit forms of $`f^f(x)`$, $`g^f(x)`$, $`h_{1,2}^f(x)`$, $`\delta f^f(x)`$, $`\delta g^f(x)`$ and $`\delta h_{1,2}^f(x)`$ for leptonic and bottom-quark final states:
$``$ For $`f=\mathrm{}`$
$`f^{\mathrm{}}(x)={\displaystyle \frac{3(1+\beta )}{2\beta W}}[\omega ^2]_\omega _{}^{\omega _+}\left({\displaystyle \frac{3(1+\beta )}{2\beta W}}(\omega _+^2\omega _{}^2)\right),`$
$`g^{\mathrm{}}(x)=f^{\mathrm{}}(x)+{\displaystyle \frac{3(1+\beta )^2}{\beta W}}x\left[\omega +\mathrm{ln}|1\omega |\right]_\omega _{}^{\omega _+},`$
$`h_1^{\mathrm{}}(x)={\displaystyle \frac{1\beta ^2}{2\beta W}}{\displaystyle \frac{1}{x}}\left[\omega ^2(32\omega )\right]_\omega _{}^{\omega _+},`$
$`h_2^{\mathrm{}}(x)={\displaystyle \frac{1}{4\beta W}}(1+\beta )(1\beta )^2{\displaystyle \frac{1}{x^2}}\left[\omega ^2(68\omega +3\omega ^2)\right]_\omega _{}^{\omega _+},`$
$`\delta f^{\mathrm{}}(x)={\displaystyle \frac{3(1+\beta )}{\beta W}}\sqrt{r}\left[\mathrm{\hspace{0.25em}2}\omega +2\mathrm{ln}|1\omega |+{\displaystyle \frac{3\omega ^2}{1+2r}}\right]_\omega _{}^{\omega _+},`$
$`\delta g^{\mathrm{}}(x)=\delta f^{\mathrm{}}(x){\displaystyle \frac{6(1+\beta )^2}{\beta W}}\sqrt{r}x[\mathrm{ln}|1\omega |`$
$`+{\displaystyle \frac{1}{1\omega }}+{\displaystyle \frac{3}{1+2r}}(\omega +\mathrm{ln}|1\omega |)]^{\omega _+}_\omega _{},`$
$`\delta h_1^{\mathrm{}}(x)={\displaystyle \frac{3(1\beta ^2)}{\beta W}}{\displaystyle \frac{\sqrt{r}}{x}}\left[\omega ^2\left(1{\displaystyle \frac{32\omega }{1+2r}}\right)\right]_\omega _{}^{\omega _+},`$
$`\delta h_2^{\mathrm{}}(x)={\displaystyle \frac{1}{2\beta W}}(1+\beta )(1\beta )^2`$
$`\times {\displaystyle \frac{\sqrt{r}}{x^2}}\left[\mathrm{\hspace{0.25em}2}\omega ^2(32\omega ){\displaystyle \frac{3\omega ^2}{1+2r}}(68\omega +3\omega ^2)\right]_\omega _{}^{\omega _+},`$ (A.10)
where $`\omega _\pm `$ are given as follows:
For $`rB`$ ($`rM_W^2/m_t^2`$ and $`B(1\beta )/(1+\beta )`$)
$$\begin{array}{cccc}\omega _+=1r,\hfill & \omega _{}=1x/B\hfill & \mathrm{for}\hfill & Brx<B\hfill \\ \omega _+=1r,\hfill & \omega _{}=0\hfill & \mathrm{for}\hfill & Bx<r\hfill \\ \omega _+=1x,\hfill & \omega _{}=0\hfill & \mathrm{for}\hfill & rx1\hfill \end{array}$$
(A.11)
For $`r<B`$
$$\begin{array}{cccc}\omega _+=1r,\hfill & \omega _{}=1x/B\hfill & \mathrm{for}\hfill & Brx<r\hfill \\ \omega _+=1x,\hfill & \omega _{}=1x/B\hfill & \mathrm{for}\hfill & rx<B\hfill \\ \omega _+=1x,\hfill & \omega _{}=0\hfill & \mathrm{for}\hfill & Bx1\hfill \end{array}$$
(A.12)
$``$ For $`f=b`$
$`f^b(x)={\displaystyle \frac{1+\beta }{2\beta (1r)}}(=\mathrm{constant}),`$
$`g^b(x)=\left(1{\displaystyle \frac{1+\beta }{1r}}x\right){\displaystyle \frac{1+\beta }{2\beta (1r)}},`$
$`h_1^b(x)={\displaystyle \frac{1\beta ^2}{2\beta x}},`$
$`h_2^b(x)={\displaystyle \frac{(1r)(1+\beta )(1\beta )^2}{2\beta x^2}},`$
$`\delta f^b(x)=\delta g^b(x)=\delta h_1^b(x)=\delta h_2^b(x)=0,`$ (A.13)
where $`x`$ is bounded as
$$B(1r)x1r.$$
The coefficients $`C(X:Y)`$ employed in the definition of the coefficient functions have been introduced through the following formulas:
$`D_V^{()}=D_V^{(0,)}+{\displaystyle \underset{v=\gamma ,Z}{}}C(D_V:A_v)\mathrm{Re}(\delta A_v),`$
$`D_A^{()}=D_A^{(0,)}+{\displaystyle \underset{v=\gamma ,Z}{}}C(D_A:B_v)\mathrm{Re}(\delta B_v),`$
$`\mathrm{Re}(D_{VA}^{()})=\mathrm{Re}(D_{VA}^{(0,)})`$
$`+{\displaystyle \underset{v=\gamma ,Z}{}}[C(D_{VA}:A_v)\mathrm{Re}(\delta A_v)+C(D_{VA}:B_v)\mathrm{Re}(\delta B_v)],`$ (A.14)
and in the analogous manner for $`E_{V,A,VA}`$, $`F_{14}`$ and $`G_{14}`$. $`D_{V,A,VA}^{(0,)}`$, $`E_{V,A,VA}^{(0,)}`$, $`F_{14}^{(0,)}`$, $`G_{14}^{(0,)}`$ could be obtained from eq.(A.17) below as a SM approximation of $`D_{V,A,VA}^{()}`$, $`E_{V,A,VA}^{()}`$, $`F_{14}^{()}`$, $`G_{14}^{()}`$. Explicit forms of the independent coefficients are given as
$`C(D_V:A_\gamma )=2C[𝒫_{}A_\gamma (𝒫_{}+v_e𝒫_{})d^{}A_Z],`$
$`C(E_V:A_\gamma )=2C[𝒫_{}A_\gamma (𝒫_{}+v_e𝒫_{})d^{}A_Z],`$
$`C(D_{VA}:A_\gamma )=C(𝒫_{}+v_e𝒫_{})d^{}B_Z,`$
$`C(E_{VA}:A_\gamma )=C(𝒫_{}+v_e𝒫_{})d^{}B_Z,`$
$`C(D_V:A_Z)=2C[(𝒫_{}+v_e𝒫_{})d^{}A_\gamma \{2v_e𝒫_{}+(1+v_e^2)𝒫_{}\}d^2A_Z],`$
$`C(E_V:A_Z)=2C[(𝒫_{}+v_e𝒫_{})d^{}A_\gamma \{2v_e𝒫_{}+(1+v_e^2)𝒫_{}\}d^2A_Z],`$
$`C(D_{VA}:A_Z)=C[\mathrm{\hspace{0.25em}2}v_e𝒫_{}+(1+v_e^2)𝒫_{}]d^2B_Z,`$
$`C(E_{VA}:A_Z)=C[\mathrm{\hspace{0.25em}2}v_e𝒫_{}+(1+v_e^2)𝒫_{}]d^2B_Z,`$ (A.15)
where $`v_e=1+4\mathrm{sin}^2\theta _W`$, $`d^{}s/[4\mathrm{sin}\theta _W\mathrm{cos}\theta _W(sM_Z^2)]`$, two polarization factors $`𝒫_{}`$ and $`𝒫_{}`$ are defined as
$$𝒫_{}P_e^{}+P_{e^+},𝒫_{}1+P_e^{}P_{e^+},$$
and the others are thereby given as
$$\begin{array}{cc}C(D_A:B_v)=2C(D_{VA}:A_v),\hfill & C(E_A:B_v)=2C(E_{VA}:A_v),\hfill \\ C(D_{VA}:B_v)=C(D_V:A_v)/2,\hfill & C(E_{VA}:B_v)=C(E_V:A_v)/2,\hfill \\ C(G_1:C_v)=C(D_V:A_v)/2,\hfill & C(G_2:C_v)=C(E_V:A_v)/2,\hfill \\ C(G_3:C_v)=C(D_{VA}:A_v),\hfill & C(G_4:C_v)=C(E_{VA}:A_v),\hfill \\ C(F_1:D_v)=C(D_V:A_v)/2,\hfill & C(F_2:D_v)=C(E_V:A_v)/2,\hfill \\ C(F_3:D_v)=C(D_{VA}:A_v),\hfill & C(F_4:D_v)=C(E_{VA}:A_v).\hfill \end{array}$$
(A.16)
As explained in the main text, they are not always independent of each other. When $`P_{}=\pm P_{}`$, i.e., $`P_e^{}=P_{e^+}=\pm 1`$, we have
$$C(\{D_V,E_V,D_{VA},E_{VA}\}:A_Z)=(1\pm v_e)d^{}C(\{D_V,E_V,D_{VA},E_{VA}\}:A_\gamma ),$$
As a consequence of the above relations one gets
$$_{\{A,B,C,D\}Z}^{f()}(x,\theta )=(1\pm v_e)d^{}_{\{A,B,C,D\}\gamma }^{f()}(x,\theta ).$$
In this case all we can determine (for the production form factors) are the following four combinations
$$\mathrm{Re}(\delta \{A,B,C,D\}_\gamma (1\pm v_e)d^{}\delta \{A,B,C,D\}_Z).$$
Finally we present here formulas for $`D_{V,A,VA}^{()}`$, $`E_{V,A,VA}^{()}`$, $`F_{14}^{()}`$, $`G_{14}^{()}`$ for completeness:
$`D_{V,A,VA}^{()}=𝒫_{}D_{V,A,VA}𝒫_{}E_{V,A,VA},`$
$`E_{V,A,VA}^{()}=𝒫_{}E_{V,A,VA}𝒫_{}D_{V,A,VA},`$
$`F_{1,\mathrm{\hspace{0.17em}2},\mathrm{\hspace{0.17em}3},\mathrm{\hspace{0.17em}4}}^{()}=𝒫_{}F_{1,\mathrm{\hspace{0.17em}2},\mathrm{\hspace{0.17em}3},\mathrm{\hspace{0.17em}4}}𝒫_{}F_{2,\mathrm{\hspace{0.17em}1},\mathrm{\hspace{0.17em}4},\mathrm{\hspace{0.17em}3}},`$
$`G_{1,\mathrm{\hspace{0.17em}2},\mathrm{\hspace{0.17em}3},\mathrm{\hspace{0.17em}4}}^{()}=𝒫_{}G_{1,\mathrm{\hspace{0.17em}2},\mathrm{\hspace{0.17em}3},\mathrm{\hspace{0.17em}4}}𝒫_{}G_{2,\mathrm{\hspace{0.17em}1},\mathrm{\hspace{0.17em}4},\mathrm{\hspace{0.17em}3}},`$ (A.17)
for
$`D_VC[A_\gamma ^22A_\gamma A_Zv_ed^{}+A_Z^2(1+v_e^2)d^2+2(A_\gamma A_Zv_ed^{})\mathrm{Re}(\delta A_\gamma )`$
$`2\{A_\gamma v_ed^{}A_Z(1+v_e^2)d^2\}\mathrm{Re}(\delta A_Z)],`$
$`D_AC[B_Z^2(1+v_e^2)d^22B_Zv_ed^{}\mathrm{Re}(\delta B_\gamma )+2B_Z(1+v_e^2)d^2\mathrm{Re}(\delta B_Z)],`$
$`D_{VA}C[A_\gamma B_Zv_ed^{}+A_ZB_Z(1+v_e^2)d^2B_Zv_ed^{}(\delta A_\gamma )^{}`$
$`+(A_\gamma v_ed^{}A_Z)\delta B_\gamma +B_Z(1+v_e^2)d^2(\delta A_Z)^{}`$
$`\{A_\gamma v_ed^{}A_Z(1+v_e^2)d^2\}\delta B_Z],`$
$`E_V\mathrm{\hspace{0.25em}2}C[A_\gamma A_Zd^{}A_Z^2v_ed^2+A_Zd^{}\mathrm{Re}(\delta A_\gamma )+(A_\gamma d^{}2A_Zv_ed^2)\mathrm{Re}(\delta A_Z)],`$
$`E_A\mathrm{\hspace{0.25em}2}C[B_Z^2v_ed^2+B_Zd^{}\mathrm{Re}(\delta B_\gamma )2B_Zv_ed^2\mathrm{Re}(\delta B_Z)],`$
$`E_{VA}C[A_\gamma B_Zd^{}2A_ZB_Zv_ed^2+B_Zd^{}(\delta A_\gamma )^{}+A_Zd^{}\delta B_\gamma `$
$`2B_Zv_ed^2(\delta A_Z)^{}+(A_\gamma d^{}2A_Zv_ed^2)\delta B_Z],`$
$`F_1C[(A_\gamma A_Zv_ed^{})\delta D_\gamma +\{A_\gamma v_ed^{}A_Z(1+v_e^2)d^2\}\delta D_Z],`$
$`F_2C[A_Zd^{}\delta D_\gamma (A_\gamma d^{}2A_Zv_ed^2)\delta D_Z],`$
$`F_3C[B_Zv_ed^{}\delta D_\gamma B_Z(1+v_e^2)d^2\delta D_Z],`$
$`F_4C[B_Zd^{}\delta D_\gamma +2B_Zv_ed^2\delta D_Z],`$
$`G_1C[(A_\gamma A_Zv_ed^{})\delta C_\gamma \{A_\gamma v_ed^{}A_Z(1+v_e^2)d^2\}\delta C_Z],`$
$`G_2C[A_Zd^{}\delta C_\gamma +(A_\gamma d^{}2A_Zv_ed^2)\delta C_Z],`$
$`G_3C[B_Zv_ed^{}\delta C_\gamma +B_Z(1+v_e^2)d^2\delta C_Z],`$
$`G_4C[B_Zd^{}\delta C_\gamma 2B_Zv_ed^2\delta C_Z]`$ (A.18)
with $`C1/(4\mathrm{sin}^2\theta _W)`$.
$``$ Note added after Publication $``$
\[ Corrigendum \] After this article has been published in Nucl. Phys. B585 (2000), 3, we have found that equation (A.3) contains an error: $`C(D_{VA}:A_v)`$ in the third line should be replaced with $`C(D_{VA}:B_v)`$ as
$`_{Bv}^{f()}(x,\theta )`$
$`={\displaystyle \frac{1}{2}}\beta ^2C(D_A:B_v)f^f(x)(3\mathrm{cos}^2\theta )+2\alpha _0^fC(D_{VA}:B_v)g^f(x)(1+\mathrm{cos}^2\theta )`$
$`{\displaystyle \frac{1}{2}}[\{C(D_A:B_v)+2\alpha _0^f(1\beta ^2)C(D_{VA}:B_v)\}f^f(x)`$
$`\{C(D_A:B_v)+2\alpha _0^fC(D_{VA}:B_v)\}\{2h_1^f(x)h_2^f(x)\}](13\mathrm{cos}^2\theta )`$
$`+2[\{\alpha _0^f(1\beta ^2)C(E_A:B_v)+2C(E_{VA}:B_v)\}f^f(x)+\alpha _0^fC(E_A:B_v)g^f(x)`$
$`\{\alpha _0^fC(E_A:B_v)+2C(E_{VA}:B_v)\}h_1^f(x)]\mathrm{cos}\theta .`$ (A.3)
Due to this correction, the two graphs expressing $`_{B\gamma }^{\mathrm{}()}`$ and $`_{BZ}^{\mathrm{}()}`$ in Figs.1 and 2 are to be replaced with those presented below:
The numerical results shown in Eqs.(4.13) and (4.14) are also no longer valid, and we have carried out re-computations. Concerning the former, i.e., Eq.(4.13), after correcting the error we find very large statistical uncertainties for measurements of the nine independent non-SM parameters, therefore, in practice it will be impossible (with no other experimental input) to determine all of them at once through the distribution that was considered, i.e., the one in Eq.(4.5).
Among those non-SM couplings, however, $`\delta A_\gamma `$ term is directly related to the top-quark electric charge and expected to be studied in various other ways. We therefore would like to give the results of an analysis without $`\delta A_\gamma `$ term and replace Eq.(4.13) with
$`\mathrm{\Delta }[\text{Re}(\delta A_Z)]=4.0\times 10^2\mathrm{for}P_e^{}/P_{e^+}=0.4/0.4,`$
$`(_{\delta C_Z:\mathrm{\hspace{0.25em}1.47},\delta D_\gamma :\mathrm{\hspace{0.25em}5.25},\delta D_Z:\mathrm{\hspace{0.25em}17.8},f_2^R:\mathrm{\hspace{0.25em}0.02}}^{\delta B_\gamma :\mathrm{\hspace{0.25em}0.08},\delta B_Z:\mathrm{\hspace{0.25em}0.06},\delta C_\gamma :\mathrm{\hspace{0.25em}0.52}}),`$
$`\mathrm{\Delta }[\text{Re}(\delta B_\gamma )]=7.2\times 10^2\mathrm{for}P_e^{}/P_{e^+}=0.2/0.3,0.3/0.2,`$
$`(_{\delta C_Z:\mathrm{\hspace{0.25em}1.17},\delta D_\gamma :\mathrm{\hspace{0.25em}1.86},\delta D_Z:\mathrm{\hspace{0.25em}14.6},f_2^R:\mathrm{\hspace{0.25em}0.03}}^{\delta A_Z:\mathrm{\hspace{0.25em}0.04},\delta B_Z:\mathrm{\hspace{0.25em}0.05},\delta C_\gamma :\mathrm{\hspace{0.25em}0.25}}),`$
$`\mathrm{\Delta }[\text{Re}(\delta B_Z)]=4.5\times 10^2\mathrm{for}P_e^{}/P_{e^+}=0.2/0.3,0.3/0.2,`$
$`(_{\delta C_Z:\mathrm{\hspace{0.25em}1.17},\delta D_\gamma :\mathrm{\hspace{0.25em}1.86},\delta D_Z:\mathrm{\hspace{0.25em}14.6},f_2^R:\mathrm{\hspace{0.25em}0.03}}^{\delta A_Z:\mathrm{\hspace{0.25em}0.04},\delta B_\gamma :\mathrm{\hspace{0.25em}0.07},\delta C_\gamma :\mathrm{\hspace{0.25em}0.25}}),`$
$`\mathrm{\Delta }[\text{Re}(\delta C_\gamma )]=1.0\times 10^1\mathrm{for}P_e^{}/P_{e^+}=0.1/0.1,`$
$`(_{\delta C_Z:\mathrm{\hspace{0.25em}1.07},\delta D_\gamma :\mathrm{\hspace{0.25em}0.81},\delta D_Z:\mathrm{\hspace{0.25em}13.9},f_2^R:\mathrm{\hspace{0.25em}0.03}}^{\delta A_Z:\mathrm{\hspace{0.25em}0.06},\delta B_\gamma :\mathrm{\hspace{0.25em}0.08},\delta B_Z:\mathrm{\hspace{0.25em}0.07},}),`$
$`\mathrm{\Delta }[\text{Re}(\delta C_Z)]=1.1\times 10^0\mathrm{for}P_e^{}/P_{e^+}=0.1/0.1,`$
$`(_{\delta D_\gamma :\mathrm{\hspace{0.25em}0.81},\delta D_Z:\mathrm{\hspace{0.25em}13.9},f_2^R:\mathrm{\hspace{0.25em}0.03}}^{\delta A_Z:\mathrm{\hspace{0.25em}0.06},\delta B_\gamma :\mathrm{\hspace{0.25em}0.08},\delta B_Z:\mathrm{\hspace{0.25em}0.07},\delta C_\gamma :\mathrm{\hspace{0.25em}0.10}}),`$
$`\mathrm{\Delta }[\text{Re}(\delta D_\gamma )]=6.9\times 10^2\mathrm{for}P_e^{}/P_{e^+}=0.1/0.2,0.2/0.1,`$
$`(_{\delta C_Z:\mathrm{\hspace{0.25em}1.08},\delta D_Z:\mathrm{\hspace{0.25em}13.9},f_2^R:\mathrm{\hspace{0.25em}0.03}}^{\delta A_Z:\mathrm{\hspace{0.25em}0.05},\delta B_\gamma :\mathrm{\hspace{0.25em}0.07},\delta B_Z:\mathrm{\hspace{0.25em}0.06},\delta C_\gamma :\mathrm{\hspace{0.25em}0.13}}),`$
$`\mathrm{\Delta }[\text{Re}(\delta D_Z)]=1.4\times 10^{+1}\mathrm{for}P_e^{}/P_{e^+}=0.1/0.2,0.2/0.1,`$
$`(_{\delta C_Z:\mathrm{\hspace{0.25em}1.08},\delta D_\gamma :\mathrm{\hspace{0.25em}0.07},f_2^R:\mathrm{\hspace{0.25em}0.03}}^{\delta A_Z:\mathrm{\hspace{0.25em}0.05},\delta B_\gamma :\mathrm{\hspace{0.25em}0.07},\delta B_Z:\mathrm{\hspace{0.25em}0.06},\delta C_\gamma :\mathrm{\hspace{0.25em}0.13}}).`$ (4.13)
On the other hand, Eq.(4.14) is simply to be replaced by
$$\mathrm{\Delta }[\text{Re}(f_2^R)]=1.5\times 10^2\mathrm{for}P_e^{}=0.9\mathrm{and}P_{e^+}=0.9.$$
(4.14)
In spite of these modifications, conclusions concerning Eq.(4.13) are not affected substantially and hold except for those on $`\delta A_\gamma `$, if only we properly adjust the parameter values used there according to the above corrected eqs.(4.13) and (4.14).
We would like to thank very much Patrick Janot for kindly pointing out that one term in eq.(A.3) seems unnatural and therefore might be a typo. This led us to rechecking that equation and finding the error mentioned here.
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# References
The ”recoil” correction $`m\alpha ^6`$ to hyperfine splitting of positronium ground state.
A.P.Burichenko
## Abstract
The ”recoil” correction of order $`m\alpha ^6`$ to the hyperfine splitting of positronium ground state was found. The formalism employed is based on the noncovariant perturbation theory in QED. Equation for two-particle component of full (many-body) wave function is used, in which effective Hamiltonian depends on the energy of a system. The effective Hamiltonian is not restricted to the nonrelativistic region, so there is no need in any regularization. To evaluate integrals over loop momenta, they are divided into ”hard” and ”soft” parts, coming from large and small momenta respectively. Soft contributions were found analytically, and hard ones are evaluated by numerical integration. Some soft terms due to the retardation cancel each other. To calculate the ”hard” contributions, a great number of noncovariant graphs is replaced by only a few covariant ones. The hard contribution was found in two ways. The first way is to evaluate contributions of separate graphs, using the Coulomb gauge. The second one is to calculate full hard contribution as a whole using the Feynman gauge. The final result for the ”recoil” correction is $`0.381(6)m\alpha ^6`$ and agrees with those of . Diagram-to-diagram comparison with the revised results of was done. All the results agree, so the ”recoil” correction is now firmly established. This means a considerable disagreement with the experimental data.
1. INTRODUCTION.
The treatment of relativistic bound states is one of challenging tasks of Quantum Electrodinamics. Positronium, the bound state of the electron and the positron, is one of the most appropriate objects for theoretical and experimental study of relativistic bound states. Because of the small masses of its constituents, effects of the strong and the weak interaction are negligible compared with the accuracy of current experiments on the positronium spectroscopy. On the other hand this experimental accuracy is good enough for comparison with results of modern theoretical investigations, which reached accuracy level $`m\alpha ^6`$ for contributions to the energy. Therefore, now positronium spectrum can be investigated within QED framework.
Now the most accurately measured positronium property is hyperfine splitting of the ground state, i.e. energy difference between the $`1^3S_1`$ and $`1^1S_0`$-states, denoted below as $`\mathrm{\Delta }\nu `$. The two best experimental results for this value are
$$\mathrm{\Delta }\nu =\mathrm{203\; 387.5}(1.6)MHz,$$
(1)
$$\mathrm{\Delta }\nu =\mathrm{203\; 389.10}(0.74)MHz;$$
(2)
obtained in and respectively.
Calculation of $`\mathrm{\Delta }\nu `$ has a long history; contributions $`m\alpha ^4`$, $`m\alpha ^5`$, $`m\alpha ^6\mathrm{ln}\alpha `$ were found in ; their sum equals
$$m\alpha ^4\left[\frac{7}{12}\frac{\alpha }{\pi }\left(\frac{8}{9}+\frac{1}{2}\mathrm{ln}2\right)\frac{5}{24}\alpha ^2\mathrm{ln}\alpha \right]=\mathrm{203\; 400.29}MHz.$$
(3)
To compare experimental results (1,2) with the theoretical one, the contributions $`m\alpha ^6`$ (without the logarithm) are to be obtained. An essential progress in their calculation was made only recently.
There are several sets of corrections $`m\alpha ^6`$ to the hyperfine splitting, that are different in origin and may be calculated independently. First, these are three sets consisting of contributions arising from the one-, two-, and three-photon annihilation; they were found in , and respectively. Second, there are contributions associated with radiative corrections to the Breit potential, i.e. those of formal order $`(Z\alpha )^4\alpha ^2m`$), found in . Third, there exist ”recoil” corrections, i.e. those coming from the graphs in which all photon lines connect two fermion ones; these corrections are of conventional order $`(Z\alpha )^nm`$ (here $`Ze`$ is the charge of one of the particles; $`Z=1`$ in positronium). Finally, there are radiative-recoil (i.e. $`(Z\alpha )^5\alpha m`$) corrections found in .
Below the sum of contributions $`m(Z\alpha )^6,m(Z\alpha )^6\mathrm{ln}\alpha `$ to $`\mathrm{\Delta }\nu `$ is denoted as $`\mathrm{\Delta }\nu _{rec}`$. Calculation of its nonlogarithmic part is the most difficult step of calculation of $`\mathrm{\Delta }\nu `$. This was the subject of works . The present work is also devoted to calculation of $`\mathrm{\Delta }\nu _{rec}`$.
The results of the three first works, namely , were all different. The result of work coincides essentially with the result of (they are equal to $`0.3763m\alpha ^6`$ and $`0.3767(17)m\alpha ^6`$ respectively). The result of the present work is $`0.381(6)m\alpha ^6`$ and also coincides, within its accuracy, with the results of . Recently the result of was corrected by its authors. Now it is equal to $`0.3764(35)m\alpha ^6`$ and agrees with the results mentioned above.
In the so–called NRQED (nonrelativistic QED), an effective field theory equivalent to QED, was formulated and then applied to calculation of $`\mathrm{\Delta }\nu _{rec}`$; nonlogarithmic part of the result equals $`0.167(33)m\alpha ^6`$. However, just now the preliminary result was obtained in the NRQED framework , which is consistent with that of . In the calculation was performed by Bethe-Salpeter formalism. In the calculation was done by an effective Hamiltonian approach; the method employed in in fact also uses an effective Hamiltonian. Difference between and is in different regularizations used; besides that, the result of was obtained by numerical integration, whereas that of is analytic. The effective Hamiltonian approach is essentially a combination of ideas of the NRQED and the old–fashioned noncovariant perturbation theory in QED. The latter was used first for QED bound state calculations in .
The present paper also employes a formalism based on old-fashioned noncovariant (”time-ordered”) perturbation theory for QED. Starting from the Shrödinger equation for full (many-body) wave function one easily obtains an equation containing only two-particle part of the wave function, in which effective Hamiltonian depends on the energy of a system. Then the effective Hamiltonian is divided into unperturbed part and perturbation, zeroth approximation for the two-particle wave function is found, and using them corrections to the energy levels are calculated by means of usual Rayleigh–Shrödinger perturbation theory. This calculation consist essentially in finding of expectation values of operators corresponding to graphs of the noncovariant perturbation theory over the ”unperturbed” wave function. To evaluate the integrals over loop momenta each of the integrands is divided into two parts, ”hard” and ”soft”; here terms ”hard” and ”soft” mean that in the desired order a ”hard” part entirely arises from the region of momenta $`m`$, whereas for a ”soft” part the region of momenta $`\alpha m`$ is also essential. This decomposition is performed in such a way that ”soft” contributions may be easily found analytically, and ”hard” ones are found by numerical integration. To calculate the ”hard” contributions with the required accuracy external legs of the graphs may be set on the mass shell, which allows to replace the sum of a great number of noncovariant graphs by the sum of only a few covariant ones.
The method of calculation described above differs from that of and in the way effective Hamiltonian is defined. In and , effective Hamiltonian is constructed so as to reproduce the scattering amplitudes, whereas in the present paper it is immediately derived from the full QED Hamiltonian. Besides that, the effective Hamiltonian used in the present work is not restricted to the nonrelativistic region, in difference with and , so there is no need in regularization of the effective Hamiltonian or matrix elements it enters.
The plan of the rest of the paper is following. In section 2, I briefly describe the formalism used in the paper, i.e. the effective ”Shrödinger equation” for two-particle component of the wave function, ”zeroth” approximation to the wave function, perturbation theory formulas for the calculation of $`\mathrm{\Delta }\nu `$; these formulas are transformed into form of expectation values over the nonrelativistic Coulomb wave function. In section 3, I describe the general method of calculation of these expectation values, namely, the explicit manner to divide them into ”soft” and ”hard” parts; correspondence between these values and covariant graphs (and also scattering amplitudes on mass shell) is established. In section 4 this general method is applied to contributions of various sets of graphs in turn. Section 5 contains the description of the procedure of numerical evaluation of ”hard” contributions, the checking of this procedure, and comparison of the results obtained with the results of . Section 6 consist of summary including comparison with experimental results.
2. FORMALISM.
The starting point of the formalism used in the present paper is the Shrödinger equation
$$H|\psi _N>=E|\psi _N>,$$
(4)
where the positronium wave function $`|\psi _N>`$ includes components with different number of particles. In (4) $`H`$ is the full QED Hamiltonian, $`H=H_0+V_N,`$ where $`H_0`$ is the free Hamiltonian, $`V_N`$ is the interaction.
Direct use of the many-body wave function is unconvenient. However in the lowest order in $`\alpha `$ the only nonvanishing component of $`|\psi _N>`$ is two-particle one, which is denoted as $`|e>`$. It is easy to obtain an equation containing only $`|e>`$. Let $`P_e`$ be projector onto the two-particle subspace, $`P_a=1P_e`$,
$$|e>P_e|\psi _N>,|a>P_a|\psi _N>.$$
From the Shrödinger equation one gets
$$P_eHP_e|e>+P_eHP_a|a>=E|e>,$$
(5)
$$P_aHP_e|e>+P_aHP_a|a>=E|a>.$$
(6)
Expressing $`|a>`$ from (6) and substituting it into (5) results in
$$(P_eHP_e+V_a(E)E)|e>=0,$$
(7)
where
$$V_a(E)=P_eHP_a\frac{1}{(EP_aHP_a)}P_aHP_e.$$
Energy levels which could be obtained from (7) are derived in a form of perturbative expansion. For this purpose ”Hamiltonian” $`P_eHP_e+V_a(E)`$ is divided into the unperturbed part $`H^{(0)}P_eH_0P_e+V_0`$ and the perturbation $`V(E)V_a(E)+P_eV_NP_eV_0,`$ where $`V_0`$ is the usual nonrelativistic Coulomb interaction. Contributions to the energy, which are of 1-st, 2-nd, and 3-rd order in $`V`$, have leading orders (besides logarithms) $`m\alpha ^4,m\alpha ^5`$ and $`m\alpha ^6`$ respectively.
Clearly $`V(E)+V_0=V_a(E)+P_eV_NP_e`$ is the two-particle irreducible kernel for the Green function appearing in the noncovariant perturbation theory (for transitions between two-particle states). Its expansion in $`\alpha `$ is
$$V(E)+V_0=P_eV_N\left(\mathrm{\hspace{0.33em}1}+\frac{P_a}{EH_0}V_N+\frac{P_a}{EH_0}V_N\frac{P_a}{EH_0}V_N+\mathrm{}\mathrm{}\mathrm{}\right)P_e$$
and corresponds to the set of two-particle irreducible graphs of the noncovariant perturbation theory (here and below ”irreducible” is meant in noncovariant sense, for instance, graph shown in fig.1 is irreducible).
Since diagram technique for the noncovariant theory is not of common use, its rules are described here. External fermion lines correspond to positive-energy spinors normalized as $`u^+u=1`$ (for convenience, the fermion-fermion channel, instead of the fermion-antifermion one, is considered). Internal fermion lines are described by the projectors
$$\mathrm{\Lambda }^\pm (𝐩)=(\epsilon _p\pm (\gamma _0+\alpha 𝐩))/(2\epsilon _p),\epsilon _p\sqrt{𝐩^2+m^2};$$
factors corresponding to external lines may be written in form
$$u(𝐩)=((2\epsilon _p)/(\epsilon _p+m))^{1/2}\mathrm{\Lambda }^+(𝐩)\left(\begin{array}{c}1\\ 0\end{array}\right)w,$$
$`\left(\begin{array}{c}1\\ 0\end{array}\right)`$ is $`(2\times 4)`$ matrix, $`w`$ is a two-spinor normalized as $`w^+w=1.`$ The two fermions are assumed to have charges of opposite sign; then Coulomb photon with the momentum $`𝐪`$ gives the factor $`4\pi \alpha /𝐪^2,`$ and magnetic (transverse) photon gives the factor $`(4\pi \alpha )(\stackrel{}{\alpha }_1\stackrel{}{\alpha }_2(\stackrel{}{\alpha }_1𝐪)(\stackrel{}{\alpha }_2𝐪)/q^2)/2q,`$ (in the noncovariant perturbation theory energy of virtual photon is equal to its momentum). Every negative-energy projector $`\mathrm{\Lambda }^{}`$ gives extra factor $`(1)`$. Factor $`(EE_k+i0)^1`$ corresponds to each intermediate state, where $`E_k`$ is energy of the intermediate state, $`E`$ is full energy of a system. Loop momenta $`𝐪_i`$ which remain undetermined lead to the integration $`d^3q_i/(2\pi )^3`$.
In the rest of the paper all operators and wave functions refer to the two-particle subspace. All the calculations are performed in the center-of-mass frame, so the only variable except spins is $`𝐩`$, i.e. momentum of particle 1 in the c.m.f.. Eigenstates of $`𝐩`$ are normalized as usual, according to $`<𝐩^{}|𝐩>=(2\pi )^3\delta ^3(𝐩𝐩^{}),`$ kernel of an arbitrary operator $`X`$ is denoted as $`X(𝐩,𝐩^{})<𝐩^{}|X|𝐩>`$.
The kernel $`V(𝐩,𝐩^{};E)+V_0(𝐩,𝐩^{})`$ is given by the sum of matrix elements corresponding to all the irreducible graphs for initial and final states having momenta $`𝐩`$ and $`𝐩^{}`$ respectively. The tree diagram with Coulomb photon corresponds to operator $`V_c`$, and that with magnetic photon corresponds to $`V_m`$. Let us write $`V`$ in the form
$$V(E)=V_1+V_2+V_3+\mathrm{}\mathrm{},$$
(8)
where
$$V_1=V_{1c}+V_m,V_{1c}V_cV_0,$$
$`V_2,V_3`$ and so on correspond to sums of the irreducible graphs having two, three, and more, photon lines respectively. Being expanded in $`v/c`$ up to the second order, $`H^{(0)}+V_1`$ gives the Breit Hamiltonian.
Consider, for example, the first-order in $`V`$ correction to energy levels. The first term in (8) gives to this correction contribution of formal order $`m\alpha ^4`$, the second term gives $`m\alpha ^5`$, and so on. However in fact this expansion does not converge: for instance, spin-independent correction to the energy, arising from $`V_2,V_3`$ and further terms of expansion (8), are all of order $`m\alpha ^5`$. This takes place due to graphs shown in fig. 2: (a), (b), (c), and so on, because these graphs produce ”ultrasoft” contributions, i.e. those coming from the region where momentum of the virtual magnetic photon is $`m\alpha ^2`$. Fortunately for the hyperfine splitting the ”rest term” of this expansion corresponds to the graphs first of which is graph shown at fig 2(c), and has order $`m\alpha ^7`$.
When deriving formulas for corrections to energy levels the ground state of positronium may be regarded as non-degenerate, because S-states with $`\sigma =0`$ and $`\sigma =1`$ (where $`\sigma `$ is the positronium full spin) do not mix. For perturbation $`V`$ depending on full energy of a system, correction to non-degenerate energy level number $`n`$ reads
$$\mathrm{\Delta }E_n=V_{nn}+\underset{mn}{}\frac{|V_{nm}|^2}{E_nE_m}+\underset{kn}{}\underset{mn}{}\frac{V_{nm}V_{mk}V_{kn}}{(E_mE_n)(E_kE_n)}+$$
$$+V_{nn}\frac{V_{nn}}{E_n}+\frac{}{E_n}\left(V_{nn}\underset{mn}{}\frac{|V_{nm}|^2}{E_nE_m}\right)+\frac{}{E_n}\left(\frac{1}{2}V_{nn}^2\frac{V_{nn}}{E_n}\right),$$
(9)
where summation over continuous part of the spectrum, as well as over discrete one, is implicit; $`E_m,E_n,E_k`$ are zeroth-order energy levels; all matrix elements and their derivatives are taken at $`E=E_n`$. It may be easily shown that for calculation of $`\mathrm{\Delta }\nu `$ to order $`m\alpha ^6`$ it is enough to use only three first terms in (9). Keeping only terms of the order desired, the recoil contribution to $`\mathrm{\Delta }\nu `$ equals (up to the order $`m\alpha ^6`$)
$$<\psi |V_1G^{}V_1G^{}V_1+V_2G^{}V_1+V_1G^{}V_2+V_3+V_1G^{}V_1+V_2+V_1|\psi >|_{\sigma =0}^{\sigma =1},$$
(10)
where $`|\psi >`$ is zeroth approximation to $`|e>`$, i.e. solution of equation
$$H^{(0)}|\psi >=E_2|\psi >,$$
(11)
and $`G^{}=G^{}(E_2)`$ is Green function of equation (11) with the ground state pole subtracted out (the Green function without the subtraction is denoted as $`G`$).
As $`H^{(0)}`$ does not depend on fermion spins, $`|\psi >`$ may be assumed to have the form
$$|\psi >=\frac{d^3p}{(2\pi )^3}\phi (𝐩)|𝐩>|\chi >,$$
where $`|\chi >`$ is spin part of the wave function. So to find $`|\psi >`$ is essentially to find $`\phi (𝐩)`$, i.e. to solve equation
$$(T+V_0E_2)\phi (𝐩)=0$$
(12)
where $`T`$ is kinetic energy of the two particles (including the mass).
Nonrelativistic approximation for (12) is
$$(T_0+V_0E_0)\phi _0(𝐩)=0,$$
(13)
where $`T_0`$ is nonrelativistic approximation for $`T`$. The ground state solution for (13) is
$$E_0=2m\frac{\gamma ^2}{m},\phi _0(𝐩)=8\gamma ^{5/2}\pi ^{1/2}\frac{1}{f_p^2}$$
(here and below notations $`\gamma \alpha m/2,f_k𝐤^2+\gamma ^2`$ are used, for an arbitrary momentum $`𝐤`$).
Let us write $`\phi `$ as
$$\phi =C_0\phi _0+\mathrm{\Delta }\phi ,$$
where $`\mathrm{\Delta }\phi `$ is orthogonal to $`\phi _0`$; then (12) is rewritten as
$$(T_0+V_0E_0)\mathrm{\Delta }\phi =(\mathrm{\Delta }T+\mathrm{\Delta }E)\phi ,$$
(14)
where $`\mathrm{\Delta }EE_2E_0,\mathrm{\Delta }TTT_0`$. Solution of (14) may be written in the form
$$\mathrm{\Delta }\phi =G_0^{}(E_0)(\mathrm{\Delta }T\mathrm{\Delta }E)\phi ,$$
(15)
where $`G_0^{}`$ is Green function of equation (13) (i.e. usual nonrelativistic Coulomb Green function) with the ground state pole subtracted out (corresponding Green function without the subtraction is denoted as $`G_0`$).
Relativistic free two-particle Green function and nonrelativistic free one are referred to as $`S_0`$ and $`S`$:
$$S_0(E^{})=(E^{}T_0)^1,$$
$$S(E^{})=(E^{}T)^1;$$
$`G`$ and $`G_0`$ are expressed as expansions
$$G_0(E^{})=(E^{}(T_0+V_0))^1=S_0+S_0V_0S_0+\mathrm{}..,$$
$$G(E^{})=(E^{}(T+V_0))^1=S+SV_0S+\mathrm{}...$$
(16)
Let us also denote
$$LG_0^{}S_0S_0V_0S_0.$$
It may be shown by iteration of (15) that with the sufficient accuracy its solution is
$$\phi (1+(SS_0)V_0+SV_0(SS_0)V_0+L\mathrm{\Delta }T)|_{E^{}=E_0}\phi _0.$$
(17)
It is convenient to use in the rest of the paper notations $`<>`$ and $`<>_{(n)},`$ defined so that for any operator $`X`$
$$<X>\phi _0^+(𝐩^{})X(𝐩,𝐩^{},E)|_{\sigma =0}^{\sigma =1}\phi _0(𝐩)\frac{d^3p}{(2\pi )^3}\frac{d^3p^{}}{(2\pi )^3},$$
and $`<X>_{(n)}`$ stands for the sum of contributions $`\alpha ^n,\alpha ^n\mathrm{ln}\alpha `$ to $`<X>`$.
If one rewrites (10) using
$$G^{}(S+SV_0S+L)|_{E^{}=E_0},$$
expression for $`\phi `$ from $`\phi _0`$ (17), and keeping only terms of the relevant order, one obtains
$$\mathrm{\Delta }\nu _{rec}=<V_3+V_2+U_{V2}+U_C+U_M+U_{MM}+U_{MCM}+U_{MMM}+U_L>_{(6)},$$
(18)
where
$$U_{V2}=(V_2SV_m+V_2(SV_cS_0V_0))+h.c.,$$
(19)
$$U_C=V_{1c}+U_{C2}+U_{C3},$$
(20)
$$U_M=V_m+U_{M2}+U_{M3},$$
(21)
$$U_{MM}=U_{MM2}+U_{MM3}$$
(22)
$$U_{MCM}=V_mSV_cSV_m,$$
(23)
$$U_{MMM}=V_mSV_mSV_m,$$
(24)
$$U_L=V_1LV_1+(V_1L\mathrm{\Delta }T+h.c.),$$
(25)
and
$$U_{C2}=V_{1c}SV_{1c}+(V_{1c}(SS_0)V_0+h.c.),$$
(26)
$$U_{M2}=V_m(SV_cS_0V_0)+h.c.,$$
(27)
$$U_{MM2}=V_mSV_m,$$
(28)
$$U_{C3}=V_0(SS_0)V_{1c}(SS_0)V_0+(V_{1c}SV_0(SS_0)V_0+h.c.)+$$
(29)
$$+(V_{1c}SV_{1c}(SS_0)V_0+h.c.)+V_{1c}SV_0SV_{1c}+V_{1c}SV_{1c}SV_{1c},$$
$$U_{M3}=(V_cSV_0S_0)V_m(SV_cS_0V_0)+(V_mSV_c(SV_cS_0V_0)+h.c.),$$
(30)
$$U_{MM3}=(V_mSV_m(SV_cS_0V_0)+h.c.).$$
(31)
3. THE METHOD OF THE CALCULATION.
In this section method of calculation of contributions to $`\mathrm{\Delta }\nu _{rec}`$, written out in (18) - (31), is presented. The main idea of this method is just the idea used in other modern works concerning bound states, though there is some difference in its implementation. Every contribution to $`\mathrm{\Delta }\nu _{rec}`$ is divided essentially into two parts, ”soft” and ”hard” ones. Here the following is undermined. Contributions to $`\mathrm{\Delta }\nu _{rec}`$ are in fact integrals over loop momenta. Any contribution under the consideration is called ”hard” if it is determined (with accuracy required) by region where all the loop momenta are of order $`m`$. Otherwise this contribution is called ”soft”. To separate ”soft” contributions from ”hard” ones, the integrands are to be expanded in powers of the momenta. In order to know whether given term of such expansion contain ”soft” contribution, all the momenta should be set to $`\alpha m`$, and then simplest power counting let us know whether this momentum area give rise contribution $`m\alpha ^6`$ to $`\mathrm{\Delta }\nu _{rec}`$ (for the term under the consideration). If this is the case, the term considered clearly must be treated as ”soft” contribution. On the other hand it can be proved that in the present problem absence of the contribution $`m\alpha ^6`$ due to the momenta region mentioned above lead to absence of ”soft” contributions of the relevant order at all, at least for the way of calculation described in the present paper.
It is clear that there is some freedom in explicit way of the decomposition into ”soft” and ”hard” parts, and ”soft” part can always be defined so that its contribution to the integral is easily evaluated analytically. On the other hand, calculating ”hard” part one can put $`E_0=2m`$ in the integrands, and the integrals obtained can be easily evaluated by means of numerical integration.
Contributions to $`\mathrm{\Delta }\nu _{rec}`$ are naturally divided into three sets, different in form and calculation procedure used. These are ”tree” contributions ($`<V_{1c},V_m>_{(6)}`$), ”one-loop” ($`<V_2,U_{C2},U_{M2},U_{MM2}>_{(6)}`$), and ”two-loop” ones ($`<V_3,U_{V2},U_{C3},U_{M3},U_{MM3},U_{MCM},U_{MMM}>_{(6)}`$) (here words ”one-” and ”two-loop” imply that the operators inside $`<>_{(6)}`$ are integrals over one and two momenta respectively), and also term $`<U_L>_{(6)}`$ that corresponds to graphs having three loops or more; the last is determined by region of momenta $`m\alpha `$ and can be easily found analytically.
Consider different kinds of the contributions one by one.
Each two-loop contribution has the form $`<X_2>_{(6)}`$ where
$$X_2(𝐩,𝐩^{},\gamma )=\frac{d^3q_1}{(2\pi )^3}\frac{d^3q_2}{(2\pi )^3}\alpha ^3Y_2(𝐩,𝐩^{},𝐪_1,𝐪_2,\gamma ).$$
(32)
It is convenient to define the value $`n`$, the ”divergency power” of the $`X_2`$; namely, $`Y_2\delta ^{n6}`$ at $`pp^{}q_1q_2\gamma \delta m`$. All the ”two-loop” contributions have $`n0`$; if $`n=0`$ then region $`q_1q_2m\alpha `$ contribute to order $`m\alpha ^6`$ and $`<X_2>_{(6)}`$ may contain terms $`m\alpha ^6\mathrm{ln}\alpha `$. It can be shown that for all the ”two-loop” terms having $`n>0`$ the only essential contribution to the integral arises at $`pp^{}\gamma ,q_1q_2|𝐪_1𝐪_2|m`$; dependence of $`Y_2`$ on $`𝐩,𝐩^{}`$ in this area can be neglected that results in
$$<X_2>_{(6)}|\phi _0(0)|^2X_2|_{𝐩=𝐩^{}=\gamma =0}|_{\sigma =0}^{\sigma =1}equiv<X_2>_p,$$
(33)
where
$$\phi _0(0)=\frac{d^3p}{(2\pi )^3}\phi _0(𝐩)=\sqrt{\frac{\gamma ^3}{\pi }}.$$
The integral in (33) does not depend on the small parameter $`\alpha `$ and so can be easily computed numerically. For the sake of brevity notation $`<>_p`$ defined in (33) is used throughout the paper.
If $`n=0`$ the integrand $`Y_2`$ is to be divided into two parts: the soft part, denoted as $`Y_{21}`$, and the hard one, $`Y_{22}`$, $`Y_2=Y_{21}+Y_{22},`$ so that $`Y_{21},Y_{22}`$ have ”divergency power” $`n=0`$ and $`n>0`$ respectively. It may be done in such a way that $`Y_{21}`$ equals approximately a uniform function at $`p,p^{},q_1,q_2,\gamma m`$ (the only exception is the contribution of graph shown at fig.2(b), treated in the section 4). Similarly $`X_2=X_{21}+X_{22}`$ where $`X_{21},X_{22}`$ are integrals of form (32) with the $`Y_2`$ replaced by $`Y_{21},Y_{22}`$ respectively. It may be shown that with such choice of $`Y_{21}`$ in all cases met in calculations below the correct recipe for the calculation is
$$<X_2>_{(6)}=<X_2X_{21}>_p+<X_{21}>_{(6)}.$$
(34)
The decomposition into $`X_{21}`$ and $`X_{22}`$ is evidently non-unique, and $`X_{21}`$ may be choosen so as readily to calculate $`<X_{21}>_{(6)}`$ analytically.
In the method described above the $`<X_{21}>_{(6)}`$ is defined by the integral which is assumed to converge. Hence the $`Y_{21}`$ falls at $`q_1,q_2m`$ faster then uniform function having $`n=0`$. As the scale of the momenta, at which the original function $`Y_2`$ begins to fall rapidly, is $`q_1q_2m`$, it is natural to demand the same behavior from the function $`Y_{21}`$ as well; in other words, there exist effective ”cutoff” of the $`Y_{21}`$ at $`q_1q_2m`$. This ”cutoff” reminds regularizations in nonrelativistic effective theories; however, in the approach described here there is no need in artificial regulators such as large regulating mass or dimension different from 4. Similarly to what explained above, in the rest of the paper all functions looking like uniform (with $`n=0`$) at small momenta fall rapidly enough at large ones.
For calculation of ”tree” and ”one-loop” contributions it is necessary to find corrections to ”leading” terms, which are of first and second relative order in the $`\alpha `$. The method of finding of these corrections is rather obvious and is shown by the following simple illustration. Consider calculation of $`<X>`$ where $`X=X(𝐩,𝐩^{})`$ is sufficiently smooth function that doesn’t depend on $`\alpha `$, and expansion of $`X(𝐩,𝐩^{})X(0,0)`$ in $`𝐩,𝐩^{}`$ starts from third-order terms. Evidently
$$<X>=(A_0+A_1+A_2)|_{\sigma =0}^{\sigma =1},$$
$$A_0=|\phi _0(0)|^2X(0,0),$$
$$A_1=\phi _0(0)\frac{d^3p}{(2\pi )^3}\phi _0(𝐩)(X(𝐩,0)+X(0,𝐩)2X(0,0)),$$
(35)
$$A_2=\frac{d^3p}{(2\pi )^3}\frac{d^3p^{}}{(2\pi )^3}\phi _0^+(𝐩^{})(X(𝐩,𝐩^{})X(𝐩,0)X(0,𝐩^{})+X(0,0))\phi _0(𝐩).$$
(36)
It is easily seen that when calculating three first terms of expansion of $`<X>`$ in $`\alpha `$, in (35,36) $`\phi _0(𝐩)`$ may be replaced by $`\mathrm{\hspace{0.33em}8}\gamma ^{5/2}\pi ^{1/2}p^4`$; hence $`A_{0,1,2}`$ exactly have orders $`\alpha ^3,\alpha ^4,\alpha ^5`$ respectively.
Each ”one-loop” contribution has the form $`<X_1>_{(6)}`$ where
$$X_1(𝐩,𝐩^{},\gamma )=\frac{d^3q}{(2\pi )^3}\alpha ^2Y_1(𝐩,𝐩^{},𝐪,\gamma );$$
(37)
Let $`n`$ to be defined so that $`Y_1\delta ^{n2}`$ at $`pp^{}q\gamma \delta m`$ (if one writes $`X_1S_0V_0`$ as the loop integral, the $`n`$ is its ”divergency power”; $`n=0`$ imply that the region $`qm\alpha `$ contributes to order $`m\alpha ^6`$, and may cause terms $`m\alpha ^6\mathrm{ln}\alpha `$) to exist.
Let us write $`Y_1`$ for any given contribution as sum of ”soft” and ”hard” parts, $`Y_1=Y_{11}+Y_{12}`$, where $`Y_{11}`$ and $`Y_{12}`$ has $`n0`$ and $`n>0`$ respectively (similarly $`X_1=X_{11}+X_{12}`$). It may be done so that $`Y_{11}`$ at $`p,p^{},q,\gamma m`$ equals approximately a uniform function (the only exception is the contribution of graph shown at fig.2(a), which is treated in the section 4). It can be proved that under such choice of $`Y_{11}`$ contributions coming from $`Y_{12}`$ can be obtained, for all terms to treat, using formula like (35), which results in
$$<X_1>_{(6)}=<((X_1X_{11})(X_1X_{11})_0)S_0V_0+h.c.>_p+<X_{11}>_{(6)};$$
(38)
in (38) the notation $`()_0,`$ is used, defined so that $`(X)_0(𝐩,𝐩^{})X(𝐩,𝐩^{})|_{𝐩=𝐩^{}=0}`$. Evidently $`Y_{11}`$ can always be choosen so that second term in (38) is easily evaluated analytically.
”Tree” contributions to $`\mathrm{\Delta }\nu _{rec}`$ are calculated in the way just like one described above. There are only two ”tree” contributions to be found: $`<V_m>_{(6)}`$ and $`<V_{1c}>_{(6)}`$. Each of them is divided into ”soft” and ”hard” parts (the explicit way of this partition is described in the section 4), and then the ”hard” contributions are evaluated according to (36).
In the section 4 separation of ”soft” contributions from $`\mathrm{\Delta }\nu _{rec}`$ is done using formulas (34),(38),(36), that results in
$$<V_3>_{(6)}=<V_3W_{V3}>_p+E_{V3}^S,$$
(39)
$$<V_2+U_{V2}>_{(6)}=<V_2SV_m+h.c.>_p+<(V_2SV_cW_{V2})+h.c.>_p+E_{V2}^S,$$
(40)
$$<U_C>_{(6)}=<V_cSV_cSV_cW_C>_p+E_C^S,$$
(41)
$$<U_M>_{(6)}=<V_cSV_mSV_cW_{M1}>_p+<(V_mSV_cSV_cW_{M2})+h.c.>_p+E_M^S,$$
(42)
$$<U_{MM}>_{(6)}=<(V_mSV_mSV_cW_{MM})+h.c.>_p+E_{MM}^S,$$
(43)
$$<U_{MCM}>_{(6)}=<V_mSV_cSV_mW_{MCM}>_p+E_{MCM}^S,$$
(44)
$$<U_{MMM}>_{(6)}=<V_mSV_mSV_m>_p,$$
(45)
here $`W_{V3},W_{V2},W_C,W_{M1},W_{M2},W_{MM},W_{MCM}`$ are operators associated with the soft contributions, which have relatively simple form, $`E_{V3}^S,E_{V2}^S,E_M^S,E_{MM}^S,E_{MCM}^S,E_C^S`$ are analytically found values of ”soft” contributions to $`\mathrm{\Delta }\nu _{rec}`$; these values and operators are written explicitly in the section 4.
Remind that operators $`V_0+V_1,V_2,V_3`$ are given by sums of all the tree, one-loop, and two-loop irreducible noncovariant graphs respectively. Hence first terms in $`<>_p`$ in (39) - (45) are also given by sums of some noncovariant graphs (reducible ones, except $`V_3`$). The sum of them is equal to the sum of expressions associated with all two-loop graphs of the recoil type. The operator corresponding to this sum is referred to as $`W`$. For convenience of references, the following notations are also used: sum of all second terms in $`<>_p`$ in (39) - (45) is denoted as $`W_0`$, and sum of all the soft contributions as $`E^S`$, hence
$$\mathrm{\Delta }\nu _{rec}=<WW_0>_p+E^S;$$
also, notations $`W^{},W_0^{}`$ are used, defined according to
$$W,W_0(𝐩,𝐩^{},\gamma )|_{𝐩=𝐩^{}=\gamma =0}=\frac{d^3p_1}{(2\pi )^3}\frac{d^3p_2}{(2\pi )^3}W^{},W_0^{}(𝐩_1,𝐩_2).$$
A sum of contributions to $`W^{}`$, arising from all noncovariant graphs with the same topological structure (taking into account the difference of magnetic and Coulomb quanta), is equal to integrand associated with usual covariant Feynman graph (after the integration over zero components of the loop momenta). Hence the $`W^{}`$ can be determined as sum of contributions of all covariant two-loop graphs of the recoil type, and there is no need to know contributions to the $`W^{}`$ arising from separate noncovariant graphs.
In the above the gauge used for noncovariant and covariant graphs was undermined to be the same. However sum of contributions of all covariant graphs to the $`W^{}`$ is independent of gauge, and so in actual calculations any gauge may be choosen. This is quite natural; in terms of the NRQED (or the effective nonrelativistic Hamiltonian approach) $`<WW_0>_p`$ is equal, besides some overall factor, to contribution $`\alpha ^3`$ to the constant of pointlike fermion-fermion interaction (if an appropriate regularization is used), which is determined from scattering amplitudes on the mass shell.
4. THE ”SOFT” CONTRIBUTIONS.
In this section the general method of separation of ”soft” contributions, described above, is applied to various terms of (18) - (31) and to contributions of various graphs.
There are some preliminary remarks concerning the calculation technique employed. First, every term contributing to $`W`$ and $`W_0`$ evidently may be replaced by its average over positronium polarizations (or, in other words, over directions of the full spin). This simplifies the calculations considerably. For contributions to $`W`$ this average is found in program way, as $`1/3`$ of sum over polarizations. For contributions to $`W_0`$ the averaging is performed using the following formulas (to the r.h.s. of which terms independent of the full spin may be added):
$$(\stackrel{}{\sigma }_1𝐚)(\stackrel{}{\sigma }_2𝐛)\frac{1}{3}(\stackrel{}{\sigma }_1\stackrel{}{\sigma }_2)(\mathrm{𝐚𝐛}),$$
(46)
$$(\stackrel{}{\sigma }_1𝐚)(\stackrel{}{\sigma }_2𝐛)\left(\stackrel{}{\sigma }_1\stackrel{}{\sigma }_2\frac{1}{𝐪^2}(\stackrel{}{\sigma }_1𝐪)(\stackrel{}{\sigma }_2𝐪)\right)\frac{1}{3}\stackrel{}{\sigma }_1\stackrel{}{\sigma }_2(\mathrm{𝐚𝐛}+(\mathrm{𝐚𝐪})(\mathrm{𝐛𝐪})/𝐪^2),$$
(47)
$$\left(\stackrel{}{\sigma }_1\stackrel{}{\sigma }_2\frac{(\stackrel{}{\sigma }_1𝐪)(\stackrel{}{\sigma }_2𝐪)}{𝐪^2}\right)\left(\stackrel{}{\sigma }_1\stackrel{}{\sigma }_2\frac{(\stackrel{}{\sigma }_1𝐤)(\stackrel{}{\sigma }_2𝐤)}{𝐤^2}\right)\stackrel{}{\sigma }_1\stackrel{}{\sigma }_2\left(1+\frac{1}{3}\frac{(\mathrm{𝐪𝐤})^2}{𝐪^2𝐤^2}\right).$$
(48)
Second, let me describe calculation of contributions associated with retardation. These contributions have ”ultrasoft” parts, for which virtual photon momenta $`m\alpha ^2`$ are essential. These contributions are contained in $`<V_m>_{(6)},<V_2>_{(6)},<V_3>_{(6)}`$. Considering them, it is convenient to replace these operators by their spin-spin parts averaged over positronium polarizations, as described above. First nonrelativistic approximations (in the sense of the power expansion in momenta) for averaged contributions of graphs 2(a) and 2(b) to $`V_2`$ and $`V_3`$ are $`V_2^{ret}`$ and $`V_3^{ret}`$,
$$V_2^{ret}(𝐩,𝐤)=\frac{d^3q}{(2\pi )^3}\alpha ^2\frac{4}{3}\frac{\pi ^2}{m^2}\stackrel{}{\sigma }_1\stackrel{}{\sigma }_2\frac{q}{q^{\mathrm{\hspace{0.33em}2}}}Y_1^{ret}(𝐩,𝐤,𝐪),$$
(49)
$$V_3^{ret}(𝐩,𝐤)=\frac{d^3q_1}{(2\pi )^3}\frac{d^3q_2}{(2\pi )^3}\alpha ^3\frac{16}{3}\frac{\pi ^3}{m^2}\stackrel{}{\sigma }_1\stackrel{}{\sigma }_2\frac{q}{q_1^2q_2^2}Y_2^{ret}(𝐩,𝐤,𝐪_1,𝐪_2),$$
(50)
where
$$Y_1^{ret}(𝐩,𝐤,𝐪)=\frac{1}{q+(f_p+f_{k_1})/2}\frac{1}{q+(f_k+f_{p_1})/2},$$
(51)
$$Y_2^{ret}(𝐩,𝐤,𝐪_1,𝐪_2)=\frac{1}{(q+(f_p+f_{k_2})/2)}\frac{1}{(q+(f_{p_1}+f_{k_1})/2)}\frac{1}{(q+(f_{p_2}+f_k)/2)}$$
(52)
(notations for the momenta are shown at fig. 2); just the same contributions arise due to graphs obtained from graphs 2(a) and 2(b) by the time reversal.
Spin-spin part of $`V_m`$, averaged over positronium polarizations, is denoted below as $`V_s`$. First nonrelativistic approximation to $`V_s(𝐩,𝐤)`$ is
$$V_m^{ret}(𝐩,𝐤)=V_{m0}\frac{q}{q+(f_p+f_k)/2},$$
(53)
where
$$V_{m0}=V_s|_{𝐩=𝐤=\gamma =0}=\frac{2}{3}\frac{\alpha \pi }{m^2}\stackrel{}{\sigma }_1\stackrel{}{\sigma }_2,𝐪𝐤𝐩.$$
It is convenient to consider $`<V_m^{ret}>_{(6)},<V_2^{ret}>_{(6)},<V_3^{ret}>_{(6)}`$ together and to calculate their overall contribution to $`\mathrm{\Delta }\nu _{rec}`$; this way of the calculation was employed, in particular, in and . The way of the calculation, used in the present paper is, in some sense, even shorter. Each of operators $`V_m^{ret},V_2^{ret},V_3^{ret}`$ is divided into two parts:
$$V_m^{ret}(𝐩,𝐤)=V_{m0}^{ret}(𝐩,𝐤)+V_{m0}^{}(𝐩,𝐤),$$
$$V_2^{ret}(𝐩,𝐤)=V_{20}^{ret}(𝐩,𝐤)+V_{20}^{}(𝐩,𝐤),$$
$$V_3^{ret}(𝐩,𝐤)=V_{30}^{ret}(𝐩,𝐤)+V_{30}^{}(𝐩,𝐤),$$
where
$$V_{m0}^{ret}(𝐩,𝐤)=\frac{1}{2}V_{m0}\left[\frac{q}{q+f_p}+\frac{q}{q+f_k}\right],$$
(54)
and $`V_{20}^{ret}`$, $`V_{30}^{ret}`$ are obtained from $`V_2^{ret}`$, $`V_3^{ret}`$ by replacing $`Y_1^{ret},Y_2^{ret}Y_{10}^{ret},Y_{20}^{ret}`$ where
$$Y_{10}^{ret}(𝐩,𝐤,𝐪)=\frac{1}{2}\left(\frac{1}{(q+f_k)(q+f_{k_1})}+\frac{1}{(q+f_p)(q+f_{p_1})}\right),$$
(55)
$$Y_{20}^{ret}(𝐩,𝐤,𝐪_1,𝐪_2)=\frac{1}{2}\left(\frac{1}{(q+f_p)(q+f_{p_1})(q+f_{p_2})}+\frac{1}{(q+f_k)(q+f_{k_1})(q+f_{k_2})}\right).$$
(56)
In the region where all the momenta are $`m\alpha `$
$$V_m^{ret}V_{m0}^{ret},Y_1^{ret}Y_{10}^{ret},Y_2^{ret}Y_{20}^{ret}.$$
If one writes $`V_{20}^{},V_{30}^{}`$ as integrals like (49,50), the integrands in them have powers, being expanded in the momenta, greater by 2 than those of $`Y_1^{ret},Y_2^{ret}`$; also $`V_{m0}^{}`$ has power, being expanded in the momenta, greater by 2 than the power of $`V_m^{ret}`$. Due to this reason $`<V_{20}^{}>_{(6)},<V_{30}^{}>_{(6)}`$ consist of ”hard” contributions only. As for $`<V_{m0}^{}>_{(6)},`$ it does contain soft contributions. However it may be easily calculated in just the way other soft contributions to $`<V_m>_{(6)}`$ are calculated, and is actually evaluated combined with them in sec.4.3.
Now turn to contributions of $`V_{m0}^{ret},V_{20}^{ret},V_{30}^{ret}`$. It is easy to see that their full contribution is equal to 0:
$$<V_{m0}^{ret}>_{(6)}+2<V_{20}^{ret}>_{(6)}+2<V_{30}^{ret}>_{(6)}=0.$$
(57)
This may be shown by a short direct calculation; however (57) is evident immediately due to the following reason. Consider calculation of the recoil contribution $`\alpha ^6m^2/M`$ to ground state hyperfine splitting of the ”hydrogen”, i.e. particle of mass $`m/2`$ moving around particle of large mass $`M`$. In this calculation $`V_m^{ret}(𝐩,𝐤)`$, for instance, must be replaced by $`2m/MV_{m0}^{ret}(𝐩,𝐤)`$; entirely $`<V_m^{ret}+2V_2^{ret}+2V_3^{ret}>_{(6)}`$ is replaced by
$`2m/M<V_{m0}^{ret}+2V_{20}^{ret}+2V_{30}^{ret}>_{(6)}`$, as may be easily seen. On the other hand, it is well known that when calculating hyperfine splitting in first order in $`m/M`$ the magnetic interaction may be regarded as instant, which immediately leads to
$$<V_{m0}^{ret}+2V_{20}^{ret}+2V_{30}^{ret}>_{(6)}=<V_{m0}>_{(6)}=0.$$
(58)
The last point to discuss is calculation of $`<U_L>_{(6)}`$, that is essentially the contribution of graphs containing Coulomb ladder with three or more loops. This contribution to $`\mathrm{\Delta }\nu _{rec}`$ arises completely, with the accuracy required, from the region where all the momenta in the loops are of order $`m\alpha `$, and can be easily evaluated analytically. There are two kinds of such contributions, namely, graphs having one or two magnetic quanta. The contribution of graphs having two magnetic quanta is
$$E_{LMM}^S<V_mLV_m>_{(6)}$$
and does not depend, in the order required, on details of formalism used; for the first time it has been found in . Method of its evaluation, used in the present paper, coincides mainly with the method used in (and involves evaluation in coordinate representation); the result obtained agrees with those found earlier and equals
$$E_{LMM}^S=\left(\frac{791}{864}\frac{\pi ^2}{18}\right)m\alpha ^60.3672m\alpha ^6.$$
(59)
The contribution of graphs having one magnetic photon does depend on details of formalism used. The value calculated in the present paper is
$$E_{LM}^S<(V_mL(V_{1c}+\mathrm{\Delta }T))+h.c.>_{(6)}.$$
This value was evaluated in two ways: in coordinate space (by method close to that of ), and by immediate integration over the momenta, using explicit form of $`L(𝐩,𝐤)`$, quoted, for instance, in :
$$L(𝐩,𝐤)=\frac{4\pi m^4\alpha ^3}{f_p^2f_k^2}\left(\frac{5}{2}4\frac{\gamma ^2}{f_p}4\frac{\gamma ^2}{f_k}+\frac{1}{2}\mathrm{ln}A+\frac{2A1}{\sqrt{4A1}}\mathrm{arctan}\sqrt{4A1}\right),$$
(60)
where
$$A=\frac{f_pf_k}{4\gamma ^2q^2}.$$
The results obtained by both methods are the same and equal to
$$E_{LM}^S=\frac{1}{64}m\alpha ^60.01562m\alpha ^6.$$
(61)
Finally, addition (59) and (61) together results in
$$<U_L>_{(6)}0.3728m\alpha ^6.$$
4.1. Irreducible Graphs.
For calculation of $`<V_3>_{(6)}`$ it is necessary first to separate from $`V_3`$ the term $`V_{30}`$ which is an integral having the ”divergency power” $`n=0`$, and to use formula (34) where $`X_{21}`$ should be set equal to $`V_{30}`$; the result has the form (39) with
$$W_{V3}=V_{30},E_{V3}^S=<V_{30}>_{(6)}.$$
There exist only three essentially different graphs contributing to $`V_{30}`$: these are graphs depicted at fig. 3(a),3(b) and 2(b). However contribution to $`V_{30}`$, due to sum of graphs 3(a) and 3(b), may be set equal to zero. The matter is that this sum in the region of small loop momenta reduces effectively to the graph 3(c), in which effective two-photon ”seagull” vertex does not depend on the spins. This vertex appears as sum of subgraphs 3(d) and 3(e) and corresponds to term $`\alpha 𝐀^2/(2m)`$ (where $`𝐀`$ is the vector potential) in the nonrelativistic Hamiltonian of particle in magnetic field. Thus $`V_{30}`$ arises from the only graph 2(b) and may be choosen to be $`2V_{30}^{ret}`$ which results in $`E_{V3}^S=2<V_{30}^{ret}>_{(6)}`$; in fact there is no need to evaluate this value, as mentioned above.
Similarly $`<V_2>_{(6)}`$ is calculated using (38) with $`X_{11}=V_{20}`$ where $`V_{20}`$ is part of $`V_2`$, having $`n=0`$. Noting that $`<U_{V2}>_{(6)}=<U_{V2}>_p,`$ as is easily seen, one finds (40) where
$$W_{V2}=(V_{20}+(V_2V_{20})_0)S_0V_0,E_{V2}^S=<V_{20}>_{(6)}.$$
Consider contributions to $`V_{20}`$. Graphs having two Coulomb quanta do not contribute to $`V_{20}`$. Among graphs having two magnetic quanta there are those with $`n=1`$ (graphs depicted in fig. 3(f),(g)), and several graphs with $`n=0`$. However their contributions to $`V_{20}`$ cancel each other, just as described above for $`V_{30}`$; for instance, sum of graphs shown in fig. 3(f) and 3(g) reduces effectively in the region of small momenta to the graph 3(h), in which effective seagull vertex does not depend on spins. Graphs having one magnetic and one Coulomb quanta, contributing to $`V_2`$, are divided into 4 sets belonging each to one of covariant graphs (A),(B), shown in fig.4, and the graphs obtained by the time reversal. It is convenient to consider graphs including electron-positron pairs and those without pairs separately.
Contributions of graphs including pairs to $`V_{20}`$ may be choosen to be
$$V_{20}^A(𝐩,𝐩^{})=V_{20}^B=\mathrm{}.=\frac{d^3q}{(2\pi )^3}\frac{2\alpha ^2\pi ^2}{3m^3}\stackrel{}{\sigma }_1\stackrel{}{\sigma }_2\frac{1}{q^2q^2}((q^2+q^2)R_qk^2R_k),$$
(62)
(the momenta notations are shown in fig. 4, $`𝐤𝐩^{}𝐩`$, and $`R_km^2/(m^2+k^2)`$ for an arbitrary $`𝐤`$ ). It is easy to find their overall contribution to $`E_{V2}^S`$, that equals
$$E_P^S=\frac{2}{3}m\alpha ^6\mathrm{ln}\alpha .$$
The only graphs without pairs, contributing to $`V_2`$, are the graph shown in fig. 2(a), discussed above, and the similar graph obtained from 2(a) by the time reversal. Their contribution to $`V_{20}`$ is choosen to be $`2V_{20}^{ret}`$. So
$$E_{V2}^S=E_P^S+2<V_{20}^{ret}>_{(6)}.$$
(63)
The second term in (63) need not in fact be evaluated, as explained above.
4.2. Graphs with one magnetic photon.
Consider calculation of $`<U_M>_{(6)},`$ i.e. contribution to $`\mathrm{\Delta }\nu _{rec}`$ due to reducible graphs with one magnetic photon, except corresponding contribution to $`<U_L>_{(6)}`$, considered above.
As discussed above, it is convenient to calculate $`<V_s>_{(6)}`$ instead of $`<V_m>_{(6)}`$, where $`V_s`$ is spin-spin part of $`V_m`$, averaged over positronium polarizations. Let us write $`V_s`$ as an ”expansion in powers of momenta”:
$$V_s=V_{m0}^{ret}+V_{m1}+V_{m2},$$
(64)
where $`V_{m0}^{ret}`$ is defined in (54), $`V_{m1}(𝐩,𝐤,\gamma )V_sV_{m0}^{ret}`$ at $`pk\gamma m`$, and $`V_{m1}`$ at $`p,k,\gamma m`$ equals approximately a uniform function of second power. It can be shown that $`<V_{m2}>_{(6)}`$ can be evaluated using (36) with $`X=V_{m2}`$. Actually there is no need to evaluate $`<V_{m0}^{ret}>_{(6)}`$, as mentioned above. The value $`<V_{m1}>_{(6)}`$ can be easily found in analytic form, if one uses appropriate choice for $`V_{m1}`$. It is convenient to choose
$$V_{m1}(𝐩,𝐤)=\frac{1}{2m^2}V_{m0}\left[\frac{p^2k^2}{q^2}(R_p+R_k)\frac{1}{q^2}(p^4R_p+k^4R_k)(p^2R_p+k^2R_k)\right],$$
(65)
where $`(𝐪𝐩𝐤).`$
In the lowest order in momenta $`V_mV_{mb}`$ where $`V_{mb}`$ is part of the Breit Hamiltonian due to magnetic photon exchange. The value $`<U_{M2}>_{(6)}`$ is calculated using (38) with
$`X_{11}=V_{mb}(S_0V_{11}^{}+S_{11}V_0)+h.c.,`$ where $`S_{11},V_{11}^{}`$ are ”first terms of expansion” of $`SS_0`$ and $`V_cV_0`$ in the momenta (namely, $`V_{11}^{}(𝐩,𝐤)V_cV_0,S_{11}(𝐩)SS_0`$
at $`pk\gamma m`$, and $`S_{11},V_{11}^{}`$ at $`p,km`$ equal approximately to uniform functions of zero power).
The value $`<U_{M3}>_{(6)}`$ is calculated using (34) with
$`X_{21}=V_{mb}S_0V_0(S_0V_{12}^{}+S_{12}V_0)+h.c.`$, where $`S_{12}(𝐩),V_{12}^{}(𝐩,𝐤)S_{11}(𝐩),V_{11}^{}(𝐩,𝐤)`$
at $`p,km`$ (but generally $`S_{12},V_{12}^{}S_{11},V_{11}^{};`$ the explicit way of the decomposition of $`SS_0`$ and $`V_cV_0`$ is determined by the purpose of convenience of subsequent calculations).
It is easy to see that in $`X_{11},X_{21}`$ defined in such manner it is enough to take into account only spin-independent parts of $`V_{11}^{},V_{12}^{}`$ (denoted as $`V_{11},V_{12}`$ respectively) and spin-spin part of $`V_{mb}`$, that can be replaced by its average over positronium polarizations, i.e. $`V_{m0}`$. So $`<U_{M2}>_{(6)}`$ and $`<U_{M3}>_{(6)}`$ can be evaluated using
$`X_{11}=V_{m0}(S_0V_{11}+S_{11}V_0)+h.c.,`$ $`X_{21}=V_{m0}S_0V_0(S_0V_{12}+S_{12}V_0)+h.c..`$
It is convenient to choose
$`V_{12}(𝐩,𝐤)=\alpha \pi /m^2R_p=\alpha \pi /m^2(1+p^2/m^2)^1`$;
$`V_{11}`$ is defined so that $`V_{11}(𝐩,𝐤)=V_{11}(𝐩,0);`$ under this condition an explicit form of $`V_{11}`$ does not enter final expressions.
It is easy to find that $`S_{11},S_{12}`$ may be set to
$$S_{1i}(p)=S_{1i}(p)|_{\gamma =0}+\frac{1}{2}\frac{\gamma ^2}{f_p}\frac{1}{4}\frac{\gamma ^4}{f_p^2},(i=1,2);$$
(66)
first term in (66) is choosen to be $`(SS_0)|_{(\gamma =0)}`$ for $`S_{11}`$ and $`R_p/4`$ for $`S_{12}`$.
After a simple calculation $`<U_M>_{(6)}`$ is found to have the form (42) with
$$E_M^S=m\alpha ^6\left(\frac{1}{3}\mathrm{ln}\alpha +\frac{53}{192}\right)+<V_{m0}^{ret}>_{(6)},$$
$$W_{M1}=((V_cSV_m)_0S_0V_0+h.c.)+V_0S_0(V_{m0}^{ret}+V_{m1})S_0V_0((V_0S_0(V_{m0}^{ret}+V_{m1}))_0S_0V_0+h.c.),$$
$$W_{M2}=V_{m0}S_0V_0(S_0V_{12}+S_{12}V_0)+(V_mSV_0+(V_mS(V_cV_0))_0)S_0V_0$$
(when deriving $`W_{M2}`$ the term
$$((V_mV_{m0})S_{11}V_0((V_mV_{m0})S_{11}V_0)_0)S_0V_0,$$
was added to it, which causes no additional contribution of order $`m\alpha ^6`$).
4.3. Coulomb graphs.
The Coulomb contribution due to reducible graphs to $`\mathrm{\Delta }\nu _{rec}`$ is equal to $`<U_C>_{(6)}`$. This value was found in the style described above: $`<V_{1c}>_{(6)},<U_{C2}>_{(6)},<U_{C3}>_{(6)}`$ (i.e. separate terms of $`<U_C>_{(6)}`$) were evaluated using (36,38,34) respectively. The calculation is simplified by the fact that in lowest (i.e. zeroth) order in momenta $`V_{1c}`$ contains only spin-independent terms, and besides that
$$V_{1c}^S(𝐩,𝐤)|_{𝐩=0}=V_{1c}^S(𝐩,𝐤)|_{𝐤=0}=0,$$
where $`V_{1c}^S`$ is spin-spin part of $`V_{1c}`$.
It is easy to see that $`<U_C>_{(6)}`$ is determined by the formula (41), in which
$$W_C=V_0S_0V_{1c}^{}S_0V_0,E_C^S=<V_{1c}^{}>_{(6)};$$
(67)
here $`V_{1c}^{}`$ is first nonvanishing term of ”expansion” of $`V_{1c}^S`$ in momenta (exactly, $`V_{1c}^{}(𝐩,𝐤)V_{1c}^S(𝐩,𝐤)`$ at $`pkm`$, and $`V_{1c}^{}`$ at $`p,km`$ equals approximately a uniform function of second power).
In (67) one can replace $`V_{1c}^{}`$ by its average over positronium polarizations (denoted as $`V_{1c}^{\prime \prime }`$); it is convenient to set
$$V_{1c}^{\prime \prime }(𝐩,𝐤)=\frac{1}{24}\stackrel{}{\sigma }_1\stackrel{}{\sigma }_2\frac{\alpha \pi }{m^2}(\frac{1}{2}(\frac{p^4}{q^2}p^2\mathrm{𝐩𝐤})R_p+(𝐩𝐤)+\frac{p^2k^2}{q^2}R_p);$$
using it, one easily obtains
$$E_C^S=\frac{1}{48}m\alpha ^6\left(\mathrm{ln}\alpha +\frac{1}{4}\right).$$
4.4. Graphs with two magnetic quanta.
Turn to calculation of $`<U_{MM2}+U_{MM3}+U_{MCM}>_{(6)}`$, i.e. contribution of reducible graphs with two magnetic quanta, except corresponding contribution to $`<U_L>_{(6)}`$, considered above. The value $`<U_{MM2}>_{(6)}`$ is found using (38) with $`X_{11}=V_{mb}S_0RV_{mb}`$, where $`R(𝐩,𝐩^{})=(2\pi )^3\delta ^3(𝐩𝐩^{})R_p`$. Similarly, the value $`<U_{MCM}>_{(6)}`$ is found using (34) with $`X_{21}=V_{mb}S_0V_0S_0RV_{mb}.`$ It is easy to see that $`<U_{MM3}>_{(6)}=<U_{MM3}>_p`$. The results have form (43),(44) with
$$W_{MM}=V_{mb}S_0RV_{mb}S_0V_0+(V_mSV_mV_{mb}S_0RV_{mb})_0S_0V_0,$$
$$W_{MCM}=V_{mb}S_0V_0S_0RV_{mb},$$
$$E_{MM}^S=<V_{mb}S_0RV_{mb}>_{(6)},E_{MCM}^S=<V_{mb}S_0V_0S_0RV_{mb}>_{(6)}.$$
Spin-orbital part of $`V_{mb}`$ does not contribute to $`E_{MM}^S`$ and $`E_{MCM}^S`$. Spin-spin and spin-independent contributions to $`V_{mb}`$ are denoted as $`V_{sb}`$ and $`V_{lb}`$ respectively; they are equal to
$$V_{sb}(𝐩,𝐩^{})=\frac{\alpha \pi }{m^2}\left(\stackrel{}{\sigma }_1\stackrel{}{\sigma }_2\frac{1}{𝐪^2}(\stackrel{}{\sigma }_1𝐪)(\stackrel{}{\sigma }_2𝐪)\right),$$
$$V_{lb}(𝐩,𝐩^{})=4\frac{\alpha \pi }{m^2}\left(\frac{(\mathrm{𝐪𝐩})^2}{q^4}+\frac{p^2}{q^2}\right),𝐪𝐩^{}𝐩.$$
As the values to calculate are average values over S-state, $`V_{lb}`$ may be replaced by any function $`V_{lb}^{}`$ such that
$$V_{lb}(𝐩,𝐩^{})𝑑\vartheta _{\mathrm{𝐩𝐩}^{}}=V_{lb}^{}(𝐩,𝐩^{})𝑑\vartheta _{\mathrm{𝐩𝐩}^{}};$$
It is convenient to choose
$$V_{lb}^{}(𝐩,𝐩^{})=4\frac{\alpha \pi }{m^2}\left(\frac{1}{2}+\frac{1}{2}\left(\frac{p^2}{q^2}+\frac{p^2}{q^2}\right)\right).$$
For separate parts of $`E_{MM}^S,E_{MCM}^S`$ one obtains
$$<V_{lb}^{}S_0RV_{sb}+V_{sb}S_0RV_{lb}^{}>_{(6)}=\left(\frac{1}{3}\mathrm{ln}\alpha \frac{1}{6}\right)m\alpha ^6,$$
$$<V_{sb}S_0RV_{sb}>_{(6)}=\left(\frac{1}{24}\mathrm{ln}\alpha \frac{7}{96}\right)m\alpha ^6,$$
$$<V_{lb}^{}S_0V_0S_0RV_{sb}+V_{sb}S_0V_0S_0RV_{lb}^{}>_{(6)}=\left(\frac{1}{6}+\frac{\pi ^2}{18}\right)m\alpha ^6,$$
$$<V_{sb}S_0V_0S_0RV_{sb}>_{(6)}=\left(\frac{5}{48}\mathrm{ln}\alpha +\frac{1}{96}\right)m\alpha ^6,$$
which results in
$$E_{MM}^S=\left(\frac{3}{8}\mathrm{ln}\alpha \frac{23}{96}\right)m\alpha ^6,$$
$$E_{MCM}^S=\left(\frac{5}{48}\mathrm{ln}\alpha \frac{5}{32}+\frac{\pi ^2}{18}\right)m\alpha ^6.$$
Now the calculation of the ”soft” contributions is completed. Their sum is
$$E^S=\left(\frac{1}{6}\mathrm{ln}\alpha +\frac{1393}{1728}\right)m\alpha ^6\left(\frac{1}{6}\mathrm{ln}\alpha +0.8061\right)m\alpha ^6.$$
To check the formalism and the calculation technique used, they were applied to calculation of the recoil contribution $`\alpha ^6m^2/M`$ to the ground state hyperfine splitting of the hydrogen. It is well known that for this particular purpose the two-particle problem reduces effectively to the problem of motion in external field, and the value desired may be easily obtained by means of coordinate-space calculations, using known solutions of the Dirac equation in the Coulomb field. The correction to the hyperfine splitting thus obtained is equal to $`4\alpha ^6m^2/M`$. The same result was found using momenta representation by the method described in the present paper. To simplify corresponding calculations, diagram technique rules were modified, because in the relevant order one can consider the magnetic field as permanent, and so the magnetic photon exchange may be assumed to be instantaneous. All the ”hard” contributions cancel each other; the soft contributions of the desired order arise only from graphs having one magnetic photon, and their calculation coincides mainly with the calculation of corresponding contributions for positronium.
5. THE ”HARD” CONTRIBUTIONS.
For the calculation of the ”hard” contribution to $`\mathrm{\Delta }\nu _{rec}`$, i.e. the value $`<WW_0>_p`$, covariant graphs are used, as explained in section 3. In covariant gauge there are only 4 essentially different two-loop ”recoil” graphs, whereas in the Coulomb gauge there are 24 ones (regarding magnetic and Coulomb quanta as depicted by different lines); here term ”essentially different” mean that these graphs cannot be obtained from each other by transposition of two particles and/or time reversal.
The hard contribution to $`\mathrm{\Delta }\nu _{rec}`$ was evaluated in two ways. The first way is to use the Coulomb gauge and to evaluate contributions of the 24 graphs separately. The second one is to calculate $`<WW_0>_p`$ as a whole using the Feynman gauge; in this case the $`W^{}`$ is sum of contributions of only 4 graphs. Besides that, these contributions have simpler form than those of the Coulomb gauge (strictly speaking, the ”simpler form” refers to simpler form of the program calculating these contributions).
The result for the ”hard” contribution is $`0.424(6)m\alpha ^6`$ for the ”separate” and $`0.426(6)m\alpha ^6`$ for the ”united” calculation. Adding the ”soft” contribution one obtains
$$\mathrm{\Delta }\nu _{rec}=m\alpha ^6\left(\frac{1}{6}\mathrm{ln}\alpha +0.382(6)\right)$$
(68)
and
$$\mathrm{\Delta }\nu _{rec}=m\alpha ^6\left(\frac{1}{6}\mathrm{ln}\alpha +0.380(6)\right)$$
(69)
for the ”separate” and the ”united” calculation methods respectively. These results are in perfect agreement with those found in .
The error estimation quoted in (68,69) is arbitrary in some extent, being subject of handwork; the reason of this is that method used in the program for calculation of $`W^{}`$ is numerically unstable, and the error quoted arises from this unstability (and not from the error of numerical integration). The direct way of the error estimation is to compare the results of ”separate” and ”united” calculations; difference between them is indeed within the estimation quoted above, as well as their differences from the results of works .
Results for contributions of separate graphs are quoted in tables 1a,1b,2 together with the revised results of paper , which agree with the results of the present paper. Errors are not given in the tables because they are not well defined; in fact, error for any contribution is $`1\%`$.
Notations for the graphs are following. If one makes no difference between magnetic and Coulomb lines, there exist 4 essentially different two-loop recoil graphs, denoted as 1,2,3,4 at fig. 5. To denote magnetic and Coulomb lines, indices ”M” and ”C” are written in the sequence the lines are attached to the lower fermion line from left to right. Values quoted in the tables refer to ”essentially different” graphs only, including in each quoted value contributions of all graphs of the similar form.
Having been averaged over the full spin direction, $`W^{}`$ depends on three variables; convenient choice for these variables are absolute values of the loop momenta, referred as $`p_1,p_2,q`$. For the ”united” method of the calculation it is convenient to symmetrize $`W^{}`$ over permutations of $`p_1,p_2,q`$, and to take into account only contributions of ”essentially different” graphs, taking some of them multiplied by 2 or 4 because of transposition of the two particles and/or time reversal.
Let us describe briefly the checking of the calculation of the ”hard” contributions. The first method of the checking is to compare $`<WW_0>_p`$ found in the two ways, the ”separate” and the ”united” (and intermediate comparison of $`W^{}(p_1,p_2,q)`$ obtained in the two ways, at several points in $`p_1,p_2,q`$).
Second, contributions to $`W_0^{}`$ were found by a relatively short analytic calculation (except contributions associated with sum of graphs 1-CMM and 2-CMM), whereas contributions to $`W^{}`$ are results of the numerical computing. An error in $`W^{}`$ or $`W_0^{}`$ leads usually to divergency of the numerical integration; clearly search for such errors is more convenient for the ”separate” calculation.
Third, the results of the ”separate” calculation are compared with the results of in which $`\mathrm{\Delta }\nu _{rec}`$ was calculated using Bethe-Salpeter formalism and the Coulomb gauge. The correction to the energy can be written as an expansion in irreducible covariant graphs (except the tree graph with Coulomb photon). Numerical value of any term of this expansion is clearly independent of formalism used, e.g. the formalism of the present paper or that of . So most of contributions found in the present paper should be equal to corresponding contributions found in . The other terms can be divided into groups, and total contribution of each group should also be equal to corresponding contribution of . Besides that, to order under the consideration the contribution of graph 1-MCM should be equal separately to corresponding contribution of .
Results for ”hard” contributions of separate graphs are given in tables 1a,1b. Part of them (namely, those quoted in table 1b)) are not uniquely defined as they depend on the choice of $`W_0`$. On the other hand, contribution of every term of the expansion in irreducible graphs, mentioned above, does not depend on formalism used, e.g. on the choice of $`W_0`$. Some terms of this expansion are trivial, i.e. they contain only ”hard” parts, and every of them consist of contribution of just one covariant graph; corresponding results are quoted in table 1a. The other (”non-trivial”) terms of this expansion involve also ”soft” contributions, and their ”hard” parts may consist of contributions of several covariant graphs; their values are given in table 2. Tables 1a,2 also contain corresponding results of . All the results of are in good agreement with those of the present work.
Contributions of graphs 3-MMC and 4-MMC were computed together, because these graphs contain noncovariant graphs shown in fig. 3(a) and 3(b), contributions of which to $`V_{30}`$ cancel each other, which simplified the calculation. Contributions of graphs 1-CMM and 2-CMM were also evaluated together due to the similar reason.
In table 2 the value $`E_{1CCC}^H`$, for instance, denotes the ”hard” part of the contribution of graph 1-CCC, i.e. one of values quoted in table 1b. The value $`E_M`$ stands for
$$E_{3CCM}^H+E_{2CMC}^H+E_{1MCC}^H+E_{1CMC}^H+E_{LM}^S+E_M^S+E_P^S,$$
and the $`E_M^{}`$ denotes
$$(ct0+tc0)+ctc0+(ctx+tcx)+(cct0+tcc0)+(ctcy+ccty)+tccz+\mathrm{\Delta }E_{MP}^{hfs}(\delta K_0T+T\delta K_0)+\mathrm{\Delta }E_d^{hfs}$$
(using notations of paper ).
Table 1a.
”Trivial” (i.e. containing only the ”hard” part) contributions to $`\mathrm{\Delta }\nu _{rec}`$
(in units $`m\alpha ^6`$); the results of are also quoted.
| Graph | contribution to $`\mathrm{\Delta }\nu _{rec}`$ | | contribution to $`\mathrm{\Delta }\nu _{rec}`$ |
| --- | --- | --- | --- |
| | | (notations of ) | (obtained in ) |
| 4-CCC | -0.0039 | cccx | -0.0039 |
| 4-CMC | 0.0042 | ctcx | 0.0043 |
| 4-MCC | -0.0486 | cctx + tccx | -0.0489 |
| 4-MCM | -0.0230 | tctx | -0.0230 |
| 1-MMM | 0.0694 | ttt0 | 0.0694 |
| 2-MMM | -0.0011 | ttty | -0.0011 |
| 3-MMM | -0.0012 | tttz | -0.0011 |
| 4-MMM | 0.0042 | tttx | 0.0041 |
| 3-CCC | 0.0064 | cccz | 0.0063 |
| 3-MCC | 0.0268 | ctcz+cctz | 0.0283 |
| 3-CMM | 0.0530 | ttcz+tctz | 0.0534 |
| 2-CCC | -0.0184 | cccy+ccx | -0.0186 |
| 2-MCC | -0.0661 | tccy | -0.0681 |
| 2-MCM | 0.0552 | tcty+ttcy | 0.0558 |
Table 1b.
”Hard” parts of ”non-trivial” contributions to $`\mathrm{\Delta }\nu _{rec}`$ (in units $`m\alpha ^6`$). Graph 1-CCC -0.0094 1-MCM -0.0745 1-CMM + 2-CMM -0.5112 4-MMC + 3-MMC -0.0209 3-CCM -0.0104 2-CMC -0.0245 1-MCC 0.1795 1-CMC -0.0092
Table 2.
”Non-trivial” contributions to $`\mathrm{\Delta }\nu _{rec}`$ (in units $`m\alpha ^6`$),
compared with the results of .
| value quoted | value quoted | coefficient | | constant |
| --- | --- | --- | --- | --- |
| | (notations of ) | at $`\mathrm{ln}\alpha `$ | constant | (obtained in ) |
| $`E_{4MMC+3MMC}^H`$ | (cttx+ttcx)+cttz | 0 | -0.0209 | -0.0209 |
| $`E_{1CCC}^H+E_C^S`$ | cc0+ccc0 | $`1/48`$ | -0.0146 | -0.0148 |
| $`E_{1MCM}^H+E_{MCM}^S`$ | tct0 | $`5/48`$ | 0.3176 | -0.3138 |
| $`E_{1CMM+2CMM}^H+E_{MM}^S`$ | ctt0+ttc0+ctty + | $`3/8`$ | -0.751 | -0.749 |
| | \+ tt0+ttx | | | |
| $`E_M`$ | $`E_M^{}`$ | 1/3 | 0.427 | 0.423 |
6. SUMMARY AND DISCUSSION.
In the present paper the ”recoil” contribution of order $`m\alpha ^6`$ to the hyperfine splitting of positronium ground state has been found. The calculation was performed using noncovariant formulation of QED perturbation theory. The result equals
$$\mathrm{\Delta }\nu _{rec}=m\alpha ^6\left(\frac{1}{6}\mathrm{ln}\alpha +0.381(6)\right)$$
and is in perfect agreement with the results of . Let us note that methods used in and in the present work are similar, whereas that of is completely different. So $`\mathrm{\Delta }\nu _{rec}`$ may be considered to be firmly established. However agreement with the experimental data is poor. Combining nonlogarithmic part of the recoil contribution with all other contributions of the relevant order (they also may be considered to be reliable), one obtains the complete nonlogarithmic contribution $`m\alpha ^6`$ to $`\mathrm{\Delta }\nu `$:
$$m\alpha ^6(0.3928)=7.33MHz$$
(70)
((70) was obtained using for $`\mathrm{\Delta }\nu _{rec}`$ the result of , as the most accurate). The complete theoretical result up to the order $`m\alpha ^6`$ is sum of (70) and (3), and equals
$$\mathrm{\Delta }\nu _{th}=\mathrm{203\; 392.96}MHz.$$
It differs from the experimental result (2) by 5 standard deviations.
The leading term of the next order in $`\alpha `$, i.e. of order $`m\alpha ^7\mathrm{ln}^2\alpha `$, was found in and equals
$$\frac{7}{8\pi }m\alpha ^7\mathrm{ln}^2\alpha =0.92MHz.$$
Taking this correction into account reduces the difference between experimental and theoretical results to 4 standard deviations. One cannot reject possibility that this difference may be explained by contributions $`m\alpha ^7\mathrm{ln}\alpha ,m\alpha ^7)`$).
I am grateful to I.B.Khriplovich, A.I.Milstein, and A.S. Yelkhovsky for useful conversations and interest to this work. I am also grateful to A.S.Yelkhovsky for informing me about the result of prior to publication, and G.S.Adkins for useful communication. I am also deeply grateful to D.Yu.Ivanov for his co-operation in the early stage of the work.
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# Renormalization of quark axial current in the chiral potential model
## Abstract
Non-conserved composite operators like the quark axial current have divergent matrix elements therefore must be renormalized. We explore how this can be done in quark model calculations where the systematic procedure of dimensional regularization and minimal subtraction is not applicable. We propose a most natural and convenient regularization scheme of cutting the intermediate quark states over which we sum in loop diagram calculations at a certain energy. We show that this scheme works perfectly for the quark axial current and we obtain the quark spin contribution to the proton spin: $`\mathrm{\Delta }_u=0.82`$, $`\mathrm{\Delta }_d=0.43`$, $`\mathrm{\Delta }_s=0.10`$, which is in excellent agreement with experiments.
PACS numbers: 12.39.Ki, 12.39.Fe, 12.39.Pn, 13.88.+e
The expression nucleon spin “crises” denotes the findings of the European Muon Collaboration (EMC) that a small proportion of the nucleon spin is carried by the quark spin and that strange quark polarizes significantly in the nucleon. This has been under hot debate for over ten years but one obtained not yet a fully satisfactory description (for a review of the nucleon spin problem, see, e.g., ). It should be emphasized that the “crises” is not for the fundamental theory of quantum chromodynamics (QCD), in the viewpoint of which the nucleon is a complicated object of quarks and gluons and quarks do not necessarily carry most of the nucleon spin. The “crises” is, however, for the naive SU(6) quark model which is quite successful in many aspects but nevertheless attributes all the nucleon spin to constituent quarks. To explore whether this “crises” is real, i.e., whether the SU(6) model could be taken as a good lowest order approximation for the nucleon, a natural way is to start from the SU(6) wavefunction and study whether we can explain the experimental result of the nucleon spin content by going to higher orders.
In the past years there has been countless work using quark models along this direction (for references see ), but the problem is not really resolved. The obstacle is that the quark axial current, which is the operator for defining the quark spin contribution to nucleon spin, is a non-conserved composite operator. Therefore when we go to higher orders divergent matrix element will be encountered and the quark axial current must be renormalized. But unfortunately, the usual renormalization schemes are not applicable in quark model calculations, where instead of divergent loop integrations over the continuous momentum, we encounter divergent summations over the discrete quark excited states whose wavefunctions are obtain numerically. To explore how one can renormalize a composite operator in quark models and what would be the result of the renormalized quark axial charge (i.e., the quark spin contribution to nucleon spin) are the aims of this paper.
In the following we will first construct the bare matrix element of the quark axial current and demonstrate its divergence, then we explore how we can renormalize the bare quantity by systematically subtracting a divergent part and obtain a (finite) physical result.
The quark spin contribution to the nucleon spin is defined as the quark axial charge of the nucleon:
$$ps|\overline{\psi _q}\gamma ^\mu \gamma ^5\psi _q|ps\overline{u}_{ps}\gamma ^\mu \gamma ^5u_{ps}\mathrm{\Delta }_q,$$
(1)
where $`q=u,d,s`$ and $`u_{ps}`$ is the nucleon spinor. An equivalent but more suitable expression for model calculation is
$$\mathrm{\Delta }_q=\frac{p+|d^3x\overline{\psi _q}\gamma ^3\gamma ^5\psi _q|p+}{p+|p+}.$$
(2)
Here $`|p+`$ is a nucleon state with positive momentum and polarization along the third direction. We are going to adopt Eq. (2) to pursue a perturbative calculation of $`\mathrm{\Delta }_q`$ in a chiral potential model. Our model Lagrangian is
$``$ $`=`$ $`\overline{\psi }[i/S(r)\gamma ^0V(r)]\psi `$ (5)
$`{\displaystyle \frac{1}{2F_\pi }}\overline{\psi }[S(r)(\sigma +i\gamma ^5\lambda ^i\varphi _i)+(\sigma +i\gamma ^5\lambda ^i\varphi _i)S(r)]\psi +`$
$`{\displaystyle \frac{1}{2}}(_\mu \sigma )^2+{\displaystyle \frac{1}{2}}(_\mu \varphi _i)^2{\displaystyle \frac{1}{2}}m_\sigma ^2\sigma ^2{\displaystyle \frac{1}{2}}m_i^2\varphi _i^2.`$
The model Lagrangian is derived from the $`\sigma `$ model in which meson fields are introduced to restore chiral symmetry . The flavor and color indices for the quark field $`\psi `$ are suppressed; the scalar term $`S(r)=cr+m`$ represents the linear scalar confinement potential $`cr`$ and the quark mass matrix $`m`$; $`V(r)=\alpha /r`$ is the Coulomb type vector potential and $`F_\pi `$=93MeV is the pion decay constant. $`\sigma `$ and $`\varphi _i`$ ($`i`$ runs from $`1`$ to $`8`$) are the scalar and pseudoscalar meson fields, respectively and $`\lambda _i`$ are the Gell-Mann matrices. The quark-meson interaction term of Eq.(5) is symmetrized since the mass matrix $`m`$ does not commute with all $`\lambda _i`$ for different quark masses.
At zeroth order the nucleon is taken as the usual SU(6) three-quark ground state of the Hamiltonian
$$H_q=d^3x\psi ^{}[\stackrel{}{\alpha }\frac{1}{i}\stackrel{}{}+\beta S(r)+V(r)]\psi .$$
(6)
The diagrams for the numerator and denominator of Eq. (2) up to second order are shown in Figs. 1 and 2 respectively.
We first discuss how to determine the renormalization constant $`Z_2`$. The mass-shell renormalization scheme is not applicable here, since our unperturbed quark basis are confined wavefunctions. But we can still use the charge renormalization condition. The conserved electromagnetic current of the Lagrangian of Eq. (5) is $`j^\mu =_qj_q^\mu +j_\varphi ^\mu `$, where $`j_q^\mu `$ and $`j_\varphi ^\mu `$ are the quark and meson current respectively:
$`j_q^\mu `$ $`=`$ $`Q_q\overline{\psi }_q\gamma ^\mu \psi _q,`$ (7)
$`j_\varphi ^\mu `$ $`=`$ $`e(\varphi _1^\mu \varphi _2\varphi _2^\mu \varphi _1+\varphi _4^\mu \varphi _5\varphi _5^\mu \varphi _4).`$ (8)
The charge renormalization condition is to require that for a quark state
$$q|d^3xj^0(x)|q=Q_q.$$
(9)
Up to second order this is shown in Fig. 3.
By computing Figs. 3C and 3D we can determine the renormalization constant $`Z_2`$, which is then to be used in Fig. 1B.
Next we remark that at zero momentum transfer, the exchange diagram Fig. 1D is actually the product of Figs. 1A and 2B, therefore the sum of Figs. 1A and 1D over the normalization Figs. 2A and 2B is just Fig. 1A. (However this is not true when we calculate the axial form factor at finite momentum transfer.) Thus only Fig. 1C is left to be evaluated together with Fig. 3C and 3D.
The essential ingredients needed for calculating these diagrams are the quark and meson propagators. The meson propagator given by the Lagrangian of Eq. (5) is the free propagator:
$$\mathrm{\Delta }_{ij}(x_1,x_2)=\frac{i}{(2\pi )^4}d^4q\frac{\delta _{ij}e^{iq(x_1x_2)}}{q^2m_i^2+iϵ}.$$
(10)
Since the non-perturbative confinement is included in $`H_q`$ the quark propagator has to be obtained numerically, and in practise we have to work with time-ordered perturbation theory. We write the solution of $`H_q`$ as
$$\psi (x)=\underset{\alpha }{}u_\alpha (x)a_\alpha +\underset{\beta }{}v_\beta (x)b_\beta ^{},$$
(11)
where $`u_\alpha (x)=e^{iE_\alpha t}u_\alpha (\stackrel{}{x})\tau _\alpha `$, $`v_\beta (x)=e^{iE_\beta t}v_\beta (\stackrel{}{x})\tau _\beta `$; $`\tau `$ is the flavor wavefunction and $`u_\alpha (\stackrel{}{x})`$ and $`v_\alpha (\stackrel{}{x})`$ are the spatial wavefunctions. The quark propagator is then
$`D(x_1,x_2)`$ $`=`$ $`\theta (t_1t_2){\displaystyle \underset{\alpha }{}}u_\alpha (x_1)\overline{u}_\alpha (x_2)`$ (13)
$`\theta (t_2t_1){\displaystyle \underset{\beta }{}}v_\beta (x_1)\overline{v}_\beta (x_2).`$
It can be shown by carrying out the time and energy integration that apart from an isospin factor, Figs. 3C and 3D yield the same expression. Therefore we define for Figs. 3C and 3D the pure space-time amplitudes:
$`B_\varphi {\displaystyle \frac{1}{F_\pi ^2}}`$ $`{\displaystyle }d^3xd^4x_1d^4x_2\mathrm{\Delta }(x_2,x_1)\times `$ (15)
$`\overline{u}_f(x_2)\mathrm{\Gamma }_\varphi D(x_2,x)\gamma ^0D(x,x_1)\mathrm{\Gamma }_\varphi u_i(x_1),`$
where $`u_i`$ and $`u_f`$ are the initial and final quark state respectively, the vertex function $`\mathrm{\Gamma }_{\pi ,K,\eta }=S(r)\gamma ^5`$ and $`\mathrm{\Gamma }_\sigma =iS(r)`$. All flavor wavefunctions are dropped out and it should be understood that for $`\varphi =K`$ the intermediate quark is the $`s`$ quark and otherwise is the $`u`$ or $`d`$ quark. Accordingly for Fig. 1C we define
$`A_\varphi {\displaystyle \frac{1}{F_\pi ^2}}`$ $`{\displaystyle }d^3xd^4x_1d^4x_2\mathrm{\Delta }(x_2,x_1)\times `$ (17)
$`\overline{u}_f(x_2)\mathrm{\Gamma }_\varphi D(x_2,x)\gamma ^3\gamma ^5D(x,x_1)\mathrm{\Gamma }_\varphi u_i(x_1).`$
Now we can express the renormalization constant $`Z_2`$ and the axial charge $`\mathrm{\Delta }_q`$ in terms of $`A_\varphi `$ and $`B_\varphi `$ multiplied by spin and isospin factors which are calculated straightforwardly:
$`Z_2^{u,d}`$ $`=`$ $`1(3B_\pi +2B_K+{\displaystyle \frac{1}{3}}B_\eta +B_\sigma ),`$ (18)
$`Z_2^s`$ $`=`$ $`1(B_K+{\displaystyle \frac{4}{3}}B_\eta +B_\sigma ),`$ (19)
$`\mathrm{\Delta }_u`$ $`=`$ $`{\displaystyle \frac{4}{3}}f_RZ_2^u+{\displaystyle \frac{2}{3}}A_\pi +{\displaystyle \frac{4}{9}}A_\eta +{\displaystyle \frac{4}{3}}A_\sigma `$ (20)
$`\mathrm{\Delta }_d`$ $`=`$ $`{\displaystyle \frac{1}{3}}f_RZ_2^d+{\displaystyle \frac{7}{3}}A_\pi {\displaystyle \frac{1}{9}}A_\eta {\displaystyle \frac{1}{3}}A_\sigma `$ (21)
$`\mathrm{\Delta }_s`$ $`=`$ $`2A_K.`$ (22)
Another useful relation is for the nucleon axial charge
$$g_A=\mathrm{\Delta }_u\mathrm{\Delta }_d=\frac{5}{3}f_RZ_2^u\frac{5}{3}A_\pi +\frac{5}{9}A_\eta +\frac{5}{3}A_\sigma $$
(23)
In Eqs. (19)-(23) we have assumed equal masses for $`u,d`$ and for $`\pi ^0,\pi ^\pm `$, therefore $`Z_2^u=Z_2^d`$ which is just a statement of SU(2) symmetry (If SU(3) symmetry is unbroken $`Z_2^s`$ would also be the same as $`Z_2^{u,d}`$); $`5/3f_R`$ is the zeroth order values of $`g_A`$ and $`f_R`$ is a relativistic reduction factor.
Since $`A_\varphi `$ and $`B_\varphi `$ correspond to loop diagrams, they would naturally be divergent. If the quark axial current were conserved, then the divergence of $`A_\varphi `$ and $`B_\varphi `$ would automatically cancel in Eq. (22) and we get a finite result for $`\mathrm{\Delta }_q`$. This is the phenomenon that conserved operators do not need extra renormalization besides the usual renormalization of mass, charge and wavefunction. However the quark axial current is not conserved, the divergences of $`A_\varphi `$ and $`B_\varphi `$ will not cancel and the naively obtained $`\mathrm{\Delta }_q`$ in Eq. (22) is divergent, which is just the general case that a composite operator has divergent matrix element and needs extra renormalization, i.e., we are going to subtract a divergent piece from $`A_\varphi `$ and $`B_\varphi `$ simultaneously and leave a finite part in Eq.(22). This finite leftover would depend on how we regularize $`A_\varphi `$ and $`B_\varphi `$ and how much is to be subtracted as the divergent part. This is the renormalization scheme dependence.
In the usual plane-wave perturbation theories, we have a most powerful and systematic renormalization scheme of dimensional regularization and minimal subtraction (MS) or modified minimal subtraction ($`\overline{\mathrm{MS}}`$), but it is evidently not applicable here. We must find a proper renormalization scheme for the quark model calculations. The divergent integration over the (continuous) momentum in the plane-wave perturbation theories is here contained in the summation over the (discrete) quark intermediate states. Thus we propose to regularize by cutting the summation at a certain energy, which is analogous to the lattice regularization by a finite lattice spacing. In the following we explore how this scheme works in practice.
Our model parameters are listed in Table I. As we demonstrated in a recent paper , $`A_\varphi `$ is actually very insensitive to the model parameters in the above described energy-cutoff regularization scheme. This property is also found for $`B_\varphi `$. Therefore we do not take too much effort in choosing the parameters except for a fit to the nucleon and $`\mathrm{\Delta }`$ masses with meson exchange potentials.
Figs. 4 and 5 give the numerical results of $`B_\varphi `$ and $`A_\varphi `$ as a function of the maximum energy up to which one sums the intermediate quark states. The divergences are clearly seen. The problem is now how to determine the cutoff. In general the cutoff might vary from one diagram to another, without any guidance we would be at lost. Fortunately, besides $`\mathrm{\Delta }_s=2A_K`$, Eq. (22) provides another clean relation:
$$\mathrm{\Delta }_u+4\mathrm{\Delta }_d=10A_\pi .$$
(24)
Thus with the experimental results:
$`\mathrm{\Delta }_u`$ $`=`$ $`0.80(6),\mathrm{\Delta }_d=0.46(6),\mathrm{\Delta }_s=0.12(4),\text{[7]}`$ (25)
$`\mathrm{\Delta }_u`$ $`=`$ $`0.82(6),\mathrm{\Delta }_d=0.44(6),\mathrm{\Delta }_s=0.10(4),\text{[8]}`$ (26)
we can determine the cutoff for $`A_K`$ and $`A_\pi `$ using Fig. 5. It is amazing to notice that the cutoffs needed for $`A_K`$ and $`A_\pi `$ are roughly the same, i.e., independent of whether the intermediate quark is $`u,d`$ (for $`A_\pi `$) or $`s`$ (for $`A_K`$) and also independent of the meson masses. This reminds us that in the MS or $`\overline{\mathrm{MS}}`$ scheme, the subtraction is independent of the mass parameters. We will therefore choose a $`\varphi `$-independent cutoff. It is also very interesting to notice that the cutoff value needed here roughly equals the inverse of the lattice spacing $`a^1`$ in lattice QCD calculation of $`\mathrm{\Delta }_q`$. To reduce the number of parameters, we choose the cutoff for $`A_\varphi `$ to be the same as $`a^1=1.74`$GeV in . But there remains still one important question to be asked: should the cutoff for $`A_\varphi `$ and $`B_\varphi `$ be the same?
It might be taken for granted that they are the same, (the simplest example is that we use the charge renormalization condition to determine the renormalization constant and then calculate the matrix element of the charge operator itself). However, we call special attention here that this is not necessarily the case. The problem is the lack of Lorentz covariance in the model.
It is actually not difficult to understand this point: If a theory does not respect Lorentz covariance, then the renormalization constant determined with the time and spatial component of the electromagnetic currents are possibly different. Therefore if we calculate the matrix element of the charge operator but use the renormalization constant determined through the renormalization condition for the spatial component of the vector current, then we might not obtain the physical charge (in some cases the result may still be finite, but this is not enough).
In the present model, Lorentz covariance is violated by the static potentials, therefore the renormalization constant determined for the charge might not be suitable for calculating the axial charge. But since we have presently no better method, we must allow different cutoffs for $`A_\varphi `$ and $`B_\varphi `$.
After describing the formalism, we are now in the position to see how this renormalization scheme can explain the experimental results. As we explained the cutoff for $`A_\varphi `$ is chosen as the inverse of the lattice spacing $`a^1=1.74`$GeV . Using Eq. (23) and Figs. 4, 5 and requiring $`g_A=1.257`$, the cutoff for $`B_\varphi `$ is determined to be $`0.734`$GeV. Then combining Figs. 4, 5 and using Eq. (22), we finally obtain:
$$\mathrm{\Delta }_u=0.823,\mathrm{\Delta }_d=0.432,\mathrm{\Delta }_s=0.104.$$
(27)
By adjusting only one free parameter (the cutoff for $`B_\varphi `$) to fit $`g_A`$, we reproduced the experimental results of $`\mathrm{\Delta }_{u,d,s}`$ perfectly. This can be regarded as a success of our model and renormalization scheme. Since we have started from the SU(6) wavefunction as the zeroth order approximation, we could say that the spin “crisis” for the naive SU(6) model is not real.
We end our discussions by emphasizing three points: (1) Non-conserved operators have divergent matrix elements and even in model calculations they must be consistently renormalized. (2) A natural and convenient renormalization scheme in quark model calculations is to cut the summation over the intermediate quark states at a certain energy; a fit to the experimental data reveals that this approach shares the same advantage of mass-parameter-independence as in the MS or $`\overline{\mathrm{MS}}`$ scheme. (3) However, due to the lack of Lorentz covariance in quark models, the cutoffs for operators of different Lorentz type are not necessarily the same.
This work is supported by the CNSF (19675018), CSED, CSSTC, the DFG (FA67/25-1), and the DAAD.
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# FUSE Spectroscopy of High Velocity Cloud Complex C
## 1 Introduction
Despite significant observational and theoretical efforts over nearly 40 years of study, the nature of Galactic H I high velocity clouds (HVCs) is still a mystery. With the launch of the Far Ultraviolet Spectroscopic Explorer (FUSE), a new and important portion of the electromagnetic spectrum is now available to study HVCs. The far-ultraviolet provides a wealth of atomic, ionic, and molecular spectral lines that can be used to probe the physical conditions in HVCs (Sembach, 1999). With its high resolution and large effective area, FUSE has the ability to probe HVCs along multiple lines of sight toward extragalactic objects. This paper presents an analysis of such a sightline and demonstrates the ability of FUSE to significantly contribute to our understanding of HVCs.
The sight line to Mrk 876 ($`l=98\stackrel{}{\mathrm{.}}27`$, $`b=40\stackrel{}{\mathrm{.}}38`$) passes through HVC complex C (Wakker & van Woerden, 1991). Wakker et al. (1999) suggest the distance to complex C is in the range of 5-25 kpc, with $`\mathrm{D}10`$ kpc as the most likely value. They derive a mass for complex C of $`6\times 10^6\mathrm{M}_{\mathrm{}}(\mathrm{D}/10\mathrm{kpc})^2`$ which, combined with the observed radial velocity, implies a mass influx of $`0.080.19(\mathrm{D}/10\mathrm{kpc})\mathrm{M}_{\mathrm{}}\mathrm{yr}^1`$. The metallicity of complex C is in the range 0.1 to 0.6 solar (Wakker et al., 1999; Gibson et al., 2000) and may be spatially variable toward multiple sightlines separated by more than 10°. The metallicity of HVCs can be used to discriminate between the current theories for their origins (van Woerden et al., 1999).
Along the Mrk 876 sightline, a 21 cm spectrum taken with the Effelsberg 100 m telescope with a 9$`\stackrel{}{\mathrm{.}}`$7 beam shows two components at $`\mathrm{V}_{\mathrm{LSR}}=172\mathrm{km}\mathrm{s}^1`$ and $`\mathrm{V}_{\mathrm{LSR}}=133\mathrm{km}\mathrm{s}^1`$ with N(H I)$`=(4.1\pm 0.8)\times 10^{18}\mathrm{cm}^2`$ and N(H I)$`=(19.2\pm 0.8)\times 10^{18}\mathrm{cm}^2`$, respectively. An NRAO<sup>1</sup><sup>1</sup>1The National Radio Astronomy Observatory is a facility of the National Science Foundation, operated under cooperative agreement by Associated Universities, Inc. 140 Foot Telescope (43 meter) observation with a 21′ beam gives a column density that is 2.2 times higher than the Effelsberg spectrum for the $``$172 $`\mathrm{km}\mathrm{s}^1`$ component while both telescopes give identical results for the $``$133 $`\mathrm{km}\mathrm{s}^1`$ component. The difference between the 9$`\stackrel{}{\mathrm{.}}`$7 and 21′ results probably arises from small scale structure in the HVC, since it is much larger than any expected statistical uncertainties or calibration differences. We have used the Effelsberg spectrum in the analysis that follows.
An H$`\alpha `$ spectrum obtained by M. Haffner (private communication) with the WHAM instrument (Reynolds et al., 1998) indicates that the H$`\alpha `$ intensity associated with complex C toward Mrk 876 is $`<0.02`$ R (1 R is 10<sup>6</sup> photons$`\mathrm{cm}^2(4\pi \mathrm{sr})^1`$) averaged over the 1° beam of the instrument. Complex C is clearly detected in other directions (Wakker et al. 1999; Tufte et al. 2000, in preparation) which implies a patchy distribution of N(H<sup>+</sup>). Since both the distance to complex C and the geometry of the emitting region are unknown, and there may exist small scale structure in the large WHAM beam, no meaningful limit can be placed on N(H<sup>+</sup>) at this time.
## 2 FUSE Observations
For a description of the FUSE satellite, its operation, and observing modes see Moos et al. (2000) and Sahnow et al. (2000). The Mrk 876 dataset (P1073101) consists of 10 consecutive exposures in the $`30\mathrm{}\times 30\mathrm{}`$ apertures, resulting in 52 ksec of on-target integration. At the time of the observations (16 October 1999) the spectrograph was not yet aligned or focused. The target appears in both LiF channels and one SiC channel. The data were passed through the standard CALFUSE pipeline, which removed the spectral motion, performed a background subtraction, removed the geometric distortions, extracted the spectra, and performed wavelength and flux calibrations. Event bursts were removed from the data by hand. Coaddition of the channels was not attempted due to the preliminary nature of the wavelength calibration. Typically, the equivalent width of an absorption line was measured separately for each channel and averaged. The typical spectral resolution was $`\lambda /\mathrm{\Delta }\lambda 12000`$. The average S/N is 14:1 per resolution element in the LiF1 channel, 17:1 in the LiF2 channel, and 8:1 in the SiC1 channel.
## 3 First Detection of an O VI HVC
The O VI lines are shown in Figure 1. Absorption from HVC complex C is clearly present at negative velocities as high as $`215\mathrm{km}\mathrm{s}^1`$. O VI is an excellent tracer of hot gas ($`25\times 10^5`$ K) and is not generally produced by photoionization. Beyond $`215\mathrm{km}\mathrm{s}^1`$ the O VI 1031.93 Å line is blended with the Galactic H<sub>2</sub>(6-0) P(3) line at 1031.19 Å. Fortunately, the H<sub>2</sub> line is narrow and was easily removed from the spectrum. Integrating the profile between $`215\mathrm{km}\mathrm{s}^1`$ and $`100\mathrm{km}\mathrm{s}^1`$ we find and equivalent width of $`146\pm 14`$ mÅ which implies N(O VI)=$`(1.5\pm 0.2)\times 10^{14}\mathrm{cm}^2`$ if the line is optically thin. High velocity absorption is also seen in the O VI 1037.62 Å line; however, it is blended with the strong Galactic H<sub>2</sub>(5-0) R(1) line at velocities more negative than $`140`$ $`\mathrm{km}\mathrm{s}^1`$, making a column density measurement impossible.
The amount of hot gas implied by the O VI detection can be significant if the metallicity of complex C is low. The models of Sutherland & Dopita (1993) give the ionization fraction of O VI as function of temperature for gas in collisional equilibrium. At the peak of the O VI ionization fraction (22% at 280,000 K), the amount of hot gas is N(H)$`>\mathrm{\hspace{0.25em}7}\times 10^{17}\mathrm{cm}^2(\mathrm{O}/\mathrm{H})_{\mathrm{}}/(\mathrm{O}/\mathrm{H})_{\mathrm{HVC}}`$.
A discussion of the possible interpretations for the detection of high-velocity O VI in complex C and other HVCs is presented by Sembach et al. (2000).
## 4 Metal lines in Complex C
Metal line absorption associated with the $`133\mathrm{km}\mathrm{s}^1`$ H I component (and, in some cases, the $`172\mathrm{km}\mathrm{s}^1`$ component) is seen in several Fe II lines, C II, N I, and N II. Upper limits can be set for Ar I, P II, and Fe III. Example absorption lines are presented in Figure 1. Table 1 lists the measured equivalent widths, column densities and derived abundances relative to the solar values. The errors are a combination in quadrature of the statistical error, based on the S/N ratio of the spectrum, and systematic errors calculated as the rms value of measurements made by different authors, by using different detector segments, different continuum placement, and by using subsets of the data (e.g. night-only).
Five HVC Fe II lines from the $`133\mathrm{km}\mathrm{s}^1`$ component are detected, from which a curve of growth can be derived. We minimize $`\chi ^2`$(N,$`b`$)=\[W<sub>λ</sub>(observed)$`{}_{}{}^{2}`$W<sub>λ</sub>(N,$`b`$)<sup>2</sup>\], and find N=3.0$`\pm `$1.2$`\times 10^{14}`$$`\mathrm{cm}^2`$ and $`b`$=12.1$`\pm `$5.6 $`\mathrm{km}\mathrm{s}^1`$. All observed values of W<sub>λ</sub>are then within 1$`\sigma `$ of the expected value. We will use this $`b`$-value below to convert equivalent width to column density for other ions (excluding O VI). The $``$172 $`\mathrm{km}\mathrm{s}^1`$ component is also seen in absorption in Fe II $`\lambda `$1144.94. The Fe III $`\lambda `$1121.98 line from the $`133\mathrm{km}\mathrm{s}^1`$ component is absent, yielding N(Fe III)/N(Fe II)$`<`$0.22.
N II $`\lambda `$1083.99 high-velocity absorption can clearly be seen in the SiC1 channel; however, the individual components are blended and the spectrum has low S/N. Integrating the line and apparent column density profiles between $``$160 and $``$100 $`\mathrm{km}\mathrm{s}^1`$ yields an equivalent width $`>`$120 mÅ and a column density of N(N II)$`>1.0\times 10^{14}\mathrm{cm}^2`$. If we calculate the column density assuming the $`b`$ value found above for iron, the result is a factor of three higher. The high-velocity component of N I $`\lambda `$1134.17 is blended with low-velocity Fe II $`\lambda 1133.67`$. Using a curve of growth for the low velocity gas, we derive W<sub>λ</sub>(Fe II $`\lambda `$1133.67)=67$`\pm `$7 mÅ. The measured width of the blend is 94$`\pm `$7 mÅ, so that 27$`\pm `$10 mÅ can be attributed to high-velocity N I. This matches, to within the errors, the column density of N I $`\lambda `$1199.55 found by Gibson et al. (2000) toward Mrk 876 using the Space Telescope Imaging Spectrograph on the Hubble Space Telescope.
C II $`\lambda `$1036.34 is clearly present in both components of the HVC and is strongly saturated. For P II $`\lambda `$1152.82 and Ar I $`\lambda `$1048.22 only $`3\sigma `$ upper limits of 25 mÅ can be set at this time.
## 5 Metal abundances in complex C
The observed ratio N(Fe II)/N(H I)$`0.5`$ (Fe/H) is unexpected given that every previously studied sightline through cool, warm, or halo gas has shown iron depleted by at least a factor 3 (Savage & Sembach, 1996). If the intrinsic metallicity of complex C is subsolar, then little, if any, iron can be depleted into dust grains. The iron abundance is also higher than the value of S/H$``$0.1 (S/H) found by Wakker et al. (1999) along the Mrk 290 sightline through complex C. Gibson et al. (2000) have measured S II absorption along the lines of sight to Mrk 817, where N(S II)/N(H I)=0.3 (S/H), and Mrk 279, where N(S II)/N(H I)$``$0.6 (S/H). They find that assuming the presence of H<sup>+</sup> along the line of sight is insufficient to reconcile their observations with the metallicity found by Wakker et al. (1999). Instead, they believe that the metallicity of complex C is spatially variable with metallicities ranging from $`0.10.6`$ solar. The low abundance of argon (N(Ar I)/N(H I)$`<0.13`$ (Ar/H)) toward Mrk 876 is probably a photoionization effect; Ar I has a photoionization cross-section about 10 times larger than that of hydrogen (Sofia & Jenkins, 1998). Upcoming FUSE observations of the Mrk 290, Mrk 817, and Mrk 279 sightlines should help resolve these issues.
## 6 Limits on Molecular Hydrogen
Previous searches for molecular gas in high velocity clouds concentrated on tracers such as CO and HCO<sup>+</sup>, which have spectral lines at millimeter wavelengths. Wakker et al. (1997) and Akeson & Blitz (1999) placed upper limits on N(H<sub>2</sub>) of $`9\times 10^{18}\mathrm{cm}^2`$ and $`5.5\times 10^{18}\mathrm{cm}^2`$ in several HVCs, respectively, assuming solar abundances and standard I(CO) to N(H<sub>2</sub>) ratios. Richter et al. (1999) have reported the discovery of H<sub>2</sub> absorption at a velocity of +120 $`\mathrm{km}\mathrm{s}^1`$ with N(H<sub>2</sub>)=$`(2.23.6)\times 10^{15}\mathrm{cm}^2`$ in an ORFEUS spectrum of the LMC star HD 269546 (Sk$``$68 82). However, this sightline is quite complicated, and has absorption associated with the Milky Way, the LMC, as well as the HVC.
FUSE can observe the far UV electronic transitions of molecular hydrogen from both the Lyman (B–X) and Werner (C–X) bands at high resolution toward many extragalactic objects behind HVCs. Toward Mrk 876, absorption lines from H<sub>2</sub> in the Milky Way are readily visible with N(H<sub>2</sub>)$`2.3\times 10^{18}\mathrm{cm}^2`$ (Shull et al., 2000). Figure 2 compares the observed spectra and the expected absorption signatures of the two components of complex C for an H<sub>2</sub> column density of $`2\times 10^{14}\mathrm{cm}^2`$ in each component and for a rotational temperature of T$`{}_{01}{}^{}=100`$ K and $`b=10\mathrm{km}\mathrm{s}^1`$. The selected lines are the strongest H<sub>2</sub> lines in the FUSE bandpass that are not blended with other atomic or molecular lines. Clearly, there is no detectable H<sub>2</sub> absorption from HVC complex C. Reasonable changes in T<sub>01</sub> or $`b`$ do not alter this conclusion. Since the most efficient mechanism for forming H<sub>2</sub> is on the surfaces of dust grains (Spitzer, 1978), the non-detection of H<sub>2</sub>, coupled with the low depletion of iron, implies that there is little or no dust in complex C along this line of sight.
This work is based on data obtained for the Guaranteed Time Team by the NASA-CNES-CSA FUSE mission operated by the Johns Hopkins University. Financial support to U.S. participants has been provided by NASA contract NAS5-32985. We thank Peter Kalberla for observing the Effelsberg H I spectrum. The Effelsberg Telescope is operated by the Max Planck Institut für Radio Astronomie at Bonn. We also thank Matt Haffner for providing the WHAM H$`\alpha `$ result.
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# RU-NHETC-2000-16 PUPT-1928 hep-th/0005047 Near Hagedorn Dynamics of NS Fivebranes or A New Universality Class of Coiled Strings
## 1 Introduction
One of the more puzzling objects in string theory is the NS fivebrane. In particular the worldvolume theory, the so-called “Little String Theory” (LST), is believed to be some kind of non-gravitational string theory. These theories were introduced in (see also , for a review see ). In particular in the decoupled theory on a cluster of NS fivebranes was defined by taking the limit
$$g_s0,M_sfixed$$
(1)
of string theory in the presence of the cluster. This definition, however, is rather indirect in the sense that it does not provide a microscopic description of the theory. Rather, it is closely related to the gravitational holographic dual (for a review see ) of the theory.
The gravitational side of the holographic duality for this case is the throat region of the CHS background . This background includes a linear dilaton direction, an $`R^6`$ component and a WZW model at a level set by the number of fivebranes $`N`$. This duality was used to analyze the observable content of the theory at the origin of its moduli space , and along its flat directions . This was done both for the type II NS fivebranes, the heterotic NS fivebranes , fivebranes at orbifolds and lower dimensional related configurations . Fivebrane thermodynamics was discussed in .
A complementary way of exploring “little string theories” was introduced in and developed further in . In this approach a discrete light cone quantization was suggested, along the lines of Matrix theory . This DLCQ description is in terms of a 1+1 sigma model on the ADHM moduli space. Even though some aspects of this model are well understood (either from AdS/CFT or from a direct field theory analysis ), it is not yet clear how to directly obtain LST quantities from it and in particular how to obtain a covariant microscopic description (we suggest a way of going to the “long strings” picture for this sigma-model in section 6). For another suggestion, involving quasi-local field theories, see .
Despite the difficulties, it is worthwhile to explore these theories for a variety of reasons. The basic point is that these are stringy theories without gravity and without a tunable string coupling. This makes them a very intriguing object to study. In addition , they are important for a Matrix description of M-theory on $`T^5`$. Moreover, in the context of holographic duality the “little string theory”/CHS background duality behaves differently from the AdS/CFT duality, and therefore might be a first step towards understanding holography in asymptotically flat spaces .
Another interesting aspect of these theories is an unusual UV/IR relation. When the theories are taken along the flat directions , such that the Higgs vev is $`M_w^2`$, then in addition to the expected massless fields ($`N`$ tensor(vector) multiplets for type IIA(IIB)), one finds additional massless states. In fact one finds an entire stringy tower of massive states with string scale $`M_s/\sqrt{N}`$, even in the limit $`M_w\mathrm{}`$. This is very puzzling because naively one expects the theory in this limit to split into $`N`$ copies of the theory of a single NS fivebrane. Rather the low energy theory is still sensitive to details of the very high energy states. This is somewhat reminiscent of effects in other non-local quantum theories, i.e., the theories on non-commutative geometries . The difference is that the extra states here are propagating particles in Minkowski space and the theory is Lorentz invariant.
More recently, these theories have appeared in the context of holographic duals of confining gauge theories . Some of the confining vacua of $`𝒩=4`$ deformed to $`𝒩=1`$ pure glue are described holographically using NS fivebranes. This suggests a role for “little strings” in the confinement picture, or more generally as an approximate description, valid at some finite energy interval, of the IR region of field theories.
In this paper we discuss the behavior of LST at high energy densities. The holographic dual to this configuration is the near horizon limit of the near extremal fivebranes , which includes the CGHS black hole . We are interested in the Euclidean black hole which means that we are discussing the canonical ensemble. From this background it is easy to extract that the theory has an Hagedorn density of states and hence a Hagedorn temperature (for discussions see ). We are interested at the behavior of the theory as the temperature approaches the Hagedorn temperature from below.
The behavior at high energies is sometimes expected to resemble weakly coupled string theory . We attempt to interpret the thermodynamics in terms of a gas of weakly coupled strings. To this end we begin, in the next section, by reviewing a class of models of free strings and the resulting thermodynamics. This is meant to be compared to qualitative features, and not necessarily precise quantitative ones, of the LST thermodynamics. We discuss possible different behaviors near the Hagedorn temperature, the validity of the canonical ensemble and other features.
In sections 3 and 4 we discuss the thermodynamics of the CGHS black hole. Section 3 is a review of known results at string tree level and section 4 discusses one loop corrections, which allow us to go off the Hagedorn temperature and study the partition function as we approach it. The results strongly suggest that the Hagedorn temperature is a limiting temperature for LST, rather than a phase transition as in weakly coupled critical string theory.
Section 5 is somewhat of an aside - we discuss how the large fluctuations of the canonical ensemble manifest themselves in the near extremal solution.
In section 6 we try to interpret our results in terms of the dynamics of an almost free string and show that there is some qualitative difference between the two systems. We interpret this as indicating a new universality class of interacting strings. We then suggest that in this universality class the strings have a strong self-attractive potential. This makes long strings want to shrink, which is a first step towards explaining the thermodynamics of the NS 5-branes. By a simple combinatoric random walk model we argue that this phase can not occur for strings above 5+1 dimensions, which explains the maximal dimensionality of “little string theory” (although it does point to the fact that a similar modification for the theory of membranes might lead to an interacting theory even in higher dimensional spacetimes). Our analysis is based both on space-time considerations and on DLCQ considerations, and points to a new way of analyzing “long strings” in the D1-D5 system.
We conclude by summarizing the main results of the paper.
As we were completing this project, we received a paper with some overlap with sections 3 and 4 in our paper. The interpretation of the result is, however, quite different.
## 2 Ensembles of Weakly Coupled Strings
Let us review some aspects of the thermodynamics of critical string theory, which are relevant for us (these would be similar to .). Of course, the notion of a thermodynamic equilibrium in a theory with gravity is ill defined, but as explained in it is justified in the weak string coupling limit, where some questions can still be addressed. The question we are interested in is that of the high energy density of states (at the sphere level). We do so with an eventual goal of examining what aspects of this high energy spectrum remain valid in “little string theory”. Note that since LST does not contain gravity, a thermal equilibrium is well defined even though there is no weak coupling.
Since the discussion of qualitative features of free strings is sufficient for us, we restrict our attention to the simplest model on the worldsheet, i.e. free bosons and fermions. However, we allow for an arbitrary central charge $`\widehat{c}_{eff}`$<sup>1</sup><sup>1</sup>1The reduced central charge is defined as usual, $`\widehat{c}=\frac{2}{3}c`$. We will refer to $`\widehat{c}`$ as the central charge in the following. (in light cone), and an unknown string tension $`M_{eff}^2`$. We allow this freedom since the immediate information we have about LST is the Hagedorn temperature which only determines the combination
$$\frac{\widehat{c}_{eff}}{M_{eff}^2}=\frac{4N}{M_s^2}$$
(2)
where $`N`$ is the number of NS 5-branes and $`M_s`$ is the spacetime string tension. This can be derived from the Euclidean near extremal solution by measuring the periodicity of the time direction, which turns out to be the same for all values of the energy density. For example, the free string model of LST in has strings of central charge 4 (4 bosonic coordinates in light cone and their supersymmetry partners) and a tension $`M_s^2/N`$. However, the Matrix model (with a single unit of null momentum) suggests a tension of $`M_s`$ and central charge $`4N`$. Hence we would like to keep the central charge and tension arbitrary (subject to the constraint (2)), and see which one better fits the “data”.
There are several ensembles that one might use. One can use either the microcanonical or the canonical ensembles, and then one can use either a compactified or uncompatified space. It is at times stated that the canonical ensemble is unreliable in a theory with a Hagedorn density because there are large fluctuations in thermodynamic quantities. For example the fluctuations of the average energy are:
$$\frac{<E^2><E>^2}{<E>^2}1$$
(3)
or much larger (as we remind the reader shortly). However, the CGHS black hole can be formulated just as well in Euclidean signature as in Lorentzian one, therefore one expects the holographic relation to hold just as well in the Euclidean/canonical ensemble case. In fact we will show that the CGHS black hole predicts large fluctuations just below the Hagedorn temperature. Henceforth we restrict our attention to Euclidean signature and to the canonical ensemble.
One can also work either in compact on non-compact space, i.e., LST either on $`R^{5,1}`$ or $`R\times T^5`$ where $`T^5`$ has some volume $`V`$. We compute the thermodynamics in both cases and examine their behavior as $`TT_H`$. In the case of compact space we will also work in the limit $`V\mathrm{}`$, however we will always take $`TT_H`$ first. Remember that these two limits do not commute in a free string theory. The limit in which $`V\mathrm{}`$ first for fixed $`T`$ converges to the uncompactified limit, whereas when $`TT_H`$ first , for fixed V, the result is different. One obtains a result which is independent of $`V`$ and different from the non-compact case.
We will incorporate an arbitrary central charge by having 4 bosons and fermions (in light cone) associated with target space coordinates of the LST, and $`D=\widehat{c}_{eff}4`$ compact bosons and fermions, which model ”internal“ degrees of freedom, in the scenario of a large central charge. The spacetime directions will be taken to be either non-compact or compact with a large volume. The compactification radia of the internal CFT will remain fixed as we vary the temperature or the space volume. To keep the Hagedorn temperature fixed, the effective string tension scales as $`M_{eff}^2=\frac{M_{s}^{}{}_{}{}^{2}}{\widehat{c}_{eff}}`$.
The mass formula for closed strings is the familiar:
$`M^2`$ $`=`$ $`2(n+n^{})+\stackrel{}{L}^2+\stackrel{}{L^{}}^2`$
$`n^{}n`$ $`=`$ $`\stackrel{}{L}\stackrel{}{L^{}}`$ (4)
where we set the effective string tension to one for now, and recover it later. The vectors $`\stackrel{}{L},\stackrel{}{L^{}}`$ are momenta and windings for the compact fields, rescaled to set the radia to one. The Jacobian of this rescaling cancels in the expression for the free energy. The resulting Free energy is (approximating the discrete sum by an integral):
$$\beta F=V_5d^5p𝑑\stackrel{}{L}𝑑\stackrel{}{L^{}}𝑑n\rho (n)\rho (n+\stackrel{}{L}\stackrel{}{L^{}})e^{\beta \sqrt{p^2+M^2}}$$
(5)
with :
$$\rho (n)=\frac{1}{2^{(D+4)}}\frac{(D+4)^{(D+5)/4}}{n^{(D+7)/4}}e^{\pi \sqrt{n(D+4)}}$$
(6)
Any potential divergence in the free energy near the Hagedorn temperature comes from a saddle point where $`n\mathrm{}`$, so we can expand the integrand in that limit. This gives:
$$\beta F=V_5d^5pdn\frac{1}{4^{(D+4)}}\frac{(D+4)^{(D+5)/2}}{n^{(D+7)/2}}e^{2\pi \sqrt{n(D+4)}\beta \sqrt{p^2+4n}}\times $$
(7)
$$\times d\stackrel{}{L}d\stackrel{}{L^{}}e^{\frac{\pi \sqrt{D+4}}{2\sqrt{n}}\stackrel{}{L}\stackrel{}{L^{}}}e^{\frac{\beta }{4\sqrt{n}}(\stackrel{}{L}+\stackrel{}{L^{}})^2}$$
One can now perform the integrals over $`\stackrel{}{L},\stackrel{}{L^{}}`$. Since these integrals are not divergent as the temperature approaches the Hagedorn temperature, one can substitute $`\beta =\beta _H=\pi \sqrt{D+4}`$.
The expression for the free energy now becomes:
$$\beta F=V_5d^5p𝑑n\frac{1}{4^{(D+4)}}\frac{(D+4)^{(D+5)/2}}{n^{(D+7)/2}}e^{2\pi \sqrt{n(D+4)}\beta \sqrt{p^2+4n}}$$
(8)
$$\times \frac{1}{2^D}\left(\frac{8\sqrt{n}}{\sqrt{D+4}}\right)^D$$
which simplifies to the following expression (omitting D independent numerical factors):
$$\beta F=V_5d^5p𝑑n(D+4)^{5/2}n^{7/2}e^{2\pi \sqrt{n(D+4)}}e^{\beta \sqrt{p^2+4n}}$$
(9)
Now, define $`m^2=4n`$. We restore dimensions by using an effective string mass $`M_{eff}`$, and use $`D+4=\widehat{c}_{eff}`$. This gives, again up to numerical factors:
$$\beta F=V_5p^4𝑑p\frac{M_{eff}^5dm}{m^6}\widehat{c}_{eff}^{}{}_{}{}^{5/2}e^{\pi \sqrt{\widehat{c}_{eff}}\frac{m}{m_{eff}}}e^{\beta \sqrt{p^2+m^2}}.$$
(10)
We find that the free energy is extensive in $`V_5`$ and finite as $`\beta \beta _H`$, and that these statements are independent of the string model used (the precise value does depend on $`\widehat{c}_{eff}`$, which might eventually help determine $`\widehat{c}_{eff}`$ for LST, if other obstacles, described in this paper, are overcome first). Even though we calculated this behavior for a free worldsheet we expect it to hold for any critical string at $`g_s=0`$ string.
The same calculations above can be repeated to find the free energy in the case of a large compact spacetime directions. One finds the qualitative behavior:
$$\beta F\mathrm{log}(\beta \beta _H)$$
(11)
and independent of $`V`$. Again, this behavior is the same for the two toy models that we are using, and we expect it to hold for every $`\sigma `$-model.
As mentioned above, in the compact space case the behavior near the Hagedorn temperature is very different from the non-compact case, and in particular the free energy is no longer extensive. The difference between the non-compact and compact case arises due to the dynamics of long strings, which are allowed to wind on any compact directions. Below we compare these results with LST thermodynamics and find that even for a compact space the free energy is extensive. This is a first hint towards the dynamics of LST - it is in a phase which suppresses long strings.
One can now discuss other thermodynamics quantities, in the two cases. In the non-compact model the energy stays finite near the Hagedorn temperature . In the compact case the energy diverges as one approaches the Hagedorn temperature:
$$E=\frac{1}{\beta \beta _H}$$
(12)
In addition energy fluctuations are large, and behave as
$$\frac{<E^2><E>^2}{<E>^2}1.$$
(13)
The large fluctuations associated with the canonical ensemble can also be attributed to the behavior of the long strings. Near the Hagedorn temperature a finite fraction of the energy is stored in a single string. This makes it increasingly difficult to equilibrate the string gas as one approaches the Hagedorn temperature.
Next, we would like to compare this behavior to the near-Hagedorn behavior of LST. We do so by studying near extremal fivebrane solution, to which we turn now.
## 3 Near Extremal NS 5-branes
We review here the gravity background holographically dual to LST at a finite temperature. The zero temperature duality was discussed in . Roughly, as one approaches the boundary (UV) the solution asymptotes (exponentially fast) to the weakly coupled throat region of the CHS background . At the interior (IR) there are several interesting regimes, described by eleven dimensional supergravity, which however will not concern us here.
At a finite temperature the solution used in has to be replaced by a near extremal version. Though the exact solution is not known for all temperatures, it is becomes very simple in the limit that the energy density becomes large, since then the entire solution resides in the throat region and can be described using string theory. The near extremal fivebrane solution was written down in . In terms of appropriately rescaled quantities (after decoupling) the solution is (using the string frame):
$`ds^2=dt^2(1{\displaystyle \frac{\mu }{u^2}})+dx_i^2+{\displaystyle \frac{N}{u^2}}(du^2(1{\displaystyle \frac{\mu }{u^2}})^1+u^2d\mathrm{\Omega }_3)`$
$`e^{2\varphi }={\displaystyle \frac{N}{u^2}}`$ (14)
where $`u`$ is related to the radial coordinate away from the brane, and $`\mu `$ is the energy density above extremality. The index $`i=1,\mathrm{},5`$ corresponds to directions along the brane, and $`d\mathrm{\Omega }_3`$ is the metric of the unit 3-sphere.
The near extremal solution can be brought into a more familiar form by a redefinition of coordinates, $`u=\sqrt{\mu }\mathrm{cosh}r`$ (we also rescale the coordinates along the brane):
$`ds^2=N\left[dt^2\mathrm{tanh}^2(r)+dx_i^2+dr^2+d\mathrm{\Omega }_3\right]`$
$`e^{2\varphi }={\displaystyle \frac{N}{\mu \mathrm{cosh}^2(r)}}`$ (15)
We identify the background to be made out of the two dimensional black hole , a level $`N`$ supersymmetric WZW and an $`R^5`$ component<sup>2</sup><sup>2</sup>2The extremal solution is a linear dilaton direction times a WZW times $`R^6`$.. We use the Euclidean version of this black hole to describe the canonical ensemble. Accordingly, the Euclidean time variable is taken to be compact, with periodicity $`2\pi `$ in the conventions of (3).
The geometry of the two dimensional black hole is that of a semi-infinite cigar. The asymptotic circle closes at the tip, which is located at $`r=0`$. The string loop expansion is controlled by the value of the dilaton at the tip of the cigar:
$$g_{s,tip}^2=\frac{N}{\mu }$$
(16)
We see therefore that the string loop expansion is expansion in inverse energy density. This is the primary reason that string loop corrections in this background can change qualitative features of the thermodynamics.
It is worth noting that issues of decoupling are not relevant here . If the asymptotic string coupling is small enough, although not strictly taken to zero, there will still be a large range of energies in which the horizon of the black hole will be in the throat region and our analysis will apply<sup>3</sup><sup>3</sup>3The distance from the tip of the cigar to the asymptotic region is proportional to $`\sqrt{N}log(g_{s,tip}/g_{s,asymp})`$..
The tree level thermodynamics of the solution (3) is well known. The system has a fixed Hawking temperature, $`T_H\frac{M_s}{\sqrt{N}}`$, regardless of the energy. We calculate the tree level free energy in the appendix and confirm that $`F=0`$ for each value of the energy density $`\mu `$ (parameterized by the value of the dilaton at the tip, $`\varphi _0=\varphi (r=0)`$).
We find then that for fixed boundary conditions, there is a family of solutions to the Euclidean equations of motion. These solutions have degenerate action. It would seem then that the parameter $`\varphi _0`$ has to be summed over, and represents a flat direction in the path integral. Summing over $`\varphi _0`$ would result in a divergence. However, in it was argued that the mode changing the parameter $`\varphi _0`$ is a non-normalizable mode, and should not be summed over in the path integral.
The next step is to study the system at finite energy, namely study string loop corrections<sup>4</sup><sup>4</sup>4 Some aspects of the quantum corrections to the CGHS black hole were considered in .. This is done in the next section.
## 4 First Contributions to the Free Energy
### 4.1 Scalings of Correction Terms
We discuss now the corrections to the leading order action, and their effect on various thermodynamic quantities. Higher derivative terms in the ten dimensional action were discussed in . They include the famous $`R^4`$ term , and many other terms, related by supersymmetry. In order to organize the corrections to the thermodynamics, we keep the number of fivebranes, $`N`$, large and fixed, and arrange the possible corrections in powers of $`N`$.
We estimate therefore the scaling with $`N`$ of a general higher derivative term, and then discuss the known terms in . The metric as written above, equation (3), is proportional, in our notations, to $`N`$ (in the string frame). Both $`H`$ and $`g_{str}^2`$ are also proportional to $`N`$. Therefore, in the Einstein frame one has:
$`g_{..}\mu ^{1/4}N^{3/4}`$
$`H\mathrm{}N`$
$`e^\varphi N^{1/2}\mu ^{1/2}`$ (17)
where the dots indicate the position of spacetime indices. As a check of this scaling, all the terms in the leading order action
$$\sqrt{g}\left[R+H^2+(\varphi )^2\right]$$
(18)
scale uniformly like $`N^3\mu `$.
We now can discuss the scaling of the correction terms. The most general putative correction term has the schematic form:
$$\sqrt{g}e^{b\varphi }R_\mathrm{}.^k_{.}^{}{}_{}{}^{l}H_{\mathrm{}}^{}{}_{}{}^{t}g_{}^{..}{}_{}{}^{p}$$
(19)
where both $`b`$ and $`p`$ can be negative or positive, and $`l,k,t`$ are positive. We take negative $`p`$ to indicate metric with lower indices.
There are some constraints on the general term (19). First, the index structure of the above term demands:
$$4k+l+3t=2p$$
(20)
This guarantees that the term in the action is a scalar under general coordinate transformations.
In addition, suppose we are interested in a particular power $`s`$ of the energy $`\mu `$. The string loop expansion is controlled by the value of the dilaton in the tip of the cigar. Therefore, the integer $`s`$ is related to the genus of the string diagram giving rise to the particular correction term (s= 1-g,where g is the genus). For the general term above the power of $`\mu `$ is:
$$4s=52b+kp$$
(21)
We now can discuss the $`N`$ scaling of the general term (19). The power of $`N`$ of such a term is denoted by $`A`$, and is given by:
$$4A=15+2b+3k+4t3p$$
(22)
In order to simplify this expression for $`A`$ we can use the two relations above to write $`A`$ in terms of the positive quantities $`k,l,t`$, giving:
$$A=5sk\frac{l+t}{2}$$
(23)
For a given genus, the power of $`N`$ is determined by $`k,l,t`$ , and the two constraints above can be used to determine the needed $`p,b`$ to complete the correction term (19).
There are possible corrections to the action already at string tree level. These $`\alpha ^{}`$ correction terms are all of $`s=1`$ form. For example, the $`R^4`$ term has $`k=4`$, and therefore scales like $`\mu N^0`$. In fact, one can show that all the other leading order tree level corrections scale like $`\mu N^0`$.
We are more interested in the leading order corrections in the energy, coming from 1-loop terms in string theory. All such terms have $`s=0`$. For example, the one-loop $`R^4`$ term has $`k=4`$, and therefore scales as $`\mu ^0N`$. By a direct check, all other one loop terms scale similarly. This includes terms coming from up to 8 point scattering.
To summarize, the leading corrections in one loop are all of the same order, and scale like $`\mu ^0N`$. Compared to the leading order action they are suppressed by $`\mu N^2`$. We assume these corrections are non-zero. In this case they represent the leading order correction to the thermodynamics.
It is of some interest to calculate the coefficient of the corrections. In addition to confirming that it is non-zero, its sign has some significance in the thermodynamics. We elaborate on this point below.
### 4.2 Corrections to $`E(T)`$
In the presence of corrections to the leading order action, the solution (3) deforms slightly. The deformation is controlled by the size of the corrections $`\eta =\frac{1}{N^2\mu }`$. We are interested in the deformed solution asymptotically in the $`r`$ direction, so we can observe a shift in the temperature. One has schematically:
$$I=I_0+\eta I_1+\mathrm{}$$
(24)
where $`I_0`$ is the leading order action (54), and $`I_1`$ are the leading order corrections at one loop, discussed above.
To keep track of the scaling of corrections to the thermodynamics, we follow an outline of the calculation. An exact calculation would require a detailed knowledge of the perturbed action $`I_1`$. Rather, we are interested in extracting the dependence of the temperature correction on the parameter $`\eta `$.
We are interested in the change in the asymptotic value of $`G_{tt}=h^2(r)`$. The equation of motion for this metric component reads:
$$h^{\prime \prime }2h^{}\varphi ^{}=\eta \frac{\delta I_1}{\delta G_{tt}}$$
(25)
We define the right hand side of this equation to be the source, denoted by $`J_1`$.
Far enough from the tip, the background becomes approximately the linear dilaton vacuum. We define the small fluctuations:
$`h=1+\delta h`$
$`\varphi =r+log2+\varphi _0+\delta \varphi `$ (26)
Here we neglect the mixing with other small fluctuations. A complete calculation would require, of course, diagonalizing the complete matrix of quadratic fluctuations in this background. Clearly the complete calculation retains the scaling with $`\eta `$ demonstrated here.
One gets then the following equation for the small fluctuations:
$$\delta h^{\prime \prime }+2\delta h^{}=J$$
(27)
We note that the dilaton fluctuations drop out asymptotically, in the linear dilaton regime. The source in this equation $`J=J_0+\eta J_1`$ consists of two terms. The first term $`J_0`$ is independent of the energy density, and is of no interest to us. The second term $`J_1`$ is the one-loop correction defined above.
The solution to equation (27) is a linear combination of the two solutions of the homogeneous equations. Therefore:
$$\delta h=a_1+a_2e^{2r}$$
(28)
The two a-priori independent coefficients are determined by demanding regularity of the metric near the origin, as usual when determining the temperature. We are interested in the coefficient of the constant solution, $`a_1`$, which affects the asymptotic radius, and hence the temperature.
The coefficient $`a_1`$ is given by an overlap integral of the constant solution with the source $`J`$. The only dependence of $`a_1`$ on $`N,\mu `$ appears explicitly in the parameter $`\eta `$. Therefore $`a_1`$ can be written as a series expansion in the parameter $`\eta `$. The $`\eta `$ independent part only renormalizes the value of the Hagedorn temperature. The leading dependence on the energy density $`\mu `$ comes through:
$$\delta h(r)\eta \text{as}r\mathrm{}$$
(29)
The shift in $`h(r)`$ gives the following correction to the inverse temperature<sup>5</sup><sup>5</sup>5There is also a correction to the relation between the parameter $`\mu `$ and the energy density, but it has a subleading effect in the string loop expansion..
$$\frac{\beta \beta _H}{\beta _H}\eta $$
(30)
The coefficient of the correction becomes significant at this point. If the proportionality constant in the last equation is positive, the system approaches the Hagedorn temperature from below, at high energies. Asymptotically the system has a positive specific heat. Therefore the canonical ensemble is well-defined, though it has large fluctuations as explained above.
The scenario of a negative coefficient is very different: the system approaches the Hagedorn temperature from above, and has a negative specific heat. The canonical ensemble does not exist for high energies. We assume this is not the case and the coefficient is positive.
This gives the following modified relation between the temperature and the energy density:
$$\mu \frac{N^{5/2}M_{s}^{}{}_{}{}^{5}}{\beta \beta _H}$$
(31)
where $`M_{s}^{}{}_{}{}^{2}`$ is the spacetime string tension. Every string model will have to reproduce this expression, including the scaling with $`N`$.
### 4.3 Corrections to the Partition Function
The modified energy-temperature relation determines the thermodynamics of the system. Such a relation is expected to arise from an effective action of the form:
$$\beta Flog(\beta \beta _H)log(\mu )$$
(32)
This kind of contribution comes from a dependence of the effective action $`I=\beta F`$ on the logarithm of $`g_{str}`$, the string coupling at the tip. Since this dependence might seem unexpected, we demonstrate how it can be generated by the one loop perturbation.
The leading perturbation to the effective action is a sum of two terms. The first comes from the substituting the deformed solution in the original action, $`I_0`$. The second term comes from the action $`\eta I_1`$, evaluated on the unperturbed solution. Logarithmic contribution can arise from terms in the integrand that are constant in the linear dilaton regime. Since the length of the linear dilaton regime is proportional the $`\mathrm{log}(g_{str})`$, as shown above, this gives the desired effect.
Terms in the integrand which are constant in the linear dilaton regime are quite natural. We demonstrate here some possible sources for them. For example, the deformed solution for the metric component $`G_{tt}`$ is given by (28). This gives the following increment to the tree level effective action:
$$\delta I=\beta 𝑑re^{2r}\left[2\delta h^{\prime \prime }+4\delta h+\mathrm{}\right]$$
(33)
The integrand in this expression contains a constant term, where the decay of the fluctuation $`\delta h`$ is compensated for by the prefactor. There are also divergent terms, proportional to the temperature shift. This terms are regulated by considering differences in the action, as demonstrated in evaluating the leading term above. This procedure can also give rise to constant terms.
Another set of contributions to the free energy comes from a direct substitution of the original solution into the new action. These contributions can also contain constant terms in the integrand. To demonstrate this we need to consider the general correction term, evaluated in the linear dilaton background. In this background the only dependence on the radial coordinate is through the dilaton profile. Furthermore, derivatives of the dilaton are constant. The only dependence on the radial coordinate comes from the exponential of the dilaton. One loop terms come with the dilaton exponential raised to zeroth power. We find therefore that the corrections terms, evaluated in the linear dilaton regime, are all constant, and give rise to the desired logarithmic dependence.
Therefore, logarithmic behavior of the free energy is natural in the linear dilaton throat. This logarithmic behavior of leads to large fluctuations in the energy. One gets, as was discussed above:
$$\frac{<E^2><E>^2}{<E>^2}1$$
(34)
The coefficient here is of order $`N^0`$. We elaborate on the issue of large fluctuations below.
In a critical string theory, the one loop contribution results in effects that are quite different from the ones found above. At temperatures near the Hagedorn temperature there is a mode, winding around the Euclidean time circle, which becomes light. The effective description of this mode is discussed in . The situation is similar to the one described by Rohm , when considering Sherk-Schwarz supersymmetry breaking in string theory. The breaking by the boundary conditions on the circle results in a dilaton tadpole, and a one loop cosmological constant. The modified equations of motion have no longer a static solution.
We now argue that the effects discussed in are subleading to the ones discussed above. The effects of SUSY breaking by the boundary conditions on the Euclidean time circle can be computed in supergravity, as the system is at a low temperature compared to the string scale. The induced cosmological constant is of order $`(\frac{\beta }{l_p})^8`$ compared to the leading order action. In the present background this gives a suppression by $`N^3\mu `$. To compare, the effects discussed above are suppressed by $`N^2\mu `$, as compared to the leading order action.
## 5 Another Look at Large Fluctuations
The expression for the free energy, or the energy as a function of the temperature, reveals a problem in using the canonical ensemble when studying the thermodynamics of the system. The statistical fluctuations in the energy density are found to be:
$$(\mathrm{\Delta }\mu )^2N^0\mu ^2$$
(35)
As one approaches the Hagedorn temperature from below, (normalized) fluctuations in the energy remain finite. This makes the canonical ensemble not very useful when calculating thermodynamic averages. Physically, this stems from the attempt to keep a system with an exponential density of states in an equilibrium with a heat bath.
However, the canonical ensemble is still holographically dual to the Euclidean black hole, and it is useful to study these large fluctuations in the holographic description. At first sight, one might expect that the mode that changes the energy of the system becomes nearly massless at high energies, leading to the aforementioned fluctuations. However, as mentioned above,this mode is found to be a non-fluctuating mode . We present here an alternative picture for the large uncertainty of the energy. A related picture has appeared before in .
The system has a large number of normalizable, fluctuating modes. Their specturm was written in , and the behavior in the linear dilaton regime was studied in . The action for small fluctuations around a linear dilaton background is:
$$I=𝑑re^{2r}\left[G^{ij}_i\psi _j\psi \right]$$
(36)
Here $`G_{ij}`$ is the string frame metric, and $`\psi `$ is a typical fluctuation, for example a fluctuation in the metric polarized along the brane directions. We study the radial profile of the fluctuations. Our radial coordinate is related to as $`z=\sqrt{N}r`$. We also set the string tension to one for now. The modes $`\psi `$ were normalized to absorb any additional factor in the action.
The modes of the action are parameterized by $`s=N\omega ^2`$. Their asymptotic behavior is:
$$\psi _\pm (\omega )=\mathrm{exp}(\beta _\pm (\omega )ri\omega t)$$
(37)
with
$$\beta _\pm =1\pm \sqrt{1N\omega ^2}$$
(38)
For energies above the gap $`\omega _0=\frac{1}{\sqrt{N}}`$, the modes are normalizable. Their fluctuation scale is set by their action. We now proceed to estimate the fluctuations of a single normalizable mode.
Assume we are working with a finite cutoff, such that the length of the tube is finite and equals $`L`$. Imposing boundary conditions at $`r=L`$, and regularity conditions at the horizon, results in a discrete spectrum. Generically, this spectrum has spacing in energy $`\omega `$ in the order of $`\frac{1}{\sqrt{N}L}`$. This assumption is invalidated if the regularity at the horizon sets a complicated, $`\omega `$ dependent, relation between the two solutions $`\varphi _\pm `$. We assume this is not the case for generic normalizable modes.
We assume therefore that the boundary conditions pick a solution of the schematic form:
$$\psi (\omega )=e^r\mathrm{sin}(\sqrt{N}\omega r)\mathrm{sin}(\omega t)$$
(39)
The exact form of the trigonometric functions is immaterial.
For large energy modes the action is:
$$I(\omega )=𝑑r\omega ^2\mathrm{sin}^2(\sqrt{N}\omega r)L\omega ^2$$
(40)
The last step comes from averaging a rapid fluctuation over the length $`L`$, which is much larger than the period. We find therefore that with a finite cutoff $`L`$ , the typical fluctuations of the mode $`\psi (\omega )`$ are given by:
$$(\mathrm{\Delta }\psi )^2=\frac{1}{L\omega ^2}$$
(41)
Each mode $`\psi (\omega )`$ couples to the mode which moves the tip of the cigar by a certain form factor. This coupling is computed at tree level , and is therefore a general function of the form $`G(s)=G(N\omega ^2)`$. The fluctuations of the location of the tip, resulting from the fluctuations in all normalizable modes, are estimated to be:
$$(\mathrm{\Delta }a)^2=L\sqrt{N}𝑑\omega G(N\omega ^2)\frac{1}{L\omega ^2}$$
(42)
where we denote the radial position of the tip of the cigar (as defined by the zero of the metric component $`G_{tt}`$) by $`a`$, and its variance by $`(\mathrm{\Delta }a)^2`$. The normalization factor is needed when converting a discrete sum with spacing $`\frac{1}{L\sqrt{N}}`$ into an integral. We note that the dependence on the cutoff $`L`$ drops off, as it should.
The form factor $`G(s)`$ is unknown, but is expected to falloff at large frequencies, in order to yield a finite result. Impose a large frequency cutoff $`\mathrm{\Lambda }`$, and rescale $`k=\omega \sqrt{N}`$. This gives:
$$(\mathrm{\Delta }a)^2=N^{\mathrm{\Lambda }\sqrt{N}}\frac{dk}{k^2}G(k^2)$$
(43)
Under the assumption of convergence of the integral, the $`N`$ dependence disappears from the integral. The fluctuations in the location of the tip of the cigar are of order $`\sqrt{N}`$, in proper distance. Therefore, in our notations, fluctuations in the coordinate $`r`$ are of order one.
This translated, via the exponential relation between the radial location of the tip and the energy density, to the following:
$$(\mathrm{\Delta }\mu )^2N^0\mu ^2$$
(44)
We see therefore that under reasonable assumptions about the behavior of the normalizable modes, their collective effects result in large variance of the energy density. This fits with the expectations of the dual configuration, namely the canonical ensemble of LST.
## 6 The Coiled Phase of Strings
### 6.1 The Discrepancies
We would now like to compare more explicitly the free string model and the results obtained from the black hole. The free string results are either:
$$TT_H,Varbitrary,<E>=\frac{1}{\beta \beta _H}$$
(45)
$$TT_H,V=\mathrm{}\frac{<E>}{V}=finite.$$
Again, these relations are expected to be robust in critical string theory. Both relations are expected to be generic in the $`g_s=0`$ string.
On the other hand in “little string theory” the relation we obtained is:
$$TT_H,Varbitrary,\frac{<E>}{V}=\frac{1}{\beta \beta _H}$$
(46)
(and we have neglected factors of $`N`$, and dimensions are corrected using $`M_s`$). We clearly see that the free string model misses qualitative features of the thermodynamics. Next we would like to see whether there is some natural modification of the string dynamics that will go towards explaining the near-extremal thermodynamics.
Let us examine the non-compact partition function to see what modification can give us the correct result. Schematically the non-compact space free energy from section 2 is (after integrating over the momenta)
$$\beta FV_5\frac{dm}{m^{7/2}}e^{(\beta \beta _c)m}.$$
(47)
We would like to examine what modification to the partition function can change the power $`m^{7/2}`$ into $`m^1`$. Of course, there is no unique answer, but for later use we poit out that a possible modification is changing the energy by a $`ln(m)`$ term. More precisely the correct mass of a state is given by
$$m_{correct}=m\frac{5}{2\beta _c}ln(m)$$
(48)
In this case, we will obtain the correct behavior for the partition function.
The general motivation will be following. The mass of the string is roughly its length. We are therefore interested in a string interaction that reduces the energy of the string by $`ln(l)`$. We will see that there are such natural interactions, and we will focus on a self-intersection interaction which does that. Because the string now self-attracts it prefers to be coiled rather then large, solving the problem of dominance of long strings, and suggesting that we should consider a new phase made out of coiled strings.
### 6.2 General Considerations
We have seen that it is difficult to understand the near Hagedorn, i.e., high string excitation, behavior of LST in terms of a familiar string theory. In this section we would like to propose a different picture which, although we can not make precise, seems to remedy the situation. The upshot will be a new phase of strongly interacting strings with new qualitative features<sup>6</sup><sup>6</sup>6Another possibility is the existence of open strings in the system, this was suggested to us by O. Aharony. The thermodynamics of open string sectors was recently discussed in ..
Let us begin by highlighting two main properties that such a solution should have:
* First of all it should suppress long strings: As we have seen above, when space is compactified on a large torus, the contribution of long winding strings makes the free energy non-extensive. In order to obtain an extensive answer, we would like to suppress the contribution of these strings to $`S(E)`$.
Actually, this aspect is closely related to the non-gravitational nature of the theory. As emphasized by Susskind , one generally expects that in gravitational theories, objects that naively contribute many states to the entropy become large, such that their contribution does not violate the Bekenstein bound (For example, the ground state of the string grows when taking into account more and more oscillators ). Hence it is natural in critical string theory for highly excited strings to tend to be very large.
In our case, since there is no gravity and there are off-shell observables, we expect the string to prefer to be coiled rather then long at large excitation number. This observation is model-independent, and relies only general aspects of the relation between the number of degrees of freedom in a large volume and the size of highly excited objects.
* Secondly, the modification should have some reasonable space-time interpretation. Suppose we have a long string, or two long strings, then a reasonable space-time interpretation requires that when pieces of the strings are far away then the force between them will fall off at least as fast as the exchange of massless particles.
The conclusion is that we would like to modify the string models that we used before, which were some CFT on a worldsheet, by a strong attractive interaction. The new string is such that it can not be written in terms of a local worldsheet action (otherwise, we do not expect to evade the problems outlined above).
As an extreme, we can try and model the string with purely local interactions in space-time. In this case it would be a self-intersection interaction. We will argue that this is indeed what happens in the light cone frame. There are various kinds of self-intersection interactions that one can write down, but one expects that the most relevant attractive self-intersection interaction is simply such that one looses a fixed amount of energy for every self-intersection (in some regulated version of the worldsheet).
This is, of course, not precise for the LST. For example, two segments of strings that intersect at a very small angle such that they are almost parallel are almost BPS and hence there is no force between them, whereas anti-parallel strings attract each other. Hence the interaction is not as simple as we suggest. Note, however, that the average over these configurations certainly gives an average attractive interaction.
Two more comments are due. First, we have emphasized the effects of self-intersections but clearly if there is a strong intersection interaction, then each string in the thermal state interacts strongly with the background. How can we take into account this interaction ? We do not have the complete answer to this question, but we can suggest the following observation. One expects that if one uses a “mean field approximation” in which one computes the single string state statistics in a homogeneous background (encoded in the values of some order parameters), then the contribution to the energy of the string will be proportional to the length of the string. This means that the “mean field approximation” parameters of the background renormalize the string tension, but may not correct other terms. In particular the self-intersection term (which we will show gives the logarithmic correction in the exponent) is left uncorrected. The string tension may indeed be renormalized, but the Hagedorn temperature already measures the physical tension, after this renormalization has been taken into account (actually, a self intersection interaction also renormalizes the string tension, and the same argument applies to this renormalization as well).
Finally, in our discussion of the interaction we have used some kind of intuition about locality. This might be dangerous in a theory with T-duality in which space-time is not a well defined object - do we mean the initial, say, torus or its T-dual ? In our discussion what we mean is that in the large volume limit, there is an interaction with the properties discussed above, and in particular with approximate locality in the large space-time. There are many other non-local interactions, coming say from winding strings, but they decay as $`V\mathrm{}`$. This is true in critical string theory, and when we turn our attention to motivations Matrix description next, we will see that the behavior is very similar in LST.
### 6.3 Matrix model considerations
It is instructive to check whether this picture is consistent with the DLCQ description of LST . We are motivated by the fact that counting the single string degeneracy is simplest in light cone, and hence we would like to use Matrix theory to analyze it.
The Matrix model for LST on $`R^{5,1}`$ is the D1-D5 system which is a 1+1 dimensional sigma model on the ADHM moduli space. The latter is parameterized by two integers: $`N_0`$ which is the number of instantons and $`N`$ which is the rank of the $`U(N)`$ gauge group. In the Matrix interpretation $`N_0`$ measures the null momentum and $`N`$ the number of NS 5-branes. The Matrix model for the LST on $`R^{1,1}\times T^4`$ is the same D1-D5 system, where the D5 is now compactified on $`T^4`$, and we will focus on this model in our discussion. This model is a sigma model on a target space which is a deformation of $`T_{}^{4}{}_{}{}^{N_0N}/S_{N_0N}`$ where $`S_l`$ stands for the group of permutations of $`l`$ elements. The model is a deformation in the sense of not being at the solvable orbifold point. Rather, it is at a singular point in the moduli space of CFT. This point is the analogue of the $`\theta =0`$ point of $`R^4/Z_2`$ . Because of these singularities it is not clear how to analyze the model, but fortunately these singularities are rather mild as far as the current question is concerned. The reader is however warned that the analysis we present is rather speculative, but at this point we would like primarily to check consistency with the picture put forward above. The analysis also implies a new way of analyzing the D1-D5 system, and it will be interesting to explore it more rigorously further.
We would like to analyze the spectrum relevant in the Matrix framework (i.e., the correct energy scaling) and see whether it fits the line of thought explained in the previous subsection. If the model was the symmetric product then the analysis would have been straightforward. We could go to a picture of “long strings” by going to the twisted sector of the $`S_{N_0}`$ part of the symmetric group and obtain a string with the correct scaling of energy, tension $`M_s`$ and central charge $`4N`$ (or we could go to even longer strings with by using the remaining $`S_N`$ symmetry to a string with central charge 4, but we will not need this stage here). This would give us a Hagedorn density, but we saw before that it does not reproduce the thermodynamics correctly. This is not surprising since the CFT is not at the orbifold point.
Nevertheless, we would like to argue that such long strings are a good starting point. The reason for that is that the entropy of strings far away from the singularity is much larger then the entropy of strings at the singularity. The reason is that the total central charge of the CFT is $`4N_0N`$, whereas the effective strings at the singularity have a much smaller central charge. This can be extracted from in which the states at the singularities were described in terms of a 1+1 field theory written in terms of a $`U(N_0)`$ vector multiplet. The states at the singularity correspond to excitations along the flat directions of this non-abelian gauge theory, and hence their central charge is much smaller.
Hence, most of the states can be thought of as bulk states, i.e., these states for the most part are traveling away from the singularities of the sigma-model and are unaware of the singularity. In this sector one can try and go to long strings first and then consider the effect of the interaction. These long strings are the strings in spacetime and the interaction is precisely an interaction which is localized when the string intersects itself - i.e., it is an interaction of the same type that we were advocating above.
Hence the Matrix theory description lends support to the proposal that the modification is a strong self-intersection interaction. This analysis also suggests a new way of analyzing the D1-D5 system.
### 6.4 A Random Walk Model
Let us now try and estimate how a self-intersection interaction in light cone can change the partition function of a string theory. We will not do so precisely but rather point to a relation between this model and a certain model of self-attractive random walks. We will perform most of the computations in the latter model.
Modelling dynamics by a random walk is familiar from polymer physics . In the context of string theory Horowitz and Polchinski analyzed corrections to the size of an excited string state using such methods in the context of the black-hole/excited string correspondence principle <sup>7</sup><sup>7</sup>7There is, however, a difference in that there if the string moves in d+1 dimensions, then the random walk is in d dimensions. In our case we are motivated by the lightcone picture to use a local intersection interaction for the 4 transverse coordinates in light cone.. Although most of the analysis there is not directly relevant to our case <sup>8</sup><sup>8</sup>8An important part of interaction there is a long-range gravitational self attraction, which we do not have here. Also, it is not at all clear what might be the interpretation of the thermal scalar in our context., it is worth noting that they also find significant effects which lead to a contraction of string states.
Also, we will not be working precisely in the DLCQ of LST but rather in a simplified model in which there are 4 transverse (to the lightcone) coordinates and a simple self intersection interaction term, i.e., the energy decreases by some fixed amount when the string self-intersects in these coordinates. Other then this interaction term, one usually takes the action for the string to be its area. We will take a simplified model in which the weight of the string in the partition function is its length. This is so because we are counting physical states in lightcone, and we expect some equipartition between the kinetic and potential term. Hence there would still be a linear relation between the average size of the string as it fluctuates in some quantum state and the energy of that state (after all in lightcone the theory is a collection of oscillators). In any case, the skeptic reader can view the following analysis as a toy model which on the one hand reproduces some aspects of the partition function, and on the other hand is convenient for the analysis of the self-intersection interaction.
Under these assumption the partition function, without self intersections, at some inverse temperature $`\beta `$ will be
$$𝑑lf(l)e^{\beta _cl\beta l}$$
(49)
where $`l`$ is the length of the string (f(l) is a function which determined the degeneracy for a given length $`l`$). In relation to the usual quantization of strings we see that in order to match to the usual free oscillator picture we need to assume that at leading order $`l\sqrt{n}`$.
We would like to regulate this partition function in a way that captures its interpretation as strings in space-time. The way to do so is to approximate it as a sum over random walks, which will also enable us to more precisely define what we mean by the number of self-intersections. The length $`l`$ will now become a discrete variable and the partition function is a sum over all random walks with $`l`$ steps with weights $`e^{\beta l}`$. We will also assume the simplest form of a random walk, i.e., a cubic lattice with nearest neighbors jumps with equal probabilities.
Adding now the attractive interaction, the partition function becomes
$$\underset{randomwalks}{}e^{\beta (lg_{rw}J)}$$
(50)
where $`J`$ is the number of self-intersection
$$J=\underset{i,j=1,..,l}{}\delta _{w(i),w(j)},$$
(51)
and $`w(i)`$ denotes the location of the random walk at time $`i`$. $`g_{rw}`$ is some definite number which is determined by micro-physics and therefore we can not estimate.
We come now to the main purpose of the random walk model. Evaluating the number of self intersections for a simple random walk of length $`l`$ and regarding it as a correction to the energy of the random walk, we will obtain a logarithmic correction to the energy. Of course, without knowing the coefficient $`g_{rw}`$ above we will not be able check whether the logarithm in the exponent gives eventually the right power of $`l`$ needed to correct the degeneracy of the states, but it is interesting to get a logarithmic correction at all. The reason is that in any higher dimensions there are no logarithmic terms in the self-intersection number. Hence 4 coordinates in lightcone is the maximum number for which one would expect to see this type of universality class of strings. Fortunately this is precisely the maximal dimension of LST.
Finally let us explain how the logarithm comes about. We would like to evaluate $`<J>_l`$ on a closed loop of length $`l`$. We will use a Feynman diagram like technique (combinatoric computations as well as diagramatic techniques are described in . An explicit formula for self-intersection number, for open random walks in various dimensions, appears in Brydges and Slade ). Since we are evaluating a single insertion of $`J`$, the dominant contributions are random walks in which start at point 0, propagate to some point X after $`i`$ steps, return to that point after additional $`j`$ steps and then return to the origin after additional $`lij`$ steps. This diagrams contributes (The factor $`l^2`$ in front is due to normalizing by the probability that the random walk will return to the initial point - i.e., will be closed):
$$<J>_ll^2_{i+j<l}𝑑i𝑑j𝑑xe^{\frac{x_1^2}{i}}\frac{1}{i^2}\frac{1}{j^2}e^{\frac{x^2}{(lij)}}\frac{1}{(lji)^2}$$
(52)
where we have made the discrete sum over steps into an integral but require $`i,j>1`$ (we are not careful with numerical coefficients). It is straightforward to evaluate this integral and which yields in the large $`l`$ limit:
$$<J>_ll+2ln(l)+O(1)$$
(53)
The first term renormalized the string tension, and the second term is precisely what we wanted above - a change in the energy of the string that is proportional to the logarithm of its space-time length.
## 7 Conclusions
In this paper we discussed the thermodynamics of “little string theory” at high energies. Studying string loop corrections to the holographic dual we were able to go slightly off the Hagedorn temperature and probe the physics as we approach that critical temperature.
The main surprise is that this physics is inconsistent with a picture of free closed strings at high energies. The discrepancies are fundamental and generic to all free ($`g_s=0`$) string models. This provides a possible clue to the dynamics of the “little strings” at high energies.
The main discrepancy has to do with extensivity of the thermodynamics when the system is put at a finite volume. Free strings do not give extensive results, due to the dominance of long strings at high energies. The holographic duality teaches us that the “little strings” do give an extensive result. This is a first hint towards understanding dynamics in “little string theory”: the “little strings” do not form a single long string at high energies, rather they prefer to clump in small coils.
We argue for a simple model which reproduces this behavior. This model involves strings (in the lightcone) which interact only when they self-intersect. This suggestion is motivated by several observations. First, it seems that the interactions of the spacetime strings as seen in the DLCQ description yields an effective interaction of this form. Second, it is consistent with a local (on large scale) dynamics in spacetime. In addition, the interaction is different enough from the free string picture, so that a different qualitative behavior is possible.
Indeed, a simple analogous random walk model shows significant changes when an attractive self-intersection interaction is added. The changes look qualitatively similar to what we need: strings tend to become “coiled”, and the partition function gets corrections similar to the ones needed to reproduce the thermodynamics of “little strings theory”.
Clearly, more work is needed to establish this claim. A direct relation to a random walk model is needed if one is to reproduce exact, quantitative features of “little string theory”. Perhaps working in the DLCQ description, this relation can be clarified. We hope to return to these issues in the near future.
## Appendix: Vanishing of the Leading Order Free Energy
We calculate here the leading order thermodynamic quantities. We use the canonical ensemble, which corresponds to a Euclidean black hole configuration with a compact time direction. The basic quantity to calculate is the free energy, which is obtained from the value of the action on-shell .
We start with the ten dimensional Einstein frame action:
$$I=\frac{1}{16\pi G_{10}}\left[_Md^{10}x\sqrt{g}(R\frac{1}{2}(_\mu \varphi )^2\frac{1}{12}e^\varphi H^2)+2_MK\right]$$
(54)
To fix the ambiguity in the total action we need a prescription for boundary terms in (54). Following , we choose the total action such that fixing the values of the metric at infinity, but not its normal derivatives, is allowed. The standard boundary term, written above, involves the trace of the second fundamental form of the boundary.
We now transform to string frame metric $`G`$ by a conformal transformation, $`g_{..}=e^{\varphi /2}G_{..}`$, where $`H_{\mathrm{}}`$ remains unchanged in the transformation. Also, to compare to notations in , we write $`\mathrm{\Phi }=2\varphi `$
In addition to the standard string frame bulk action, there is an additional boundary term, which is, for the sphere at infinity:
$$(9/4)\sqrt{G}e^\mathrm{\Phi }G^{rr}_r\mathrm{\Phi }$$
(55)
where $`r`$ is the radial direction.
Now we perform dimensional reduction. Denote by $`V_5`$ the volume of the noncompact 5 dimensions, and by $`V_{sph}`$ the volume of the 3-sphere. We choose the following ansatz for the fields:
$`ds^2=\left(ds_{2}^{}{}_{}{}^{2}+dx_{i}^{}{}_{}{}^{2}+Nd\mathrm{\Omega }_3\right)`$
$`\mathrm{\Phi }=\mathrm{\Phi }(r)`$
$`H=NdV`$ (56)
Here $`i=1,\mathrm{},5`$ are flat directions along the brane, and $`d\mathrm{\Omega }_3`$ denotes the standard metric on $`S^3`$. $`dV`$ is the volume element on the 3-sphere. This ansatz covers both the extremal and non-extremal fivebrane solution, therefore it is suitable when calculating the effective action.
One then gets
$$I=\frac{V_5V_{sph}}{16\pi l_{s}^{}{}_{}{}^{8}}d^2x\sqrt{G}e^\mathrm{\Phi }\left[R+G^{\mu \nu }(_\mu \mathrm{\Phi })(_\nu \mathrm{\Phi })+2/N\right]$$
(57)
The prefactor is normalized in to be 1, by a shift of the dilaton $`\mathrm{\Phi }`$. We set the prefactor to 1, and recover it later.
We work in the ansatz for the two dimensional metric:
$$ds_{2}^{}{}_{}{}^{2}=N\left(dr^2+h^2(r)dt^2\right)$$
(58)
Here $`t`$ is Euclidean time, compactified with a period $`\beta l_s`$, where $`\beta `$ is dimensionless. The physical temperature is then $`\frac{1}{\sqrt{N}\beta l_s}`$.
The action should include also boundary terms discussed above. The boundary terms in the ten dimensional string action, written in the present ansatz, give the following terms:
$$I_1+I_2=\frac{2}{N}\beta e^\mathrm{\Phi }h^{}$$
(59)
Both are to be evaluated at the boundary at infinity only, where the Dirichlet boundary conditions have to be imposed.
The bulk two dimensional action (57) depends on second derivatives of the function $`h`$. In order to find the equations of motion we integrate by parts. This results in a boundary term we denote $`I_3`$:
$$I_3=\beta \sqrt{G}e^\mathrm{\Phi }\left[G^{\mu \nu }\mathrm{\Gamma }_{\mu \nu }^r\mathrm{\Gamma }_{r\mu }^\mu \right]$$
(60)
The bulk action is then:
$$I=\beta 𝑑re^\mathrm{\Phi }\left[2h^{}\mathrm{\Phi }^{}+4h+h\mathrm{\Phi }_{}^{}{}_{}{}^{2}\right]$$
(61)
where prime denoted derivative with respect to $`r`$.
The value of the action on shell is also a total derivative. This is easily shown using the equations of motion of the field $`\mathrm{\Phi }`$. Therefore the two dimensional bulk action (57), when evaluated on shell, equals the following terms, evaluated at the boundary:
$$I_3+I_4=\frac{1}{N}\beta e^\mathrm{\Phi }h\mathrm{\Phi }^{}$$
(62)
The boundary here includes the black hole horizon as well, due to the different origin of this boundary term. This is still to be regarded as a bulk action, though for this particular solution it reduces to boundary terms.
It is easy to show that the following solves the equations of motion:
$`h(r)=\mathrm{tanh}(r)`$
$`\mathrm{\Phi }(r)=2log\mathrm{cosh}(r)+a`$ (63)
Where $`a`$ is an arbitrary constant. This is the near extremal solution (3). The action of this solution should be compared with the vacuum configuration:
$`h(r)=1`$
$`\mathrm{\Phi }=2r+a2log2`$ (64)
The additive constant in $`r`$ is chosen so that the field $`\mathrm{\Phi }`$ has the same asymptotic behavior.
The effective action of the black hole configuration is found by a direct substitution in the above expressions. On general grounds one gets:
$$I_{total}=\beta F=\beta ES$$
(65)
Here $`F`$ is the free energy of the configuration, $`E`$ is the average energy and $`S`$ is the entropy. The first term $`\beta E`$ comes from boundary contribution, and the second one from a bulk contribution.
The bulk contribution comes from equation (62). This leads to a divergent result for the black hole (Appendix: Vanishing of the Leading Order Free Energy), but has to be compared to the action of the linear dilaton vacuum. This gives:
$$I_3+I_4=\frac{1}{N}\frac{V_5V_{sph}}{8\pi l_{s}^{}{}_{}{}^{8}}\beta e^a$$
(66)
where we have restored the prefactor set to 1 above.
The terms that get a contribution from the boundary at infinity are given in (59). They give the action:
$$I_1+I_2=\frac{1}{N}\frac{V_5V_{sph}}{8\pi l_{s}^{}{}_{}{}^{8}}\beta e^a$$
(67)
Combining the bove gives $`F=0`$ for every $`a`$. Also one can read off the energy density and the entropy:
$`\mu ={\displaystyle \frac{E}{V_5}}{\displaystyle \frac{N}{g^2}}{\displaystyle \frac{1}{l_{s}^{}{}_{}{}^{6}}}`$
$`S=\beta _HE\text{with}\beta _H=l_s\sqrt{N}`$ (68)
This gives the expected Hagedorn behavior, and the correct relation between the energy density and the string coupling, as obtained by other methods.
## 8 Acknowledgments
We ate happy to thank O. Aharony, M. Aizenman, C. Bachas, T. Banks, M. Douglas, J. Harvey, G. Horowitz, I. Klebanov, A. Lawrence, E. Martinec, G. Moore, B. Pioline, A. Rajaraman, S. Shenker, L. Susskind, E. Verlinde and H. verlinde for useful discussions. MR would like to thank the theory group at the University of Chicago for hospitality while this work was being completed. The work of MB is supported by NSF grant 98-02484. The work of MR is supported by DOE grant DOE-FG02-96ER40959.
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# Critical Resistivity along the Quantum Hall Liquid–Insulator Transition Line
\[
## Abstract
The critical resistivity $`\rho _{xx}^c`$ measured along the quantum Hall liquid–insulator transition line indicates a pronounced peak for a critical filling factor $`\nu _c1`$. The origin of this behavior, which marks the crossover from the high to low magnetic field regime in the phase diagram, is explained in the framework of classical transport in a puddle network model. The proposed scenario is also consistent with the behavior of the critical Hall resistance along the transition line. In addition, a functional form is suggested as a fit for isotherms $`\rho _{xx}(\nu )`$, to be compared with experimental data in the moderately high field regime.
\]
The rich phase diagram of a disordered two–dimensional electron system in a perpendicular magnetic field consists of two prominent types of phases, distinguishable by their transport properties at low temperatures ($`T`$). These include the quantum Hall (QH) liquid states, characterized by a quantized Hall resistivity $`\rho _{xy}`$ accompanied by $`\rho _{xx}0`$, and insulating states in which $`\rho _{xx}\mathrm{}`$. The transitions between adjacent phases have been extensively studied . More recently, there has been much progress in the experimental effort to map the phase diagram in the integer QH regime , and thus test the theoretically predicted structure . In particular, a transition line separating the conducting phases from the insulator was identified in the $`n`$$`B`$ plane (where $`n`$, the carrier density and $`B`$, the magnetic field, are the two independent control parameters). Along this line, in the moderately low $`B`$ regime, direct transitions are observed from $`N>1`$ QH states ($`N`$ integer) to the insulator . This behavior, and the overall structure of the diagram, are consistent with a numerical result based on the tight binding model for non–interacting electrons .
Comparison of the transport data along the conductor–insulator transition line ($`n_c(B_c)`$) in Refs. and indicates significant differences, particularly at low $`B`$ (where the GaAs hole system studied in Ref. exhibits a $`B=0`$ metal–insulator transition at a finite $`n_c`$). These can be attributed to sample dependent properties. However, the critical resistivity $`\rho _{xx}^c`$ measured along this transition line exhibits a striking common feature (observed in other samples as well ): it has a pronounced peak at $`\nu _c1`$. Here $`\nu _c`$ is the critical value of the Landau level (LL) filling factor ($`\nu n\varphi _0/B`$, with $`\varphi _0`$ the flux quantum), so that $`\nu _c1`$ marks the crossover from high to low $`B`$, where in the latter, more than one LL is significantly occupied. This crossover is also signified by the critical Hall resistivity $`\rho _{xy}^c`$ : for $`\nu _c1`$ it is quantized at $`\rho _{xy}^c=h/e^2`$, while above $`\nu _c1`$ it exhibits a classical Hall effect $`\rho _{xy}^ch/\nu _ce^2`$. $`\rho _{xx}^c`$ saturates to the universal value $`h/e^2`$ in the high $`B`$ limit where $`\nu _c1/2`$. However, the peak value of $`\rho _{xx}^c`$ is significantly elevated from $`h/e^2`$ and sample dependent ($`4h/e^2`$ in Ref. , $`2h/e^2`$ in Ref. ).
Since this puzzling behavior of $`\rho _{xx}^c`$ appears to be generic, it calls for a simple explanation of its physical origin. In this paper I show that the observed dependence of $`\rho _{xx}`$ on $`\nu `$, and in particular the peak of $`\rho _{xx}^c`$ at $`\nu _c1`$, follow from a phenomenological model for the transport, formerly introduced to interpret data in the high $`B`$ limit . The essential assumption underlying this model is that near the transition from a QH liquid to the insulator, the transport is dominated by narrow junctions connecting QH puddles; the puddles are defined by regions encircled by current carrying channels (edge states), and their typical size is assumed to be larger than a characteristic dephasing length. As a result, the longitudinal resistance of the sample is that of a classical random resistor network, and the Hall resistance is quantized provided all the QH puddles are at the same filling factor. This scenario successfully explains the experimentally observed quantized Hall insulator phenomenon , as well as the $`\nu `$–dependence of $`\rho _{xx}`$ . More recently, a similar model for the transport was shown to be consistent with the transport properties at low $`B`$, in samples indicating a $`B=0`$ metal–insulator transition .
It is useful to start by showing that in the high $`B`$ limit, this model implies $`\rho _{xx}^ch/e^2`$, regardless of the details of the disorder potential. In this limit only the lowest LL is occupied, hence all the QH puddles contain a single conducting channel. The longitudinal resistance of the junction $`j`$ between adjacent puddles in the network is therefore given by the Landauer–Büttiker formula
$$R_{xx}^j=\frac{h}{e^2}\frac{_j}{𝒯_j},$$
(1)
where $`𝒯_j=1_j`$ is the transmission probability across the junction; the Hall resistance $`R_{xy}^j=h/e^2`$ for any $`j`$. As proven in Ref. , the global resistivity tensor of the sample conveniently separates into Hall and longitudinal components: $`\rho _{xy}=h/e^2`$, and $`\rho _{xx}`$ is set by the random network of the resistors $`\{R_{xx}^j\}`$. Namely, $`\rho _{xx}R_{xx}^p`$, $`p`$ being the resistor at the percolation threshold of the distribution . The dependence of $`\rho _{xx}`$ on $`T`$ and $`\nu `$ is therefore determined by $`𝒯_p(T,\nu )`$. For sufficiently high $`T`$ such that quantum tunneling is neglected ,
$$𝒯_p(T,\nu )=f(ϵ_pϵ_F(\nu )),f(ϵ)\frac{1}{1+e^{ϵ/k_BT}};$$
(2)
here $`ϵ_F(\nu )`$ is the Fermi level, and $`ϵ_p`$ is the energy of the saddle point at which the junction $`p`$ is located. Employing Eq. (1), one obtains
$$\rho _{xx}\frac{h}{e^2}\frac{1𝒯_p}{𝒯_p}=\frac{h}{e^2}\mathrm{exp}\left[\frac{ϵ_pϵ_F(\nu )}{k_BT}\right].$$
(3)
It follows immediately that the critical filling factor $`\nu _c`$, defined such that $`\rho _{xx}(\nu _c)`$ is independent of $`T`$, obeys $`ϵ_F(\nu _c)=ϵ_p`$ so that the exponent in Eq. (3) vanishes. Consequently,
$$\rho _{xx}^c=\rho _{xx}(\nu _c)\frac{h}{e^2}.$$
(4)
Note that this argument can be extended to account for quantum tunneling across the junctions. Eq. (2) is then generalized to a Sommerfeld expansion of the form
$$𝒯_p(T,\nu )=f(ϵ_pϵ_F(\nu ))+\underset{n=1}{\overset{\mathrm{}}{}}\alpha _n(T)f^{(2n)}(ϵ_pϵ_F(\nu )),$$
(5)
where $`f^{(2n)}(x)`$ is the $`2n`$-th derivative of the Fermi function. This yield a $`T`$–independent expression for $`R_{xx}^p`$ provided all terms in the sum vanish. Since $`f^{(2n)}(x)=f^{(2n)}(x)`$, this occurs for $`ϵ_F(\nu )=ϵ_p`$, and once again $`\rho _{xx}^ch/e^2`$. It is also important to notice that the above argument for the universality of $`\rho _{xx}^c`$ is not violated in the case where $`\nu _c`$ is different from $`1/2`$ due to lack of particle–hole symmetry in the disorder potential. The only essential assumption is the neglect of mixing with higher LL. An interesting example for a remarkable universality of $`\rho _{xx}^c`$ at arbitrary $`\nu _c`$ is provided by the transitions from a fractional QH liquid to the insulator ; this observation will be discussed towards the end of the paper.
The above scenario changes once the magnetic field $`B`$ (at a given $`n`$) is reduced to a point where states in the second LL are occupied. In terms of the configuration of conducting channels in real space, this corresponds to the formation of closed loops of a second channel at the Fermi level inside the QH puddles. The overall filling factor of the system is given by
$$\nu =p[(1p_h)+\nu _hp_h],$$
(6)
where $`pA_l/A`$ is the total area fraction of the conducting liquid, $`p_hA_h/A_l`$ is the relative area fraction occupies by higher LL states and $`\nu _h2`$ their local filling factor. For a smoothly varrying disorder potential, $`p_h`$ can be estimated noting that $`p_h=r_2^2/r_1^2`$, where $`r_{1(2)}`$ is the average radius of the region surrounded by the first (second) LL channel in a typical QH puddle, and
$$\frac{1}{2}Kr_1^2\frac{1}{2}Kr_2^2+\mathrm{\Delta }(B);$$
(7)
$`\mathrm{\Delta }(B)`$ is either the cyclotron gap $`\mathrm{}\omega _c`$ or the Zeeman energy gap for spin–resolved Landau bands, and $`K`$ is the average curvature of the potential local minima. The critical filling factor $`\nu _c`$ for a transition to the insulator is then expressed as
$$\nu _c=p_c\left[1+\frac{(\nu _h1)r_2^2}{r_1^2}\right]p_c\left[\nu _h\frac{B_c}{B_0}\right],$$
(8)
where $`p_c`$ is the percolation threshold of the liquid, $`B_c`$ the critical field, and $`B_0`$ is a sample dependent parameter: $`B_0=m^{}cKA_p/he(\nu _h1)`$, $`B_0=KA_p/2\pi g\mu _B(\nu _h1)`$ for the spin–unresolved and spin–resolved cases, respectively ($`m^{}`$ is the effective mass of the carriers, $`A_p`$ is the typical puddle area at percolation). Note that Eq. (8) is valid for $`0r_2r_1`$, i.e. for $`p_c\nu _c\nu _hp_c`$; for a particle–hole symmetric potential $`p_c=1/2`$, hence the expression describes an interpolation between $`\nu _c1/2`$ (at high $`B`$) and $`\nu _c\nu _h/2`$. Neglecting the occupation of LL with $`N>2`$, $`\nu _h2`$ and the upper limit of this regime is $`\nu _c1`$. The critical density is then given by
$$n_c\frac{1}{\varphi _0}\left[B_c\frac{B_c^2}{2B_0}\right],$$
(9)
which is indeed a reasonable fit to the shape of the critical line $`n_c(B_c)`$ in Ref. , for $`\nu _c<1`$.
More interesting is the effect of occupied higher LL states on $`\rho _{xx}`$ as $`\nu _c1`$. In the regime $`\nu _c<1`$, corresponding to $`p_h<1`$, these states are localized, forming closed loops inside the QH puddles. A typical junction in the network of conducting channels is schematically depicted in Fig. 1. As $`p_h1`$, the localized high LL states are coupled to the primary conducting channels via tunneling through occasional weak links (dashed lines in the figure), and thus play the role of resonant scatterers which induce backscattering across the conducting regions. Denoting this backscattering rate close to the $`j`$-th junction by $`\mathrm{\Gamma }_j`$, and neglecting quantum interference between the two reflection paths , the effective transmission through the junction is reduced by a factor of $`(1\mathrm{\Gamma }_j)`$. The junction resistance becomes
$$R_{xx}^j=\frac{h}{e^2}\frac{1(1\mathrm{\Gamma }_j)𝒯_j}{(1\mathrm{\Gamma }_j)𝒯_j}.$$
(10)
The entire distribution of resistors in the random network is thus shifted up, and consequently the network resistance (determined as usual by the percolative resistor $`R_{xx}^p`$) is
$$\rho _{xx}\frac{h}{e^2}\frac{1(1\mathrm{\Gamma }_p)f(ϵ_pϵ_F(\nu ))}{(1\mathrm{\Gamma }_p)f(ϵ_pϵ_F(\nu ))}.$$
(11)
Here I assumed that the dependence of $`\mathrm{\Gamma }_p`$ on $`T`$ is negligible compared to the transmission across the junction, which involves an energy barrier. Eq. (11) clearly implies that the $`T`$ dependence disappears at a single value of the filling factor $`\nu _c`$ where $`ϵ_F(\nu _c)=ϵ_p`$, as in the high $`B`$ limit. However, the value of the critical resistivity thus defined is modified:
$$\rho _{xx}^c\frac{h}{e^2}\frac{1+\mathrm{\Gamma }_p}{1\mathrm{\Gamma }_p}>\frac{h}{e^2}.$$
(12)
For $`\nu _c1`$, the area fraction $`p_h`$ increases, implying that $`\mathrm{\Gamma }_p`$ is obviously an increasing function of $`\nu _c`$ (the explicite behavior of $`\mathrm{\Gamma }_p(\nu _c)`$ depends on details of the disorder potential). As a result, $`\rho _{xx}^c`$ grows with $`\nu _c`$ towards the limit case $`\nu _c1`$.
Another characteristic of the $`\nu _c<1`$ (high $`B`$) regime is that in the entire range of parameters where the classical transport scenario is valid, the Hall resistance maintains the quantized value $`\rho _{xy}=h/e^2`$. This follows from the fact that the Hall resistance of each junction as depicted in Fig. 1 is determined by the chemical potential difference between parallel right and left moving conducting channels, which is not affected by the modified transmission probability and yields $`R_{xy}^j=h/e^2`$ for any $`j`$.
The case where $`\nu _c1`$ marks a crossover to a low $`B`$ regime. From Eq. (8) (assuming $`p_c1/2`$, $`\nu _h2`$), $`\nu _c=1`$ corresponds to $`r_2r_1`$, which means that the average size of regions with a local filling fraction $`\nu 2`$ coincides with the average size of the conducting regions . Beyond this threshold value of $`\nu _c`$, higher LL states are delocalized, adding conducting channels to the percolating (or nearly percolating) network. As a consequence, the resistance of a typical junction $`R_{xx}^j`$ acquires a form similar to Eq. (10), reduced by an overall factor of $`1/N`$, where $`N`$ is the number of conducting channels that are clamped together. Note that $`\mathrm{\Gamma }_j`$ then describing scattering to partially occupied LL higher then $`N`$. The resulting critical value of the global longitudinal resistivity is then reduced compared to its value at $`\nu _c1`$, and is given by an expression of the form
$$\rho _{xx}^c\frac{h}{Ne^2}\frac{1+\mathrm{\Gamma }_p}{1\mathrm{\Gamma }_p}.$$
(13)
Here both $`N`$ and $`\mathrm{\Gamma }_p`$ depend on $`\nu _c`$ in a complicated manner, hence the actual behavior of $`\rho _{xx}^c`$ at $`\nu _c1`$ ($`B0)`$ is expected to be strongly sample dependent. Indeed, in Ref. $`\rho _{xx}^c0`$ in this limit, while in Ref. it saturates back to $`\rho _{xx}^ch/e^2`$. However, in proximity to the threshold $`\nu _c1`$ the generic trend is a reduction of $`\rho _{xx}^c`$ when $`\nu _c`$ is increased beyond unity. It therefore follows that this crossover between the high and low $`B`$ regimes is indicated by a peak in the critical resistance.
The increasing number of conducting channels ‘clamped’ together in the low $`B`$ regime also modifies the Hall resistance, since $`\{R_{xy}^j\}`$ are no longer identical. As argued in Ref. , the pure transverse component of the resistivity measured across the sample can be expressed as
$$\rho _{xy}=\frac{h}{e^2}\frac{\underset{i}{}(I_i/N_i)}{_iI_i},$$
(14)
where $`I_i`$ is the local current through the $`i`$-th QH puddle between the voltage probes, and $`N_i`$ its local filling factor. This yields a quantized value only provided all the $`N_i`$’s are identical. This is no longer the case for $`\nu >1`$: Eq. (14) generally describes a weighted average in which the weight of terms with $`N_i>1`$ is monotonically increasing with $`\nu `$. As a result $`\rho _{xy}h/\nu e^2`$, as in the classical Hall regime. Along the critical line, this implies $`\rho _{xy}^ch/\nu _ce^2`$. The transition points in the vicinity of integer values $`N>1`$ of $`\nu _c`$ correspond to the cases where a majority of the QH puddles consists of a unique filling factor $`N`$. They therefore characterize direct transitions from $`N>1`$ QH phases to the insulator.
The role of mixing with high LL does not only lead to violation of the universality of $`\rho _{xx}^c`$: naturally, the behavior of $`\rho _{xx}(\nu ,T)`$ across the transition is also modified. Indeed, in the high $`B`$ limit, isotherms $`\rho _{xx}`$ vs. $`\nu `$ can be fitted to a simple formula
$$\rho _{xx}=\frac{h}{e^2}\mathrm{exp}\left[\frac{\mathrm{\Delta }\nu }{\nu _0(T)}\right],$$
(15)
where $`\mathrm{\Delta }\nu =\nu \nu _c`$ and $`\nu _0(T)\alpha T+\beta `$. This formula exhibits a duality symmetry $`\rho _{xx}(\mathrm{\Delta }\nu )=\rho _{xx}^1(\mathrm{\Delta }\nu )`$, which is eventually violated in the lower $`B`$ regime. The analysis presented in this paper implies a specific modification of the above formula, which can be used as a test for the proposed transport scenario. As shown in Ref. , the simple exponential dependence on $`\mathrm{\Delta }\nu `$ in Eq. (15) can be derived from the energy dependent expression Eq. (3) (with an effective activation temperature) assuming a parabolic shape of the barriers, such that $`(ϵ_pϵ_F(\nu ))(\mathrm{\Delta }\nu )`$. Employing the same approximation in the moderately high $`B`$ regime, where the resistivity is given by Eq. (11), and relating the parameter $`\mathrm{\Gamma }_p`$ to $`\rho _{xx}^c`$ through Eq. (12), the isotherms acquire the form
$$\rho _{xx}(\nu )\frac{(\rho _{xx}^ch/e^2)}{2}+\frac{(\rho _{xx}^c+h/e^2)}{2}\mathrm{exp}\left[\frac{\mathrm{\Delta }\nu }{\nu _0(T)}\right].$$
(16)
In Fig. 2, a few such isotherms are plotted for a case where $`\rho _{xx}^c`$ is considerably elevated from $`h/e^2`$. Note that for $`\mathrm{\Delta }\nu >0`$, $`\mathrm{log}\rho _{xx}`$ vs. $`\mathrm{\Delta }\nu `$ curve upward compared to the straight line expected from the symmetric formula Eq. (15). A preliminary comparison with experimental data (e.g. Ref. ) confirms this behavior.
Finally, it is interesting to comment on the application of the transport model described here for transitions from the fractional QH state $`1/3`$ to the insulator. Experimentally , the transport data exhibit the same behavior as near the integer transitions: e.g., Eq. (15) is obeyed replacing $`\nu `$ by an effective filling factor of composite Fermions. This suggests that a similar model describes here the d.c. transport properties of composite Fermions. In particular, high LL of Fermions correspond to the hirarchy fractions $`2/5,3/7`$ etc. In analogy with the integer case, occupation of localized high LL Fermion states would imply the formation of incompressible regions in the disordered sample at these filling fractions. However, in contrast with the integer case, the energy gap supporting such regions is tiny. As a result, a mechanism for enhancing the resistance by backscattering through localized high LL states is expected to be practically absent. This possibly explains the observation that the duality symmetry in $`\rho _{xx}`$, as well as the universality of $`\rho _{xx}^c`$ that comes with it, are better manifested by data near the fractional transitions.
I thank A. Auerbach, M. Hilke, A. Kamenev, D. Shahar, and S. Sondhi for useful conversations. This work was supported by grant no. 96–00294 from the United States–Israel Binational Science Foundation (BSF), Jerusalem, Israel.
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# Automatic Closure of Invariant Linear Manifolds for Operator Algebras
## 1. CSL Algebras
###### Definition 1.
A projection $`P`$ in $``$ is *hyperatomic* if there are finitely many atoms $`A_1,\mathrm{},A_k`$ from $``$ such that $`P=E(A_1,\mathrm{},A_k)`$, the smallest projection in $``$ containing $`A_1,\mathrm{},A_k`$. We say that $`P`$ is *generated* by $`A_1,\mathrm{},A_k`$ if $`P=E(A_1,\mathrm{},A_k)`$. If each non-zero projection in $``$ is hyperatomic, we say that the lattice is *hyperatomic*.
###### Remark.
A projection $`P`$ is hyperatomic if, and only if, each ascending sequence $`F_1F_2F_3\mathrm{}`$ with $`P=F_n`$, $`F_n`$, is eventually constant. Indeed, assume that $`P`$ is generated by the atoms $`A_1,\mathrm{},A_k`$. Given an increasing sequence $`F_n`$ of elements of $``$ with $`P=F_n`$, there is, for each $`j=1,\mathrm{},n`$, a projection $`F_{n_j}`$ such that $`A_jF_{n_j}`$. If $`m\mathrm{max}\{n_jj=1,\mathrm{},k\}`$ then $`P=E(A_1,\mathrm{},A_k)F_mP`$; hence $`F_m=P`$, for all large $`m`$.
On the other hand, assume the ascending chain condition. Let $`\mathrm{SS}`$ be the set of atoms contained in $`P`$. If $`E(\mathrm{SS})<P`$, then $`PE(\mathrm{SS})`$ contains no atoms, i.e., no non-zero subinterval which is minimal. In this case it is easy to construct a strictly increasing sequence $`F_1<F_2<\mathrm{}`$ in $``$ such that $`P=F_n`$, contradicting the ascending chain condition. Thus we may assume that $`P`$ is generated by the atoms which it contains. If $`P`$ is not generated by finitely many atoms, let $`A_1,A_2,\mathrm{}`$ be an infinite sequence of atoms which generates $`P`$. Then $`E(A_1)E(A_1,A_2)\mathrm{}E(A_1,\mathrm{},A_k)<P`$, for all $`k`$, and $`P=_kE(A_1,A_2,\mathrm{},A_k)`$, again contradicting the ascending chain condition. Thus we may conclude that $`P`$ is generated by finitely many atoms; i.e., $`P`$ is hyperatomic.
The version of this remark appropriate to the whole lattice appeared in .
For $`x`$ and $`y`$ in $``$, $`xy^{}`$ denotes the rank one operator on $``$ given by $`zz,yx`$. Also, if $`P`$, then $`P_{}`$ denotes $`\{L:LP\}`$.
###### Lemma 2.
Let $`A`$ be an atom from $``$, let $`x`$ be a non-zero vector in $`A`$, let $`yE(A)`$ and put $`T:=x^2yx^{}`$. Then $`T\mathrm{Alg}`$, $`Tx=y`$ and $`TA=T`$.
###### Proof.
Since $`xA`$, we clearly have $`TA=T`$. Notice that if $`L`$ satisfies $`LE(A)`$, then $`AL=0`$; for otherwise $`AL`$, whence $`E(A)L`$. Therefore $`AE(A)_{}=0`$, so $`x(E(A)_{})^{}`$. Thus, from , we see that $`T\mathrm{Alg}`$. ∎
###### Proposition 3.
Let $``$ be a commutative subspace lattice on $``$ and let $`x`$. The following are equivalent:
1. $`x`$ is a closed vector for $`\mathrm{Alg}`$, i.e., $`\{\mathrm{Alg}x\}`$ is closed.
2. The projection onto $`[\mathrm{Alg}x]`$ is hyperatomic.
###### Proof.
Let $`P`$ be the orthogonal projection onto $`[\mathrm{Alg}x]`$. First assume that $`P`$ is not hyperatomic. Let $`F_1<F_2<\mathrm{}`$ be a strictly ascending sequence of projections in $``$ such that $`P=F_n`$. Let $`a_n`$ be a sequence of positive real numbers such that
(1)
$$\underset{n=1}{\overset{\mathrm{}}{}}n^2a_n^2<\mathrm{}.$$
Let $`k_n`$ be a sequence of positive integers such that for all $`n`$,
$`k_{n+1}k_n+2,\text{and}`$
$`(PF_{k_n})x=F_{k_n}^{}xa_n.`$
For each $`n`$, let $`y_n(F_{k_n+1}F_{k_n})`$ be a vector with $`y_n=na_n`$. By (1), the sum $`_{n=1}^{\mathrm{}}y_n`$ converges to an element $`yP`$.
For every $`n`$, we have
$$F_{k_n}^{}y(F_{k_n+1}F_{k_n})y=y_n=na_n\text{and}F_{k_n}^{}xa_n.$$
Hence for all $`n`$,
$$\frac{F_{k_n}^{}y}{F_{k_n}^{}x}\frac{na_n}{a_n}=n.$$
It follows from that $`yTx`$ for any $`T\mathrm{Alg}`$; i.e., $`y\{\mathrm{Alg}x\}`$. Thus $`\{\mathrm{Alg}x\}[\mathrm{Alg}x]`$, so $`x`$ is not a closed vector.
Now suppose that $`P`$ is hyperatomic. Let $`A_1,\mathrm{},A_n`$ be a finite set of atoms of $``$ such that
$$P=E(A_1,\mathrm{},A_n)=\{FA_jF,j=1,\mathrm{},n\}.$$
By deleting some atoms, if necessary, we may assume that $`A_1,\mathrm{},A_n`$ is a minimal set which generates $`P`$. Thus, if $`\mathrm{SS}`$ is any proper subset of $`\{A_1,\mathrm{},A_n\}`$, then $`E(\mathrm{SS})<P`$.
Let $`x_k=A_kx`$, $`k=1,\mathrm{},n`$ . Note that $`x_k0`$, for each $`k`$. (Otherwise, we have $`\overline{\{\mathrm{Alg}x\}}P=[\mathrm{Alg}x]`$ and $`x`$ is in $`E(A_1,\mathrm{},A_{k1},A_{k+1},\mathrm{},A_n)`$, a contradiction.) Now let $`y`$ be any vector in $`[\mathrm{Alg}x]=P`$. Then there exist vectors $`y_1,\mathrm{},y_n`$ (not necessarily unique) such that $`y_kE(A_k)`$, for each $`k`$, and $`y=y_1+\mathrm{}+y_n`$. By Lemma 2, there is, for each $`k`$, an element $`T_k\mathrm{Alg}`$ such that $`T_kx_k=y_k`$ and $`T_k=T_kA_k`$.
Let $`T=T_1+\mathrm{}+T_k`$. Clearly, $`T_kx_j=0`$ whenever $`kj`$. So $`Tx={\displaystyle \underset{k=1}{\overset{n}{}}}T_kx_k={\displaystyle \underset{k=1}{\overset{n}{}}}y_k=y`$. This shows that $`[\mathrm{Alg}x]\{\mathrm{Alg}x\}`$, and it follows that $`x`$ is a closed vector for $`\mathrm{Alg}`$. ∎
A von Neumann algebra $`𝔐`$ is also a CSL algebra precisely when $`𝔐^{}`$ is abelian. In this sense, one can view CSL algebras as generalizations of the von Neumann algebras with abelian commutant. The discussion in the introduction shows that every invariant manifold for a von Neumann algebra $`𝔐`$ with abelian commutant is closed exactly when $`𝔐`$ is finite dimensional, or equivalently, when $`\mathrm{Lat}(𝔐)`$ is hyperatomic. The next result, Theorem 4, generalizes this characterization to the class of all CSL algebras.
We also remark that Theorem 4 extends work of Froelich in . Motivated by operator theory, Froelich introduced the notions of strictly cyclic operator algebras (those for which there is $`x`$ with $`\{𝒜x\}=`$) and of strongly strictly cyclic operator algebras (those for which the compression of $`𝒜`$ to each invariant projection is strictly cyclic). He showed that strict cyclicity is equivalent to the ascending chain condition for the identity and the analogous result for strong strict cyclicity, essentially $`14`$ in Theorem 4.
###### Theorem 4.
Let $``$ be a commutative subspace lattice. The following statements are equivalent.
1. $``$ is a hyperatomic CSL.
2. Every invariant manifold for $`\mathrm{Alg}`$ is closed, i.e., $`\mathrm{Man}(\mathrm{Alg})=`$.
3. Every singly generated invariant manifold for $`\mathrm{Alg}`$ is closed.
4. Every element of $``$ is singly generated; that is, for $`P`$ there exists a vector $`xP`$ such that $`\{\mathrm{Alg}x\}=P`$.
5. Every invariant manifold for $`\mathrm{Alg}`$ is singly generated.
###### Proof.
($`12`$) This is proved in .
($`23`$) Obvious.
($`34`$) If $`P`$ is any element of $``$ then, since $``$ is separable, there is a vector $`x`$ such that $`P`$ is the projection onto $`[\mathrm{Alg}x]`$. But our hypothesis is that $`\{\mathrm{Alg}x\}`$ is already closed, so $`\{\mathrm{Alg}x\}=P`$.
($`41`$) Given $`P`$ we may find $`x`$ such that $`\{\mathrm{Alg}x\}=P`$, thus $`x`$ is a closed vector. By Proposition 3, $`P`$ is a hyperatomic projection. Since $`P`$ is arbitrary, every projection is hyperatomic and so $``$ is hyperatomic.
($`54`$) Obvious.
($`25`$) If $``$ is any invariant linear manifold, then $``$ is closed by hypothesis, so by the equivalence of (2) and (4), we see that there exists a vector $`x`$ such that $`=\{\mathrm{Alg}x\}`$. ∎
We next turn our attention to operator algebras $`𝒜`$ which are not weakly closed, but are subalgebras of CSL algebras. Theorem 4 shows that a necessary condition for $`\mathrm{Man}(𝒜)`$ to coincide with $`\mathrm{Lat}(𝒜)`$ is that $`𝒜`$ be a subalgebra contained inside the algebra of a hyperatomic CSL, so we will restrict our attention to this setting.
###### Definition 5.
Let $`()`$ and $`𝒞()`$ be operator algebras. We will say that $``$ is $`𝒞`$-transitive if $`\{x\}=\{𝒞x\}`$ for every $`x`$. Our primary interest is when $`𝒞`$. Notice that when $`𝒞=()`$, then the statement that $``$ is $`𝒞`$-transitive is simply the statement that $``$ is a transitive operator algebra.
The next proposition shows how $`\mathrm{Alg}`$-transitivity, closed vectors, and automatic closure of invariant manifolds for an algebra $`𝒜\mathrm{Alg}`$ are related under the mild assumption of a “local approximate unit,” i.e., when $`x[𝒜x]`$ for every $`x`$.
###### Proposition 6.
Let $``$ be a hyperatomic CSL and let $`𝒜()`$ be an algebra such that $`\mathrm{Lat}𝒜=`$ and such that $`x[𝒜x]`$, for every $`x`$. The following statements are equivalent.
1. $`𝒜`$ is $`\mathrm{Alg}`$-transitive.
2. Every vector $`x`$ is closed for $`𝒜`$.
3. Every invariant manifold for $`𝒜`$ is closed.
Moreover, when any of these conditions hold, every invariant manifold for $`𝒜`$ is singly generated; i.e., if $`\mathrm{Man}(𝒜)`$, then there exists $`x`$ such that $`=\{𝒜x\}`$.
Before beginning the proof, we remark that while every element of $``$ is singly generated as an $`\mathrm{Alg}`$ module, it is not a priori clear that every element of $``$ is singly generated as an $`𝒜`$ module.
###### Proof.
($`12`$) Let $`x`$. Since $``$ is hyperatomic, Theorem 4 shows that $`[\mathrm{Alg}x]=\{\mathrm{Alg}x\}`$ which, by assumption, is $`\{𝒜x\}`$. Thus statement (2) holds.
($`21`$) By assumption, for all $`x`$, we have $`\{𝒜x\}=[𝒜x]`$. Since $`\{\mathrm{Alg}x\}=[\mathrm{Alg}x]`$ is the smallest element of $``$ which contains $`x`$ and $`x\{𝒜x\}`$ by hypothesis, we see $`\{\mathrm{Alg}x\}\{𝒜x\}\{\mathrm{Alg}x\}`$. Thus, $`\{𝒜x\}=\{\mathrm{Alg}x\}`$, for all $`x`$.
($`13`$) Suppose that $`x_1,x_2`$. We claim that there is a vector $`x`$ in $``$ such that
(2)
$$\{𝒜x\}=\{𝒜x_1\}\{𝒜x_2\}$$
Write $`P_1`$ and $`P_2`$ for the projections onto $`\{\mathrm{Alg}x_1\}=\{𝒜x_1\}`$ and $`\{\mathrm{Alg}x_2\}=\{𝒜x_2\}`$. These projections are in $``$ and hence in $`\mathrm{Alg}`$. Let $`x=x_1+P_1^{}x_2`$. We will show that $`\{\mathrm{Alg}x\}=\{\mathrm{Alg}x_1\}\{\mathrm{Alg}x_2\}`$.
Since $`x_1=P_1x`$, it follows immediately that $`\{\mathrm{Alg}x_1\}\{\mathrm{Alg}x\}`$. Now let $`y\{\mathrm{Alg}x_2\}\{\mathrm{Alg}x_1\}^{}`$. Then there is $`T\mathrm{Alg}`$ such that $`y=Tx_2`$. Since $`x_2=P_1^{}x_2+P_1x_2`$, we have
$$y=Tx_2=TP_1^{}x_2+TP_1x_2=TP_1^{}x_2+P_1TP_1x_2.$$
Since $`P_1^{}y=y`$,
$$y=P_1^{}TP_1^{}x_2=P_1^{}TP_1^{}(x_1+P_1^{}x_2)=P_1^{}TP_1^{}x.$$
This shows that $`\{\mathrm{Alg}x_2\}\{\mathrm{Alg}x_1\}^{}\{\mathrm{Alg}x\}`$. Combining this with $`\{\mathrm{Alg}x_1\}\{\mathrm{Alg}x\}`$ gives $`\{\mathrm{Alg}x_1\}\{\mathrm{Alg}x_2\}\{\mathrm{Alg}x\}`$. The reverse inequality follows from the fact that $`x=x_1+P_1^{}x_2`$ and the claim is verified.
Now let $``$ be an arbitrary invariant linear manifold for $`𝒜`$. We need to show that $``$ is closed. Let $`Q`$ be the projection onto the closure of $``$. Clearly $`Q`$ and hence is a hyperatomic projection. Let $`A_1,\mathrm{},A_n`$ be a family of atoms of $``$ such that $`Q=E(A_1,\mathrm{},A_n)`$. Then $`Q=E(A_1)E(A_2)\mathrm{}E(A_n)`$. Notice that if $`x_j`$ is a non-zero vector in $`A_j`$, then $`[\mathrm{Alg}x_j]=\{\mathrm{Alg}x_j\}=\{𝒜x_j\}=E(A_j)`$. Inductively applying (2), we see that there exists $`y`$ such that $`=\{𝒜y\}`$, whence $``$ is singly generated. Since $`𝒜`$ is $`\mathrm{Alg}`$-transitive, $`y`$ is a closed vector, whence $``$ is closed.
($`32`$) Obvious.
It remains to show that when $`\mathrm{Man}(𝒜)=`$, then every invariant manifold for $`𝒜`$ is singly generated. Let $``$ be an invariant linear manifold for $`𝒜`$. Then by hypothesis, the orthogonal projection $`Q`$ onto $``$ belongs to $``$, hence there is a vector $`x`$ such that $`=\{\mathrm{Alg}x\}`$. Clearly $`x`$. Now $`x[𝒜x]=\{𝒜x\}`$ and since $`\{\mathrm{Alg}x\}`$ is the smallest element of $``$ containing $`x`$, we conclude that $`\{𝒜x\}\{\mathrm{Alg}x\}\{𝒜x\}`$. ∎
The following theorem requires the Kadison transitivity theorem for its proof and is (partially) an extension of that theorem.
###### Theorem 7.
Let $``$ be a hyperatomic CSL. Suppose that $`𝒜\mathrm{Alg}`$ is a norm closed operator algebra such that $`\overline{𝒜}^{\text{wot}}=\mathrm{Alg}`$. Assume that $`E𝒜F𝒜`$ for all atoms $`E`$ and $`F`$ of $``$ and that $`E𝒜E`$ is a $`\text{C}^{}`$-algebra for each atom. Then $`\mathrm{Man}(𝒜)=\mathrm{Lat}(𝒜)=`$, every element of $`\mathrm{Man}(𝒜)`$ is singly generated, and $`𝒜`$ is $`\mathrm{Alg}`$-transitive.
###### Proof.
Observe that $`E𝒜E`$ is a $`\text{C}^{}`$-algebra which is weakly dense in $`E\mathrm{Alg}E=(E)`$; thus $`E𝒜E`$ is an irreducible $`\text{C}^{}`$–subalgebra of $`(E)`$.
We first assume that the identity operator $`I`$ is generated by a single atom $`E_0`$ of $``$. We shall prove that the invariant manifold for $`𝒜`$ generated by a unit vector in $`E_0`$ is all of $``$. So fix a unit vector $`\xi E_0`$ and let $`x`$ be any vector. Let $`\{Q_n\}_{n=0}^{\mathrm{}}`$ be a sequence of projections in $`^{\prime \prime }`$ such that each $`Q_n`$ is a finite sum of atoms of $``$, $`E_0Q_0`$, $`_{n=0}^{\mathrm{}}Q_n=I`$ and $`_{n=1}^{\mathrm{}}Q_nx<\mathrm{}`$.
Fix $`n0`$. Write $`Q_n=_{j=1}^{k_n}E_{n,j}`$ as a finite sum of atoms of $``$ and let $`x_n=Q_nx`$. Since $`E_{n,j}\mathrm{Alg}E_0=(E_0,E_{n,j})`$, and $`E_{n,j}𝒜E_0`$ is weakly dense in $`E_{n,j}\mathrm{Alg}E_0`$, we may find a norm one operator $`Y_{n,j}𝒜`$ such that $`Y_{n,j}=E_{n,j}Y_{n,j}E_0`$. Hence we may find a unit vector $`u_{n,j}E_0`$ such that $`Y_{n,j}u_{n,j}>1/2`$. By Kadison’s transitivity theorem, there exist unitary operators $`Z_{n,j}E_{n,j}𝒜E_{n,j}`$ and $`W_{n,j}E_0𝒜E_0`$ such that
$$\frac{E_{n,j}x_n}{Y_{n,j}u_{n,j}}Z_{n,j}Y_{n,j}u_{n,j}=E_{n,j}x_n\text{ and }u_{n,j}=W_{n,j}\xi .$$
Writing
$$A_{n,j}=\frac{E_{n,j}x_n}{Y_{n,j}u_{n,j}}Z_{n,j}Y_{n,j}W_{n,j},$$
we see that
$$A_{n,j}𝒜,A_{n,j}<2E_{n,j}x_n,\text{ and }A_{n,j}\xi =E_{n,j}x_n.$$
Therefore, if $`B_n=_{j=1}^{k_n}A_{n,j}`$, we find $`B_n𝒜`$ and $`B_n\xi =x_n`$. Moreover, since $`E_{n,j}A_{n,j}=A_{n,j}`$, we find that for any $`\eta `$, $`B_n\eta ^2=_{j=1}^{k_n}E_{n,j}A_{n,j}\eta ^2`$, so
$$B_n\left\{\underset{j=1}{\overset{k_n}{}}A_{n,j}^2\right\}^{1/2}<2x_n.$$
Notice also that $`B_n=B_nE_0`$ by construction.
The fact that $`_{n=1}^{\mathrm{}}x_n<\mathrm{}`$ shows that the series $`_{n=0}^{\mathrm{}}B_n`$ converges uniformly to an element $`B𝒜`$. Clearly, $`B\xi =_{n=0}^{\mathrm{}}B_n\xi =_{n=0}^{\mathrm{}}x_n=x`$. Thus we have shown that the invariant manifold generated by $`\xi `$ is all of $``$.
Furthermore, notice that our construction shows the following:
1. $`B=BE_0`$;
2. if $`E`$ is an atom of $``$ such that $`Ex=0`$, then $`EB=0`$; and
3. $`B=lim_n\mathrm{}R_n`$, where for each $`n`$, $`R_n=_{j=1}^{p_n}C_{n,j}`$ is a finite sum of elements $`C_{n,j}𝒜`$ which satisfy $`C_{n,j}=E_{n,j}C_{n,j}F_{n,j}`$ for some atoms $`E_{n,j}`$ and $`F_{n,j}`$ of $``$.
Returning to the general case, if $`E`$ is any atom from $``$, we may compress to $`P(E)`$ (i.e., replace $`𝒜`$ by $`P(E)𝒜P(E)`$ acting on $`P(E)`$) and apply the argument above to obtain the following: if $`\xi `$ is any non-zero vector in $`E`$ and if $`xP(E)`$, then there is $`B𝒜`$ such that $`B\xi =x`$, $`B=BE`$, and $`Fx=0`$ implies $`FB=0`$ for all atoms $`F`$. (There is one delicate point: our hypotheses do not guarantee that $`P(E)𝒜P(E)𝒜`$, but in the construction of $`B`$, $`B`$ is a norm limit of elements which are finite sums of elements of the form $`F_1XF_2`$ with $`F_2`$ and $`F_2`$ atoms of $`𝒜`$. Such elements are in $`𝒜`$, by our hypotheses.)
Now let $``$ be an invariant linear manifold under $`𝒜`$. Let $`P`$ be the projection onto $`\overline{}`$. Then $`P`$ is invariant under $`𝒜`$ and, hence, under $`\mathrm{Alg}`$. So, $`P`$.
Let $`E_1,\mathrm{},E_n`$ be independent atoms which generate $`P`$. So, $`P=P(E_1,\mathrm{}E_n)`$ and $`E_i\mathrm{Alg}E_j=0`$ whenever $`ij`$. There is a vector $`x`$ such that $`E_ix0`$ for all $`i`$. (Let $`y_iE_i`$ with $`y_i=1`$ and approximate $`y_i`$ in norm by an element of $``$.)
Clearly $`𝒜x`$. We will prove that $`P𝒜x`$; this implies that $`=P`$, whence $``$ is closed and singly generated. Let $`yP`$ be arbitrary. Write $`y=y_1+\mathrm{}+y_n`$, where $`y_iP(E_i)`$ for each $`i`$. This can be done since $`P=_iP(E_i)`$.
For each $`i`$, there is an element $`B_i𝒜`$ such that $`B_ix_i=y_i`$, $`B_i=B_iE_i`$, and $`Fy_i=0`$ implies $`FB_i=0`$, for all atoms $`F`$. Let $`B=B_1+\mathrm{}+B_n`$. Then
$`Bx`$ $`=B_1x+\mathrm{}+B_nx`$
$`=B_1E_1x+\mathrm{}+B_nE_nx`$
$`=B_1x_1+\mathrm{}+B_nx_n`$
$`=y_1+\mathrm{}+y_n=y.`$
Thus, $`y𝒜`$ and $`P𝒜x`$.
Finally, since $`𝒜`$ is weakly dense in $`\mathrm{Alg}`$, for every $`y`$ we have $`\{𝒜y\}=[𝒜y]=[\mathrm{Alg}y]`$. Since $``$ is hyperatomic, we have $`\{\mathrm{Alg}y\}=[\mathrm{Alg}y]`$, so $`\{𝒜y\}=\{\mathrm{Alg}y\}`$ for every $`y`$. It follows that $`𝒜`$ is $`\mathrm{Alg}`$-transitive. ∎
This theorem implies immediately a result tacit in the proof of the automatic closure theorem in :
###### Corollary 8.
Let $`𝒦`$ be the algebra of compact operators and suppose $``$ is a hyperatomic CSL. Then every invariant linear manifold for $`𝒦\mathrm{Alg}`$ is closed.
## 2. TAF algebras
We turn now to representations of strongly maximal triangular AF (TAF) algebras. These are subalgebras of AF $`\text{C}^{}`$-algebras arising as limits of triangular digraph algebras and have been extensively studied; see, for example, . If $`𝒜`$ is a closed subalgebra of an AF $`\text{C}^{}`$-algebra $`𝒞`$, then $`𝒜`$ is triangular AF or TAF if $`𝒜𝒜^{}`$ is a canonical masa in $`𝒞`$. A masa $`𝒟`$ is a canonical masa in $`𝒞`$ if the closed span of $`N_𝒟(𝒞)`$ is $`𝒞`$, where
$$N_𝒟(𝒞)=\{f𝒞:f\text{ is a partial isometry, }fdf^{},f^{}df𝒟\text{ for }d𝒟\}.$$
A triangular algebra $`𝒜`$ is strongly maximal if $`\overline{𝒜+𝒜^{}}=𝒞`$.
Let $`𝒜`$ be a strongly maximal triangular AF subalgebra of the AF $`\text{C}^{}`$-algebra $`𝒞`$ with $`𝒟=𝒜𝒜^{}`$ a canonical masa in $`𝒞`$. For reasons we will explain momentarily, we consider representations $`\pi :𝒜()`$ satisfying the following conditions:
1. $`\pi `$ is the restriction to $`𝒜`$ of a $``$-representation $`\rho `$ of $`𝒞`$ on $``$;
2. $`\pi (𝒟)`$ is weakly dense in a masa in $`()`$; and
3. $`\mathrm{Lat}(\pi (𝒜))`$ has order type $``$ and is multiplicity free.
Representations satisfying the first two conditions are called *masa preserving* \[21, p. 130\]. Since $`\mathrm{Lat}(\rho (𝒞))\mathrm{Lat}(\pi (𝒜))\mathrm{Lat}(\pi (𝒜^{}))=\{0,I\}`$, the $``$-representation $`\rho `$ is necessarily irreducible. We will occasionally call a representation which satisfies conditions 1, 2, and 3 an *admissible representation*.
If $`𝒜`$ is non-unital, so is $`𝒞`$. By $`𝒜^+`$ we mean the obvious subalgebra of the unitization $`𝒞^+`$ of $`𝒞`$, and it is easy to see that $`𝒜^+`$ is a strongly maximal TAF subalgebra of $`𝒞^+`$. Since $`\mathrm{Man}\pi (𝒜)=\mathrm{Man}(\pi (𝒜^+)),`$ we lose no generality by assuming that all algebras and representations are unital, and thus we make this assumption in the sequel.
The simplest example of a representation satisfying the three conditions above is the Smith representation of the standard embedding algebra acting on $`\mathrm{}^2()`$ \[21, Example I.2\]. In fact, for standard embedding algebras, \[21, Theorem III.2.1\] shows that $`\mathrm{Lat}(\pi (𝒜))`$ is multiplicity free for representations $`\pi `$ satisfying all the other conditions above.
A more general class of strongly maximal TAF algebras, the $``$-analytic algebras considered in , also admit representations of this form. However, not all strongly maximal TAF algebras have representations satisfying conditions 1, 2, and 3; for example, the refinement embedding algebras (see ) have no such representations. Further, for a masa preserving representation of a refinement embedding algebra, there is a non-closed invariant linear manifold.
We have previously observed that the second condition implies that $`\pi (𝒜)`$ is $`\sigma `$-weakly dense in the CSL algebra $`\mathrm{Alg}\mathrm{Lat}(\pi (𝒜))`$. However, since $`𝒜`$ is a strongly maximal TAF algebra, more is true: \[21, Proposition 0.1\] shows that the second condition implies that $`\mathrm{Lat}(\pi (𝒜))`$ is a nest. (A nest is a totally ordered CSL.) Moreover, for many of the examples in , $`\mathrm{Alg}\mathrm{Lat}(\pi (𝒜))`$ is multiplicity free. If every invariant manifold for $`\pi (𝒜)`$ is closed, then necessarily the nest $`\mathrm{Lat}(\pi (𝒜))`$ is hyperatomic.
Furthermore, if $`𝒜`$ is a $``$-analytic subalgebra of a simple AF $`\text{C}^{}`$–algebra and if $`\pi `$ is an irreducible representation of $`C^{}(𝒜)`$ which satisfies condition 2, then by \[21, Proposition III.3.2\] $`\mathrm{Lat}\pi (𝒜)`$ is a nest whose order type is a subset of the integers. Since a nest is hyperatomic if, and only if, the complementary nest is well-ordered, the hyperatomic nests with order type a subset of the integers are just the finite nests and the nests of order type $``$. Automatic closure for invariant manifolds is trivial when $`\mathrm{Lat}\pi (𝒜)`$ is a multiplicity free finite nest. If the nest is finite but not multiplicity free (the nest may be the trivial nest $`\{0,I\}`$, for example), the automatic closure question is open. For nests with order type $``$, Theorem 10 below gives an affirmative answer to the question.
Some motivation for, in effect, restricting to irreducible representations can be found in the following fact, although it does not reduce the study of masa preserving representations to the study of irreducible masa preserving representations.
###### Lemma 9.
Let $`\pi `$ be a representation of $`𝒞`$ such that $`\pi (𝒟)`$ is weakly dense in a masa in $`()`$. Every invariant linear manifold for $`\pi (𝒜)`$ is closed if, and only if, $`\pi `$ decomposes as a direct sum of finitely many irreducible representations $`\pi =_{i=1}^n\pi _i`$ and each invariant linear manifold for $`\pi _i(𝒜)`$ is closed $`(i=1,\mathrm{}n)`$.
###### Proof.
If $`\mathrm{Man}(\pi (𝒜))=\mathrm{Lat}(\pi (𝒜))`$, then since every invariant manifold for $`\pi (𝒞)`$ is also an invariant manifold for $`\pi (𝒜)`$, the discussion in the introduction shows that $`\pi `$ decomposes as required. Then every linear manifold invariant for $`\pi _i(𝒜)`$ is also invariant for $`\pi (𝒜)`$. Conversely, since $`\pi (𝒟)^{\prime \prime }`$ is a masa in $`()`$, every invariant manifold $``$ for $`\pi (𝒜)`$ decomposes as finite orthogonal sum of invariant manifolds for $`\pi _i(𝒜)`$, whence $``$ is closed. ∎
We can now state the main theorem of this section.
###### Theorem 10.
Let $`𝒜`$ be a strongly maximal triangular subalgebra of an AF $`\text{C}^{}`$-algebra $`𝒞`$. If $`\pi :𝒜()`$ is a masa preserving, order type $``$, multiplicity free representation, then every invariant linear manifold for $`\pi (𝒜)`$ is closed and singly generated.
While the proofs of Theorem 10 and of Theorem 7 employ similar methods, this theorem is not subsumed by Theorem 7 as $`\pi (𝒜)`$ is not a bimodule over the algebra generated by the atoms of the nest.
To prove Theorem 10, we need to describe admissible representations in terms of coordinates. The full development of such coordinates is technical, and the reader is referred to for more complete treatments. Associated to each AF $`\text{C}^{}`$-algebra $`𝒞`$ there is a unique AF groupoid $`𝔊`$ so that the $`\text{C}^{}`$-algebra of $`𝔊`$, $`C^{}(𝔊)`$, and $`𝒞`$ are isomorphic as $`\text{C}^{}`$-algebras. The elements of $`C^{}(𝔊)`$ can be identified with continuous functions on $`𝔊`$. With this identification, $`C(𝔊_0)`$ embeds in $`C^{}(𝔊)`$ and is a canonical masa in $`C^{}(𝔊)`$. In particular, we may identify $`𝒟`$ with $`C(𝔊_0)`$ and $`𝒞`$ with $`C^{}(𝔊)`$. Given a unit $`e𝔊_0`$, its orbit is the set
$$[e]:=\{f𝔊_0:\text{ for some }x𝔊,e=x^1x\text{ and }xx^1=f\}.$$
Given a triangular algebra $`𝒜`$ with $`𝒟𝒜𝒞`$, there is an anti-symmetric subset of $`𝔊`$ containing $`𝔊_0`$, denoted $`\mathrm{Spec}(𝒜)`$, so that $`𝒜`$ is isomorphic to $`\{fC^{}(𝔊):\mathrm{supp}f\mathrm{Spec}(𝒜)\}`$. If $`𝒜`$ is strongly maximal, then $`\mathrm{Spec}(𝒜)`$ totally orders each orbit in $`𝔊_0`$. Similar coordinates can be defined for other groupoid $`\text{C}^{}`$-algebras; see .
By Theorem II.1.1 in , each representation satisfying conditions 1 and 2 is unitarily equivalent to a representation of the type constructed below. Recall that a 1-cocycle is a groupoid homomorphism $`\alpha :𝔊G`$, where $`G`$ is an abelian group; for us $`G`$ is $`𝕋=\{x:|x|=1\}`$. As these are the only cocycles we consider, we abbreviate this to cocycle.
For $`vN_𝒟(𝒞)`$, let $`\sigma _v`$ denote the partial homeomorphism on $`𝔊_0=\widehat{D}`$ induced by the map $`dvdv^{}`$. A measure $`\mu `$ on $`𝔊_0`$ is $`𝔊`$-quasi-invariant if, for each $`vN_𝒟(𝒞)`$, the measures $`\mu `$ and $`\mu \sigma _v`$ are equivalent, as measures on the domain of $`\sigma _v`$. Given such a measure $`\mu `$, we say a cocycle $`\alpha :𝔊𝕋`$ is $`\mu `$-measurable if, for each $`vN_𝒟(𝒞)`$, the function, denoted $`\alpha _v`$, from domain of $`\sigma _v`$ to $``$ that sends $`x`$ to $`\alpha (x,\sigma _v(x))`$ is measurable.
Since $`𝒞`$ is generated by the diagonal $`DC(𝔊_0)`$ and $`N_D(𝒞)`$, we can build a representation $`\rho `$ of $`𝒞`$ on $`L^2(𝔊_0,\mu )`$ by defining the action of $`\rho `$ on $`𝒟`$ and on $`N_𝒟(𝒞)`$ and then extending by linearity to $`𝒞`$. For $`f𝒟C(𝔊_0)`$ and $`vN_𝒟(𝒞)`$, define respectively
$$\rho (f)\eta =f\eta ,\rho (v)\eta =\alpha _v\left[\frac{d(\mu \sigma _v)}{d\mu }\right]^{1/2}(\eta \sigma _v).$$
###### Theorem 11.
\[21, Theorem II.1.1\] Every representation satisfying conditions 1 and 2 is unitarily equivalent to one arising as above from a $`𝔊`$-quasi-invariant measure $`\mu `$ and a $`\mu `$-measurable cocycle $`\alpha `$, for some choice of $`\mu `$ and $`\alpha `$.
Suppose now that $`\pi `$ is an admissible representation. Since $`\pi `$ is multiplicity free, the support of $`\mu `$ is a countable set $`S`$. The irreducibility of $`\pi `$ implies that $`S`$ is the orbit of a single point of $`𝔊_0`$ and because $`\mathrm{Lat}(\pi (𝒜))`$ has order type $``$, $`S`$ is ordered by $`\mathrm{Spec}(𝒜)`$ as $``$. Thus $`L^2(𝔊_0,\mu )`$ may be identified with with $`\mathrm{}^2()`$. Using $`\{e_j:j\}`$ for the basis vectors of $`\mathrm{}^2()`$ and letting $`P_n`$ be the projection onto $`\overline{\mathrm{span}\{e_k:k<n\}}`$, the lattice of $`\pi (𝒜)`$ is $`\mathrm{Lat}\pi (𝒜)=\{0,I\}\{P_n:n\}`$.
Given a finite subset $`Y𝔊_0`$, we associate a digraph algebra (an algebra isomorphic to a finite-dimensional CSL algebra) to $`𝔖=\mathrm{Spec}(𝒜)(Y\times Y)`$, namely the span of the rank one operators $`e_x(e_y)^{}`$ for $`(x,y)𝔖`$ acting on the space $`\mathrm{}^2(\{e_y:yY\})`$.
###### Lemma 12.
Given a finite subset $`Y𝔊_0`$, let $`S`$ be the digraph algebra associated to $`𝔖=\mathrm{Spec}(𝒜)(Y\times Y)`$. There is an isometric inclusion $`\zeta :S𝒜`$ so that $`s`$ is in the graph of $`\zeta (e_s)`$ for each $`s𝔖`$.
Lemma 12 was proved in \[7, Lemma 4.2\]; we need only observe that the inclusion constructed there is isometric.
### Proof of Theorem 10
Let $`\{P_n:n\}`$ be the projections onto the elements of $`\mathrm{Lat}(\pi (𝒜))`$, listed in decreasing order; thus $`0<\mathrm{}<P_2<P_1<P_0=I`$. For $`n\{1,2,\mathrm{}\}`$ we let $`e_n`$ be a unit vector in the range of $`P_{n+1}P_n`$; since $`\mathrm{Lat}(\pi (𝒜))`$ is multiplicity free, $`\{e_n\}_{n=\mathrm{}}^1`$ is an orthonormal basis for $``$.
We first show that the singly generated invariant manifolds are closed. Consider the manifold $`M`$ generated by $`e_1`$. Clearly, $`M`$ is dense in $``$. We shall show that if $`x`$ and $`x,e_10`$, then there exists $`T\pi (𝒜)`$ such that $`Te_1=x`$ and, moreover, that $`T`$ can be taken so that $`T^1\pi (𝒜)`$.
Fix an element $`x`$ with $`x,e_1=1`$ and choose a decreasing sequence of positive numbers $`\epsilon _k`$ such that $`_{k=1}^{\mathrm{}}\epsilon _k=\delta `$, where $`\delta <(1+x)^1`$. Since $`P_n^{}`$ are finite rank and converge strongly to $`I`$, we may choose an increasing sequence $`n_k`$ so that $`P_{n_1}x<1`$ and $`P_{n_k}xP_{n_{k+1}}x<\epsilon _k`$, for $`k>1`$.
Let $`x_1=xP_{n_1}x`$ and, for $`k>1`$, let $`x_k=P_{n_k}xP_{n_{k+1}}x`$. Clearly, $`_{k>1}x_k<\delta `$. Since there is a natural identification between $`(IP_{n_k})`$ and $`^{n_k}`$, we may regard $`x_k`$ as an element of $`^{n_k}`$.
Now let $`X_1M_{n_1}()`$ be defined by $`e_1^{}e_1^{}+x_1^{}e_1^{}I`$. Here, $`e_1`$ denotes the “last standard basis vector” in $`^{n_1}`$. Since $`x,e_1=1`$, we find that relative to the decomposition $`I=(Ie_1^{}e_1^{})+e_1^{}e_1^{}`$, $`X_1`$ has the upper triangular form
$$X_1=\left[\begin{array}{cc}I_{n_11}& v\\ 0& 1\end{array}\right].$$
Thus $`X_1^2=I_{n_1}`$. For $`k>1`$, let $`X_k=x_k^{}e_1^{}T_{n_k}()`$.
Let $`𝔒𝔊_0`$ be the support of the measure $`\mu `$. Since $`𝔒`$ has a natural identification with $``$, for each $`k`$, let $`Y_k𝔒`$ be that part identified with $`\{n_k,\mathrm{},2,1\}`$. Let $`\zeta _n:T_{n_k}()𝒜`$ be the isometric embedding associated to $`Y_k`$ given by Lemma 12. Since $`\zeta _n`$ is isometric and $`_{k>1}X_k=_{k>1}x_k<\delta `$, we see that the sum
$$\underset{k=1}{\overset{\mathrm{}}{}}\zeta _{n_k}(X_k)$$
converges uniformly to an element $`X𝒜`$. Notice also that if we let $`Z=_{k>1}\zeta _{n_k}(X_k)`$, then $`X=\zeta _{n_1}(X_1)+Z`$. Since $`\zeta _{n_1}(X_1)`$ is a square root of $`I`$, $`\zeta _{n_1}(X_1)<1+x`$, and $`Z<\delta `$, we find $`X=\zeta _{n_1}(X_1)(I+\zeta _{n_1}(X_1)Z)`$ is invertible and $`X^1𝒜`$.
Let $`T=\pi (X)`$. Then $`T`$ is an invertible element of $`\pi (𝒜)`$ and an examination of the construction shows that $`Te_1=x`$. Note that if $`x,e_1=0`$, the same construction still gives an operator $`T`$ in $`\pi (𝒜)`$ such that $`Te_1=x`$; in this case $`T`$ is no longer invertible.
We conclude that if $`y_1`$ and $`y_2`$ are vectors in $``$ with $`y_1,e_10`$, then there exists $`T\pi (𝒜)`$ such that $`Ty_1=y_2`$. (Indeed, find $`S_i\pi (𝒜)`$ such that $`S_ie_1=y_i`$ and $`S_1`$ is invertible; then take $`T=S_2S_1^1`$.)
It follows from our work so far that if $`x`$ has $`x,e_10`$, then the invariant manifold generated by $`x`$ is $``$, which is obviously closed.
Now let $`x`$ be an arbitrary unit vector and let $`M`$ be the invariant manifold generated by $`x`$. The closure of $`M`$ is an element of the nest, so $`\overline{M}=P_n`$ for some $`n`$. Clearly, $`x,e_n0`$ and by “compressing” the argument above to $`P_n`$ we see that $`M=P_n`$, so $`M`$ is closed. Hence all singly generated invariant manifolds are closed.
The result now follows from Proposition 6. ∎
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# Hopping in Disordered Media: A Model Glass Former and A Hopping Model
### 1.1 Abstract
Two models involving particles moving by “hopping” in disordered media are investigated:
A model glass-forming liquid is investigated by molecular dynamics under (pseudo-) equilibrium conditions. “Standard” results such as mean square displacements, intermediate scattering functions, etc. are reported. At low temperatures hopping is present in the system as indicated by a secondary peak in the distribution of particle displacements during a time interval $`t`$. The dynamics of the model is analyzed in terms of its potential energy landscape (potential energy as function of the $`3N`$ particle coordinates), and we present direct numerical evidence for a 30 years old picture of the dynamics at sufficiently low temperatures. Transitions between local potential energy minima in configuration space are found to involve particles moving in a cooperative string-like manner.
In the symmetric hopping model particles are moving on a lattice by doing thermally activated hopping over energy barriers connecting nearest neighbor sites. This model is analyzed in the extreme disorder limit (i.e. low temperatures) using the Velocity Auto Correlation (VAC) method. The VAC method is developed in this thesis and has the advantage over previous methods, that it can calculate a diffusive regime in finite samples using periodic boundary conditions. Numerical results using the VAC method are compared to three analytical approximations, including the Diffusion Cluster Approximation (DCA), which is found to give excellent agrement with the numerical results.
## Chapter 2 Introduction
This Ph.D. thesis deals with two models characterized by particles moving by “hopping” in disordered media. The first model is a simple model glass-former, investigated here by computer simulations using molecular dynamics . Glass forming liquids are liquids that upon cooling falls out of equilibrium and form a glass, see figure 2.1. For general reviews on glass-forming liquids see . For reviews on computer simulations of glass-forming liquids see .
The normal liquid behavior is to crystallize upon cooling below the melting temperature, $`T_L`$. In a large class of liquids the crystallization can be avoided (e.g. by fast cooling) and the liquid becomes “supercooled”. A supercooled liquid is in a pseudo-equilibrium state; The crystal is the real equilibrium state, but the supercooled state is typically stable on very long time scales. The supercooled state is characterized by relaxation times increasing strongly with decreasing temperature. If cooled with a constant cooling rate this means that the liquid at some temperature falls out of equilibrium and undergoes a glass-transition and forms a glass, which is a disordered solid. The glass-transition temperature, $`T_g`$, depends on the cooling rate; a lower cooling rate results in a lower glass transition temperature. By convention the “laboratory” glass transition is taken to be where the relaxation time $`\tau `$ is on the order of 100 sec (i.e. the “typical” time scale in the laboratory)<sup>1</sup><sup>1</sup>1The relaxation time is related to the viscosity $`\eta `$ by $`\eta =G_{\mathrm{}}\tau `$, where $`G_{\mathrm{}}`$ is the instantaneous shear modulus. $`\tau 100s`$ typically corresponds to $`\eta 10^{13}\text{poise}`$, which is sometimes used as the definition of the laboratory glass transition..
Binary systems with particles differing about 20% in diameter have been found to be good candidates for simple models of glass-forming liquids in computer simulations . At sufficiently low temperatures particles in these systems are found to move by “hopping”, i.e. they are localized for a period of time, and then move more or less directly to another position, where they again become localized . Since hopping of some form is expected to play a increasingly dominant role as temperature is lowered , systems exhibiting hopping at temperatures which can be reached under equilibrium conditions in computer simulations are of special interest. In this thesis we investigate the dynamics of a binary Lennard-Jones liquid, which has earlier been shown to exhibit hopping . This is done under (pseudo-) equilibrium conditions, i.e. above the glass transition temperature.
The model glass-former investigated in this thesis is *not* a model of any particular liquid existing in the laboratory (or the real world for that matter). Like most models glass formers studied in computer simulations, it should be thought of as a (very) simple model exhibiting some of the complex behavior found in real glass forming liquids. The reason why it is interesting to investigate such a simple model is two-fold: I) The glass-transition is found to occur in liquids that are chemically very different, but the related phenomenology is found to be strikingly similar. Since liquids with different chemical details behave in a similar way, it makes sense to study simple models ignoring the chemical details, in an attempt to understand the fundamental mechanisms. II) Most theories for glass-forming liquids also treat these in a highly simplified manner, ignoring the chemical details (for the same reasons given above). Consequently computer simulations makes it possible to test these theories on their own terms, i.e. without any additional approximations. One particular example of this is the Mode Coupling Theory (MCT) . When doing computer simulations of glass forming liquids it is rather natural to compare the results with the predictions of the MCT since these are plenty and detailed, *and* MCT deals with the temperature range accessible (at the present) to computer simulations under equilibrium conditions. It is not a goal in it self to test MCT in the present thesis, but some of the results are compared to the predictions of the MCT.
The second model investigated in the present thesis is a so-called “hopping-model”, where the hopping behavior of particles is “built in”; In the symmetric hopping model particles are moving on a lattice by doing thermally activated hopping over random energy barriers. This model is found to have interesting universal features in the extreme disorder limit, i.e. when the temperature goes to zero. The symmetric hopping model has earlier been treated as a model for frequency dependent conduction in glasses , with particles representing non-interacting charge carriers (ions or electrons). In that context the universality manifests itself as follows: At low temperatures the shape of the conduction vs. frequency curve becomes independent of temperature, which is also what is found in experiments (see eg. ). In the present thesis the symmetric hopping model is treated in the slightly more general context of diffusion in disordered media.
### 2.1 Outline
This thesis consists of 3 main chapters:
In chapter 3 we report results from the simulations of the model glass-former mentioned above, with emphasis on the dynamics. In this chapter “standard” results are reported, i.e. the measures commonly used to describe (glass-forming) liquids. All of the features found for this particular model has been reported earlier for other models, and this chapter thus constitutes a “review” of the behavior of model glass-formers, illustrating the similarities in the behavior of these.
In chapter 4 the dynamics of the model glass-former described in chapter 3 is investigated in term of its “potential energy landscape”, which is simply the potential energy as a function of the $`3N`$ particle coordinates. The new concept of “inherent dynamics” is introduced, and the results from applying this concept to the model described in chapter 3 are discussed.
Chapter 5 deals with the symmetric hopping model. A new numerical method for investigating the model in the extreme disorder limit is developed. Numerical results are presented and compared with 3 analytical approximations.
The three main chapter (3-5) contains individual conclusions.
### 2.2 Acknowledgments
The present Ph.D. thesis is the result of work done at Department of Mathematics and Physics (IMFUFA) at Roskilde University (RUC). I would like to thank the people at IMFUFA, especialy Jeppe C. Dyre (my supervisor), Niels B. Olsen, Tage E. Christensen, Johannes K. Nielsen and Ib Høst Pedersen.
The spring of 1998 I spent as a guest researcher at Center for Theoretical and Computational Materials Science (CTCMS) at National Institute of Standards and Technology (NIST), Maryland, USA. I would like to thank NIST and CTCMS for the hospitality, and I would like to thank the people I met and worked with there: Sharon C. Glotzer, Srikanth Sastry, Jack F. Douglas, Claudio Donati, and Paulo Allegrini.
### 2.3 Papers
The following papers contain results obtained for the two models discussed in the present thesis:
* Paper I: *Effective one-dimensionality of universal ac hopping conduction in the extreme disorder limit.* J.C. Dyre and T.B. Schrøder. Phys. Rev. B. 54 14884 (1996).
* Paper II: *Hopping in a supercooled binary Lennard-Jones liquid.* T.B. Schrøder and J.C. Dyre. J. Non-Cryst. Solids. 235-237 331 (1998).
* Paper III: *Crossover to potential energy landscape dominated dynamics in a model glass-forming liquid* . T.B. Schrøder, S. Sastry, J.C. Dyre and S.C. Glotzer. J. Chem. Phys. 112 (2000) in press.
(http:$`\backslash \backslash `$xxx.lanl.gov: cond-mat/9901271)
* Paper IV: *Scaling and Universality of ac Conduction in Disordered Solids.* T.B. Schrøder and J.C. Dyre. Phys. Rev. Lett. 84 310 (2000).
* Paper V: *Universality of ac conduction in disordered solids.* J.C. Dyre and T.B. Schrøder. Rev. Mod. Phys. 72 (2000) in press.
Paper I reports numerical results for the symmetric hoping model, which was done mainly during my master thesis . Paper II-V report results obtained as part of the present Ph.D. thesis. Paper IV and V was written after the conclusion of my Ph.D. (summer 1999).
## Chapter 3 A Model Glass-former
A glass-forming binary Lennard-Jones liquid is investigated by molecular dynamics under (pseudo-) equilibrium conditions. In this chapter we report ’standard’ results for this model liquid, i.e. pair-correlation functions, mean square displacements, intermediate scattering functions, etc. The main result is, that at low temperatures hopping is present in the system as indicated by a secondary peak in $`4\pi r^2G_s(r,t)`$, where $`G_s(r,t)`$ is the van Hove self correlation function. It has not been possible to identify a temperature range, where the asymptotic predictions of the ideal mode coupling theory are fulfilled. Some of the results reported in this chapter are contained in paper II and paper III.
### 3.1 Model and Method
The system investigated in the present work, has earlier been shown to exhibit hopping by Wahnström . By “hopping” is here meant, that particles behave like illustrated in figure 3.1; The particle shown stays relatively localized for a period of time, and then move some distance, where it again becomes localized. The presence of hopping was the main motivation for investigating this system, and this kind of dynamics will be discussed in more detail in this and the following chapter.
In the work of Wahnström, the system was investigated both above and below the glass-temperature, $`T_g`$, i.e. the system was allowed to fall out off equilibrium as it was cooled. As a consequence, it was unclear whether the hopping seen was a feature of the equilibrium liquid, or if it was a consequence of non-equilibrium dynamics. Here we attempt to keep the system equilibrated at all temperatures, i.e. the equilibration time is made longer and longer as the temperature is lowered.
The “Wahnström system” is a binary mixture of $`N=500`$ particles, with 50% particles of type A, and 50% particles of type B<sup>1</sup><sup>1</sup>1The exact number of particles used in the simulations are: $`N_A=251`$, and $`N_B=249`$. The particles interact via the pair-wise Lennard-Jones potential, where the parameters depend on the types of the two particles involved ($`\alpha `$ and $`\beta `$):
$$V_{\alpha \beta }(r)=4ϵ_{\alpha \beta }\left(\left(\frac{\sigma _{\alpha \beta }}{r}\right)^{12}\left(\frac{\sigma _{\alpha \beta }}{r}\right)^6\right)$$
(3.1)
The forces are given by the negative gradient of the potential:
$`𝐅_{\alpha \beta }(𝐫)`$ $`=`$ $`V_{\alpha \beta }={\displaystyle \frac{V_{\alpha \beta }}{r}}{\displaystyle \frac{𝐫}{r}}`$ (3.2)
$`=`$ $`{\displaystyle \frac{48ϵ_{\alpha \beta }}{(\sigma _{\alpha \beta })^2}}\left(\left({\displaystyle \frac{\sigma _{\alpha \beta }}{r}}\right)^{14}{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{\sigma _{\alpha \beta }}{r}}\right)^8\right)𝐫`$ (3.3)
$`V_{\alpha \beta }(r)`$ is characterized by a minimum at $`V_{\alpha \beta }(2^{1/6}\sigma _{\alpha \beta })=ϵ_{\alpha \beta }`$, steep repulsion at shorter distances, and a weaker attraction at longer distances.
The length-parameters used in the Wahnström system are<sup>2</sup><sup>2</sup>2We follow and term the large particles “A”, and small particles “B”: $`\sigma _{AA}=1`$, $`\sigma _{BB}=1/1.20.833`$ and $`\sigma _{AB}=(\sigma _{AA}+\sigma _{BB})/2`$. The energy-parameters are all identical: $`ϵ_{AA}=ϵ_{AB}=ϵ_{BB}=1`$. The masses of the particles are given by $`m_A=2`$, and $`m_B=1`$. The length of the sample was $`L=7.28`$, which gives a (reduced) density of $`\rho =N/L^3=1.296`$. Times are reported in units of $`\tau (m_B\sigma _{AA}^2/48ϵ_{AA})^{1/2}`$ (This expression contains an error in paper II). When comparing simulations of a model-liquid like the one described here with experimental results for real liquids, it is customary to use “Argon units”, i.e. parameters used when modeling Argon atoms by the Lennard-Jones potential; $`\sigma =3.4\AA `$, $`ϵ=120Kk_B`$, and $`\tau =3\times 10^{13}sec`$ .
The potential is “cut and shifted” at $`r=2.5\sigma _{\alpha \beta }`$, which means that it is set to zero for $`r2.5\sigma _{\alpha \beta }`$, and $`|V_{\alpha \beta }(2.5\sigma _{\alpha \beta })|0.016`$ is added to the potential for $`r2.5\sigma _{\alpha \beta }`$. This makes the potential continuous at $`r=2.5\sigma _{\alpha \beta }`$, while the force is not. The equations of motion are integrated with periodic boundary conditions using the Leap-Frog algorithm (which is simply a discretizaition of Newton’s second law) with a time step of $`0.01\tau `$.
Three independent samples were used, each initiated by generating a random configuration followed by equilibration at a high temperature ($`T=5.0`$). The cooling was done by controlling the total energy of the system, which was done by scaling the velocities. An equilibration run of length $`t_{eq}`$ was then performed, followed by a production run of the same length. If the samples was determined not to be equilibrated properly, $`t_{eq}`$ was doubled, by “degrading” the production run to be part of the equilibration run, and making a new production run twice as long. This procedure was continued until the samples was determined to be equilibrated. The criteria used for determining if the samples were equilibrated will be discussed later. Note, that by this procedure, it is the total energy ($`E_{tot}=E_{pot}+E_{kin}`$) that is controlled. The reported temperatures and pressures are computed as time-averages over the instantaneous temperature and pressure respectively :
$`T(t)`$ $``$ $`{\displaystyle \frac{2E_{kin}(t)}{3N3}}`$ (3.4)
$`P(t)`$ $``$ $`\rho T(t)+W(t)/V,W(t){\displaystyle \frac{1}{3}}{\displaystyle \underset{i}{}}{\displaystyle \underset{j>n}{}}𝐫_{ij}𝐅_{ij}`$ (3.5)
where $`W(t)`$ is the virial and the summing is over all pair of particles.
Some of the results for the system described above, will be compared with results from a different binary Lennard-Jones mixture (the “Kob & Andersen system”), which has been investigated in a number of papers . This system is a 80:20 mixture, with $`\sigma _{AA}=1`$, $`\sigma _{AB}=0.8`$, $`\sigma _{BB}=0.88`$, $`\sigma _{AB}=0.8`$, $`ϵ_{AA}=1.0`$, $`ϵ_{AB}=1.5`$, $`ϵ_{BB}=0.5`$, and $`m_A=m_B=1`$.
### 3.2 Static Results
Figure 3.2 shows as a function of temperature a) the potential and total energy per particle, and b) the pressure vs. Temperature. The error-bars are estimated from deviations between the three independent samples, which are found to be in reasonable agreement. This is a necessary (but not sufficient) condition for the system(s) to be equilibrated. Note that its the total energy which was set to be equal for the three samples, and consequently there are small deviations in the measured temperatures.
The specific heat capacity at constant volume, $`C_v`$, is given by :
$$C_v=\frac{1}{N}\frac{dE_{tot}}{dT}$$
(3.6)
From figure 3.2a we see that $`C_v`$ increases as the temperature is lowered. Note that there is no indication of a glass-transition, since this would be indicated by a sharp decrease in $`C_v`$ when cooling below the glass transition temperature, $`T_g`$, as seen in the results by Wahnströhm. The heat capacity at constant volume will be discussed further in section 4.2.
In figure 3.3 is shown the pair-correlation functions, $`g_{AA}(r)`$, $`g_{AB}(r)`$, and $`g_{BB}(r)`$, at three of the temperatures simulated. The pair-correlation function, $`g_{\alpha \beta }(r)`$, is the average relative density of particles of type $`\beta `$ around particles of type $`\alpha `$ . It is normalized to be 1 at large distances, $`r`$, were the correlations disappears. At the highest temperature ($`T=1.06`$) all three pair-correlation functions is seen to look like ’typical’ high temperature pair-correlation functions, with a sharp first neighbor peak, a more rounded second neighbor peak, etc. As the temperature is lowered the first neighbor peak becomes sharper and the second peak splits into two peaks. The splitting of the second peak upon cooling is often seen in super-cooled liquids (see eg.).
The parameters used in the potential does not energeticly favor the mixture of A and B particles ($`ϵ_{AB}`$ is not larger than $`ϵ_{AA}`$ and $`ϵ_{BB}`$), which might cause the system to phase separate at sufficiently low temperatures. In the event of phase separation, we would expect the area of the first peak in $`g_{AB}(r)`$ to decrease at the temperature where phase separation starts occurring. There is no indication of this in figure 3.3.
In figure 3.4 is shown the static structure factors, $`S_{AA}(q)`$, $`S_{AB}(q)`$, and $`S_{BB}(q)`$, corresponding to the data shown in figure 3.3. The static structure factor is the (3 dimensional) Fourier transform of pair correlation function, and thus contains the same information. However the finite sample size introduces some features in $`S_{\alpha \beta }(q)`$ which are clearly unphysical (see eg. ). By recalculating $`S_{\alpha \beta }(q)`$ from $`g_{\alpha \beta }(r)`$ with a cut-off in $`r`$ that is less than $`L/2`$ one can get an idea about which features of $`S_{\alpha \beta }(q)`$ are due to finite-size effects. Doing this shows, that at $`q`$-values to the left of the first peak ($`q6`$, depending on type of correlation) $`S_{\alpha \beta }(q)`$ is dominated by finite-size effects, which is why most of it is not shown. Also the small “wiggles” seen at $`q10`$ is a finite-size effect. The splitting of the peaks at $`q15`$ is however not a finite-size effect.
### 3.3 Mean Square Displacement
In figure 3.5 we show the mean square displacement of the a) A particles, $`r^2(t)_A`$, and b) B particles, $`r^2(t)_B`$. At short times ($`t<1`$) a ballistic regime ($`r^2(t)_\alpha t^2`$) is seen for both A and B particles, at all temperatures. The ballistic regime is simply a consequence of the velocities being constant at these short time scales:
$$r^2(t)_\alpha =(Vt)^2_\alpha =V^2_\alpha t^2$$
(3.7)
In the insets in figure 3.5 is shown $`V^2_\alpha `$ calculated by eq. 3.7. $`V^2_\alpha `$ can also be calculated directly from the temperature:
$$T=\frac{2E_{kin}}{3(N1)}=\frac{Nm_\alpha (V/\tau )^2_\alpha }{3(N1)}=\frac{48m_\alpha V^2_\alpha }{3(11/N)}$$
(3.8)
In the insets in figure 3.5 is also shown $`V^2_\alpha `$ as calculated from eq. 3.8. The excellent agreement between eq. 3.7 and 3.8 (with no fitting involved) is not a big surprise, but it gives confidence in the argument leading to eq. 3.7, and acts as a consistency check.
In figure 3.5 a diffusive regime ($`r^2(t)_\alpha t`$) is seen at long times, for all temperatures. As the temperature is lowered the time scale at which the diffusive regime sets in increases, and a plateau is seen to evolve between the ballistic and diffusive regimes. This behavior is typical for what is seen in simulations of super-cooled liquids, and the plateau is argued to be associated with particles being trapped in local “cages” consisting of their neighbors .
The vertical dashed lines in figure 3.5 identifies the time $`t_1`$, defined by $`r^2(t_1)_\alpha =1`$. The significance of this time will be discussed in connection with figure 3.6.
Fitting $`r^2(t)_\alpha `$ in the diffusive regime to the Einstein relation:
$$r^2(t)_\alpha =6Dt$$
(3.9)
we determine the diffusion coefficient, $`D`$, as a function of temperature (and particle type). We postpone the discussion of the temperature dependence of $`D`$ to section 3.6, where it will be discussed together with other measures of the time scales involved (eg. $`t_1`$). Following Kob and Andersen we in figure 3.6 present the data for $`r^2(t)_\alpha `$ as a function of $`6Dt`$. In the diffusive regime the data for all temperatures should, according to eq. 3.9, fall on the line $`6Dt`$, which is seen to be the case for $`r^2(t)_\alpha 1`$. In other words the time $`t_1`$, defined by $`r^2(t_1)_\alpha =1`$, can be viewed as marking the onset of the diffusive regime. The fact that the diffusive regime is reached at all temperatures (compare fig. 3.5), is a necessary (but not sufficient) condition, for the system(s) to be in equilibrium.
As is the case for the data presented by Kob and Andersen, the data in figure 3.6 seems to indicate a universal behavior; As the temperature is decreased $`r^2(t)_\alpha `$ follows the thick dashed curve in figure 3.6 to lower and lower scaled times. The thick dashed curves are fits to the fitting function used by Kob and Andersen (see table 3.1):
$$r^2(t)_\alpha =r_c^2+A(Dt)^b+6Dt$$
(3.10)
The parameter $`b`$ is the so-called “von-Sweidler” exponent, which is related to the dynamics at intermediate times, i.e. between the plateau and the diffusive regime. According to the asymptotic predictions of the ideal mode coupling theory (MCT) , the von-Sweidler exponent, $`b`$, should be independent of the particle type. Kob and Andersen argue that the two values they find, 0.48 for the A particles and 0.43 for the B particles, are within reasonable agreement. This is *not* the case for the $`b`$-values found here; $`0.17\pm 0.03`$ for the A particles and $`0.27\pm 0.01`$ for the B particles.
### 3.4 The van Hove Correlation Function
To characterize the dynamics in more detail, we calculate the van Hove self correlation function (where the sum is over particles of type $`\alpha `$):
$$G_{s\alpha }(𝐫,t)=\frac{1}{N_\alpha }\underset{i=1}{\overset{N_\alpha }{}}\delta (𝐫_i(t)𝐫_i(0)𝐫),$$
(3.11)
which is the probability density of a particle of type $`\alpha \{A,B\}`$ being displaced by the vector $`𝐫`$, during the time interval $`t`$. In an isotropic system $`G_{s\alpha }(𝐫,t)`$ does not depend on the direction of $`𝒓`$, and the probability distribution of a particles being displaced the *distance* $`r`$ is $`4\pi r^2G_{s\alpha }(r,t)`$. The mean square displacement, $`r^2(t)_\alpha `$, is given by the second moment of $`4\pi r^2G_{s\alpha }(r,t)`$, and the later thus gives a more detailed view of the dynamics.
In figure 3.7 is plotted $`4\pi r^2G_{sA}(r,t_1)`$ and $`4\pi r^2G_{sB}(r,t_1)`$, i.e. the distribution of distances moved by A and B particles respectively in the time interval $`t_1`$, which in the previous section was demonstrated to be at the onset of the diffusive regime. The thick curve is the Gaussian approximation:
$$G_s(r,t)=\left(\frac{3}{2\pi r^2(t)}\right)^{3/2}\mathrm{exp}\left(\frac{3r^2}{2r^2(t)}\right)$$
which is seen to be reasonably fulfilled at high temperatures. As the temperature is lowered the results starts graduately deviating from the Gaussian approximation, and a shoulder builds at the average inter particle distance, $`r1.0`$, which at T=0.59 becomes a well-defined second peak. The second peak, observed also in other model liquids, is interpreted as single particle hopping (see figure 3.1), i.e. the particles stay localized in their local cages a certain time (first peak), then moves more or less directly out to distances approximately equal to the inter particle distance (second peak) (This behavior was also reported in ). To what extent the hopping behavior is related to activated hopping over energy barriers will be discussed in chapter 4.
Note that when analyzing the dynamics in more detail by means of the van Hove self correlation function (fig. 3.7), the universality seen in the mean square displacement (fig. 3.6) is completely gone; The way the system “achieves” the diffusive behavior (i.e. the dynamics at $`t_1`$) changes qualitatively from high temperature (Gaussian) to low temperature (Hopping).
The deviation from Gaussian behavior is often analyzed in terms of the Non-Gaussian parameter :
$$\alpha _2(t)=\frac{r^4(t)r^4(t)_{Gauss}}{r^4(t)_{Gauss}}=\frac{3r^4(t)}{5r^2(t)^2}1$$
(3.12)
which is the relative deviation between the measured $`r^4(t)`$ and what it would be for a given value of $`r^2(t)`$, if $`G_s(r,t)`$ was Gaussian: $`r^4(t)_{Gauss}=5r^2(t)^2/3`$. In figure 3.8 is shown $`\alpha _2`$ for the A and B particles respectively.
Like in we see a universal behavior in the way $`\alpha _{2A}(t)`$ increases, before it reaches its maximum. Following a power-law fit is done in this regime, from which is found exponents of $`0.64`$ and $`0.67`$ for the A and B particles respectively. The exponent found in is $`0.4`$. As noted in the universality seen in $`\alpha _2(t)`$ is different from the one found in the behavior of $`r^2(t)_\alpha `$ (fig. 3.6); The universality seen in $`\alpha _2(t)`$ does not involve any scaling.
The time $`t^{}`$ which is defined as the position of the maximum in $`\alpha _2(t)`$, can be interpreted as the time where the dynamics deviates most from Gaussian behavior . In figure 3.9 is shown for three temperatures, the time development of $`4\pi r^2G_{sB}(r,t)`$. At each temperature $`4\pi r^2G_{sB}(r,t)`$ is shown at a time close to<sup>3</sup><sup>3</sup>3The reason that $`t^{}`$ itself is not used here, has to do with at what time-steps configurations are stored, and thus which time differences are easy accessible $`t^{}`$ (indicated by the arrows), and 2,4,8, and 16 times that time. At $`T=1.06`$ the distribution of particle displacements is seen to be characterized by a single peak, which “spreads out” without any qualitative change in the shape, as time increases. This is the typical behavior for liquids at high temperatures . At $`T=0.66`$ the behavior is seen to be qualitatively similar, except for a weak indication of a shoulder at $`r1.0`$. The time development of $`4\pi r^2G_{sB}(r,t)`$ at $`T=0.59`$ is qualitatively different from what is found at higher temperatures; The first peak decreases while the position is almost constant, and at the same time as the second peak starts building up. This supports the “hopping” interpretation of the second peak given above; The particles escape their local “cages” (the first peak), by “suddenly appearing” at approximately the inter particle distance (second peak). After the last time shown here, the second peak starts decreasing. This is interpreted as a consequence of particles hopping again.
With the regards to the interpretation of $`t^{}`$ mentioned above, one should note that for $`T=0.59`$ the second peak in $`4\pi r^2G_{sB}(r,t)`$ builds up *after* the time $`t^{}`$, i.e. after $`\alpha _2(t)`$ has its maximum. A more detailed analysis of the non-Gaussian behavior should include the higher order analogues of $`\alpha _2(t)`$ (involving higher moments of $`4\pi r^2G_{s\alpha }(r,t)`$) , but this approach will not be pursued further here.
The van Hove distinctive correlation function, $`G_{D\alpha \beta }(r,t)`$, is the time dependent generalization of the pair correlation function :
$$G_{D\alpha \beta }(𝐫,t)=\frac{1}{\sqrt{N_\alpha N_\beta }}\underset{i=1}{\overset{N_\alpha }{}}\underset{j=1}{\overset{N_\beta }{}}\delta (𝐫_j(t)𝐫_i(0)𝐫),$$
(3.13)
where the first sum is over particles of type $`\alpha `$, the second sum is over particles of type $`\beta `$, and $`ij`$ if $`\alpha =\beta `$ (this last term is contained in $`G_{s\alpha }(r,t)`$, see equation 3.11). $`G_{D\alpha \beta }(r,t)`$ is the relative density at time $`t`$ of $`\beta `$-particles at $`𝐫`$, given that there at time $`t=0`$ was a $`\alpha `$-particle at $`𝐫=0`$ (normalised to be 1 if there is no correlation).
In figure 3.10 is shown $`G_{DBB}(r,t)`$ for the same temperatures and times as in figure 3.9. Included in figure 3.10 is also $`G_{DBB}(r,t=0)=g_{BB}(r)`$. Except for the lowest values of $`r`$, the features of the pair correlation function is seen to approach the long-time value $`G_{D\alpha \beta }(r,t)=1`$ in a smooth manner, which is the typical high-temperature behavior . As the temperature is decreased, a significant “overshoot” develops at the small $`r`$-values<sup>4</sup><sup>4</sup>4the data for $`r<0.1`$ is removed, since the statistical noise dominates the data in this range, indicating that there is an excess probability that when a particle jumps, it does so to a position previously occupied by another particle. Thus, while $`G_{sB}(r,t)`$ (fig. 3.9a) shows that particles has a tendency to jump a distance approximately equal to the inter-particle distance, $`G_{D\alpha \beta }(r,t)`$ give the further information, that they have a tendency to jump to a position previously occupied by another particle. Similar (but less pronounced behavior) was found in .
In a scenario where the dynamics is dominated by particles jumping to positions previously occupied by other particles, one should expect that moving particles make up correlated string-like objects in the liquid. This kind of behavior is found by Muranaki and Hiwatari in a soft sphere system , and by Donati et. al. in the Kob & Andersen system . This type of behavior will be discussed further in chapter 4.
### 3.5 The Intermediate Scattering Function
The self part of the intermediate scattering function, $`F_s(q,t)`$, is the (3 dimensional) Fourier transform of the van Hove self correlation function :
$`F_s(𝐪,t)`$ $``$ $`{\displaystyle G_s(𝐫,t)e^{i𝐪𝐫}𝑑𝐫}=\mathrm{cos}𝐪(𝐫_j(t)𝐫_j(0))`$ (3.14)
$`F_s(q,t)`$ is a relaxation function, normalized to be 1 at $`t=0`$, and it goes to zero for $`t\mathrm{}`$. Figure 3.11 shows the self part of the intermediate scattering function for the A particles and B particles; $`F_{sA}(q_{maxA}=7.5,t)`$ and $`F_{sB}(q_{maxB}=8.1,t)`$. $`q_{maxA}=7.5`$ and $`q_{maxA}=8.1`$ are here the positions of the first peak in the static structure factor for the A-A correlation ($`S_{AA}(q)`$), and the the B-B correlation ($`S_{BB}(q)`$), respectively (see figure 3.4). As the temperature is lowered a two-step relaxation develops, which is what is typically seen<sup>5</sup><sup>5</sup>5 In most experiments, and some simulations, this is seen as two peaks (one for each relaxation step) in the imaginary part of the generalized susceptibility: $`\chi ^{\prime \prime }(𝐪,\omega )\omega \pi F(𝐪,\omega )`$, where $`F(𝐪,\omega )`$ is the Fourier transform of $`F(𝐪,t)`$, see eg. . in glass-forming liquids, both in experiments and in simulations . The initial relaxation is a consequence of the particles oscillating in their cages, while the the long-time (alpha) relaxation is a consequence of particles escaping their cages . The alpha relaxation is often found to be well approximated by stretched exponentials, $`f(t)=f_c\mathrm{exp}((t/\tau _\alpha )^\beta )`$ , which is also the case in figure 3.11, where fits to stretched exponentials are shown as dashed lines.
In Fig. 3.12 we show as a function of temperature the three fitting parameters used in figure 3.11; a) relaxation times, $`\tau _\alpha `$, b) stretching parameters, $`\beta `$, and c) non-ergodicity parameters $`f_c`$. As expected the relaxation times, $`\tau _{\alpha A}`$ and $`\tau _{\alpha B}`$ is seen to increase dramatically as the temperature is lowered. The asymptotic prediction of the ideal mode-coupling theory (MCT) is $`\tau _\alpha =\tau _0(TT_c)^\gamma `$ . The divergence of $`\tau _\alpha `$ at $`T_c`$ predicted by the ideal MCT is never seen in practice. This is argued to be consequence of relaxation by hopping taking over close to $`T_c`$ . This means that close to $`T_c`$ the power-law is expected to break down and it should not be fitted in this regime. Unfortunately we have no independent estimate of $`T_c`$, and furthermore it is not known how close to $`T_c`$ the power-law is expected to hold. Since we have already demonstrated, that hopping is present in the system at the lowest temperatures, we can expect that these are close to $`T_c`$, which means that we might run into the problem described above. We note also, that the power-law is an *asymptotic* prediction, i.e. it is expected to break down at high temperatures, but again it is not known how far from $`T_c`$ this is supposed to happen.
As discussed above, the temperature range (if any) where the asymptotic predictions of ideal MCT is supposed to hold, is not known. Consequently the following approach for fitting $`\tau _\alpha =\tau _0(TT_c)^\gamma `$ has been applied; First the fit was done for all 8 temperatures ($`T_{cut}=0`$), then by excluding the lowest temperature ($`T_{cut}=0.60`$), and then by excluding the two lowest temperatures ($`T_{cut}=0.63`$).
The results of the fitting-procedure described above, is seen in table 3.2. For both the A and B particles the parameters found for $`T_{cut}=0.60`$ and $`T_{cut}=0.63`$ agree within the error-bars, while they disagree with the parameters for $`T_{cut}=0`$. Thus one might argue that the parameters found for $`T_{cut}=0.60`$ and $`T_{cut}=0.63`$ describes the “true” power-law, while the fit found using $`T_{cut}=0`$ deviates as a consequence of fitting too close to $`T_c`$. Thus our best attempt at a power-law fit is the one achieved for $`T_{cut}=0.60`$ (i.e. by excluding the lowest temperature), which is the fit shown in figure 3.12 for the A and B particles respectively. The values estimated in this manner for $`T_c`$ and $`\gamma `$ for the A and B particles are identical within the error-bars, as predicted by the ideal MCT. The temperature dependence of $`\tau _\alpha `$ will be discussed further in section 3.6.
The numerical data for $`\tau _\alpha (T)`$ does not by itself give strong evidence for a dynamical transition at the estimated $`T_c=0.592\pm 0.004`$; nothing special seems to happen at that temperature. However, the agreement with the lowest temperature ($`T=0.591\pm 0.002`$), where we clearly see hopping (figure 3.7), gives us confidence that there *is* a dynamical transition close to the estimated $`T_c`$.
The stretching parameters, $`\beta _A`$ and $`\beta _B`$, are seen in figure 3.12b to decrease from values close to 1 (i.e. exponential relaxation) at high temperatures, to values close 0.5 at the lowest temperatures. The values found at high-temperatures are somewhat dependent on the time-interval used for the fitting, and the real error-bars are thus bigger than the ones shown (which are estimated from deviations between the three independent samples, with the fitting done in the same time-intervals). However focusing on the medium to low temperatures ($`T0.73`$), there can be no doubt that $`\beta `$ is increasing as function of temperature. This is in contradiction with the asymptotic predictions of the ideal MCT, which predicts $`\beta `$ to be constant. The ideal MCT does not directly predict the long-time relaxation to be stretched exponentials, but it predicts it to exhibit time-temperature super-position (TTSP) . This mean that the *shape* of the long-time relaxation should be independent of temperature, which for stretched exponentials means that $`\beta `$ should be constant.
The question of whether $`\beta `$ is constant or increases towards 1 has led to controversy about what is the right fitting procedure . Another way of checking for TTSP, is by scaling the data appropriately and look for approach to an universal curve. This is done in figure 3.13 where $`F_{sA}(q_{max},t)/f_{cA}`$ is plotted versus $`t/\tau _{\alpha A}`$ which should be identical for any temperature range where TTSP holds. No such temperature range can be identified. One might argue, that the scaling approach used in figure 3.13 relies on the values of $`f_c`$ estimated from the fits to stretched exponentials, which at the high temperatures is dependent on the time-range chosen for the fitting. Scaling $`F_{sA}(q_{max},t)`$ to agree at $`F_{sA}(q_{max},t)=e^1`$ (i.e. without dividing by $`f_c`$), like done in does *not* change the conclusion above; there is no indication of TTSP.
The fact that $`F_{s\alpha }(q,t)`$ decays to zero (figure 3.11), and that the parameters describing the alpha relaxation is in reasonably agreement for the three independent samples (as illustrated by the error-bars in figure 3.12), are necessary conditions for the liquid to be in equilibrium. Of all tests for equilibrium applied, this was the one that was found to be most sensitive, i.e. requiring the longest equilibration times.
In figure 3.14 is plotted $`F_{sA}(q,t)`$ at $`T=0.59`$ for $`q`$-values $`1,2,\mathrm{},20`$, and $`q=q_{max}=7.5`$ (thick dashed line). Also shown in figure 3.14 are fits to stretched exponentials (dashed lines). The fitting parameters used in 3.14 is plotted in figure 3.15. At the lowest $`q`$-values the fits are not perfect, but except for that the fits are reasonable. Note that in the time-range were the second peak builds up in $`4\pi r^2G_{sA}(r,t)`$, i.e. from $`t^{}610^2`$ to $`t_1710^3`$, figure 3.14 shows no clear indication of this. Since $`F_{s\alpha }(q,t)`$ is the Fourier transform of $`G_{s\alpha }(r,t)`$, and thus contains the same information, it *is* possible to extract the information about the hopping from $`F_{s\alpha }(q,t)`$. However since we already have an excellent indication of the hopping in $`4\pi r^2G_{s\alpha }(r,t)`$, and $`G_{D\alpha \beta }(r,t)`$, this will not be pursued further here.
### 3.6 Time Scales
In this section, we compare the measures of time-scales found in the previous sections. In figure 3.16 is plotted $`(6D)^1`$ (squares, from figure 3.5), $`t_1`$ (plus, from figure 3.5), $`t_\alpha `$ (diamonds, from figure 3.12), and $`t^{}`$ (circles, from figure 3.8). If $`t_1`$ is in the diffusive regime, the relation $`(6D)^1=t_1`$ should hold, see equation 3.9. This is seen in figure 3.16 to hold to a good approximation; the largest deviation is 8%, with $`t_1`$ being smaller than $`(6D)^1`$, thus confirming that $`t_1`$ is a good estimate of the onset of the diffusive regime, as argued in section 3.3. $`t^{}`$ is found to be roughly a factor $`10`$ smaller than $`t_1`$, and $`(6D)^1`$ (the ratio changes from roughly 6 at high temperatures, to roughly 11 at low temperatures). Also shown in figure 3.16 are attempts to fit $`(6D)^1`$ to a power-law. According to the asymptotic predictions of ideal MCT, $`(6D)^1`$ should have the same temperature dependence as $`\tau _\alpha `$ (apart from a constant factor), which means that we should find the same $`T_c`$ and $`\gamma `$. Consequently the fitting was done in the temperature range which gave the “best” power-law fit to $`\tau _\alpha `$, i.e. by excluding the lowest temperature. This fit is shown, for the A and B particles respectively, as the full curves in figure 3.16, and the parameters and error-bars are given in table 3.3. The fit to the data are reasonable, but the estimated $`T_c`$ deviates from the one found from $`\tau _\alpha `$, which is in contradiction with MCT. The dashed curves in figure 3.16 are the results of power-law fits where $`T_c`$ was set to have the value found for $`\tau _\alpha `$, $`T_c=0.592`$. This also fits the data reasonable, but now the exponents $`\gamma `$ are different from the ones found from $`\tau _\alpha `$, thus illustrating that $`\tau _\alpha `$ and $`(6D)^1`$ has different temperature dependence (as can be seen directly in figure 3.16). Also Kob and Andersen find different temperature dependence for the diffusion and the relaxation .
Figure 3.17 shows the alpha relaxation times $`\tau _{\alpha A}`$ and $`\tau _{\alpha B}`$, in the so-called “Arrhenius-plot”, i.e. as a function of $`1/T`$ and with a logarithmic y-axis. In this plot, a Arrhenius dependence of the relaxation time, $`\tau _\alpha =\tau _0\mathrm{exp}(E/T)`$ would give a straight line, which is clearly not the case. Or in the terminology of Angel , the liquid is “fragile” (non-Arrhenius), as opposed to “strong” (Arrhenius). Included as full lines in figure 3.17 are fits to the Vogel-Fulcher law, $`\tau _\alpha =\tau _0\mathrm{exp}(A/(TT_0))`$, which is often seen to fit the relaxation times over a range of temperatures in experiments. The fits are extrapolated, to include the time $`\tau _\alpha 310^{15}`$ which in Argon units (see section 3.1) corresponds to $`\tau _\alpha 10^3`$ seconds, i.e. including the time scales involved in the laboratory glass-transition (see introduction). This extrapolation over 12 decades in time (from information covering 2.5 decades), should of course not be taken too seriously, but is included here to illustrate the large differences in time scales. Included as dashed lines in figure 3.17 are the power-law fits from figure 3.12.
### 3.7 Finite Size Effects
To estimate finite size-effects we in this section present results from a sample with $`N=1000`$ particles which is otherwise equivalent with the samples described in the previous sections (i.e. same density and total energy per particle).
Figure 3.18 shows the self part of the intermediate scattering function for the A particles for the N=1000 system, together with fits to stretched exponential (dashed lines), using the same fitting procedure as for the N=500 samples (figure 3.11). The resulting fitting parameters are shown in figure 3.19 together with those found for N=500 (figure 3.12). The results for the two system sizes are seen to be in reasonable agreement, indicating that the results in the previous sections do *not* depend strongly on system size. In a decrease in $`\tau _\alpha `$ by a factor of $`30`$ was found at low temperatures for a soft sphere system by going from $`N=108`$ to $`N=10000`$. At the moment it is unclear if this much larger change in relaxation time is due to the larger difference in system size, starting from a smaller system, or other differences between the two systems.
### 3.8 Conclusions
The original motivation for investigating the binary Lennard-Jones system described here, was that it was already demonstrated to exhibit hopping . The first goal of the present work was to test if the hopping was still there if the cooling was done under (pseudo-) equilibrium conditions. Unfortunately there is not a single condition that is known to be sufficient for the system to be equilibrated. The different possibilities include: No drift in static properties (temperature, potential energy, pressure, etc.), long time dynamics being diffusive, relaxation functions such as $`F_s(q,t)`$ decaying to zero. These are all necessary conditions, but none of them are (known to be) sufficient. In the present work, it was found that the most sensitive condition (i.e. requiring the longest equilibration time), was that $`F_s(q,t)`$ decays to zero, *and* that it does so in the same manner for the three independent samples, i.e. with the relaxation time and stretching exponent being within reasonable agreement. Assuming that this condition is sufficient, it is concluded that the liquid is in equilibrium at all the temperatures presented here, and thus that the hopping found *is* a feature of the equilibrium liquid.
Although it was not a goal in itself to test the ideal mode coupling theory (MCT), the results of the simulations was compared to the asymptotic predictions of ideal MCT. This was done for two reasons; I) It provides a convenient way of comparing results with other simulations (e.g. Kob and Andersen ). II) It is “common” belief that ideal MCT breaks down when hopping dynamics takes over, and the critical mode coupling temperature $`T_c`$ thus constitutes an estimate of when hopping should start to dominate the dynamics. At first hand the last point seems to work fine; the best attempt at a power-law fit to $`\tau _\alpha `$ gives $`T_c=0.592\pm 0.004`$, which is close to the lowest of the temperatures simulated, where we clearly see hopping (figure 3.7). The fact that hopping also is present at the second lowest temperature, might then be taken as an indication that dynamical transition from the “MCT regime” to the “hopping regime” is not a sharp transition, but takes place over some temperature interval. However, the big problem with this scenario is the failure to identify any temperature range, where the asymptotic predictions of ideal MCT holds; There is no indication of TTSP, and the temperature dependence of the time-scales for diffusion and relaxation are clearly different.
In the full version of ideal MCT (as opposed to the asymptotic predictions used in the present work), the exponents of the power-laws discussed here4 are not free parameters, but are determined from the static structure factor . Calculating the exponents from the static structure factor is however a difficult numerical task, which has only been done for a few systems . This has not been attempted for the system used here. In an attempt to include the hopping dynamics in MCT, the so-called extended mode coupling theory has been developed introducing an “hopping-parameter”, as a extra fitting parameter . It has not been attempted to apply extended MCT to the system investigated here.
## Chapter 4 Inherent Dynamics
The dynamics of the model glass-forming liquid described in the previous chapter, is here analyzed in terms of its “inherent structures”, i.e. local minima in the potential energy. In particular, we compare the self part of the intermediate scattering function, $`F_s(q,t)`$, with its inherent counterpart $`F_s^I(q,t)`$ calculated on a time series of inherent structures. $`F_s^I(q,t)`$ is defined as $`F_s(q,t)`$ except that the particle coordinates at time $`t`$, are substituted with the particle coordinates in the corresponding inherent structure, found by quenching the equilibrium configuration at time $`t`$. We find that the long time relaxation of $`F_s^I(q,t)`$ can be fitted to stretched exponentials, as is the case for $`F_s(q,t)`$. Comparing the fitting parameters from $`F_s(q,t)`$ and $`F_s^I(q,t)`$ we conclude, that below a transition temperature, $`T_x`$, the dynamics of the system can be separated into thermal vibrations around inherent structures and transitions between these.
The main conclusions of this chapter can be found in paper III. In the present chapter we introduce the concepts of “energy landscape” and “inherent dynamics”, followed by a summary of paper III. Following this the data shown in paper II and some additional data is discussed, and a conclusion is given.
### 4.1 The Potential Energy Landscape
The potential energy landscape of an atomic system (i.e. particles without internal degrees of freedom) is simply the potential energy, as a function of the $`3N`$ particle coordinates. Letting $`𝐑`$ denote the $`3N`$ dimensional vector describing the state point of the system (i.e. its position in the $`3N`$ dimensional configuration space), we write the energy landscape as $`U(𝐑)`$, see figure 4.1. The behavior of the system can be viewed in terms of the state point, $`𝐑(t)`$, moving on the energy landscape surface, $`U(𝐑)`$. This surface contains a large number of minima, termed “inherent structures” by Stillinger and Weber . The inherent structures are characterized by zero gradients in the potential, and they are thus mechanically stable configurations. The inherent structures are separated by saddle points acting as energy barriers. Each inherent structure has a basin of attraction, in which a local minimization of the potential energy (a “quench”) will map the state point $`𝐑(t)`$ to the corresponding inherent structure $`𝐑^I(t)`$.
So far all we have done is defining the energy landscape. For a system like the binary Lennard-Jones mixture investigated in the previous chapter we know $`U(𝐑)`$ exactly; its simply the sum of the pair-potentials (equation 3.1)<sup>1</sup><sup>1</sup>1The fact that we know $`U(𝐑)`$ does not necessarily mean that we can find all the inherent structures. This can only be done for small systems ($`N`$ less than approxemately $`30`$ depending on other factors such as density and the potential), see eg. . This is not the approach we take here.. Of course having a well-defined quantity only helps, if it can tell us something about what we are interested in, i.e. in this case the dynamics of a glass-forming liquid. This is the main question we deal with in this chapter. First we note that for the simulations presented in the previous chapter, the dynamics is governed by Newton’s second law, which using the concepts introduced above can be written as:
$`{\displaystyle \frac{d^2}{dt^2}}𝐑(t)=𝐌^1U(𝐑(t))`$ (4.1)
where $`𝐌`$ is a ($`3N`$x$`3N`$) diagonal matrix, with the appropriate masses on the diagonal. In this sense the dynamics is *defined* by the energy landscape. This is however obviously not telling us anything new. What we want to know is, can we think of the dynamics in terms of the following scenario; the dynamics of the liquid is separated into (thermal) vibrations around inherent structures and (thermally activated) transitions between these. Presented with this scenario, one might very well ask: How can a system with constant total energy (e.g. the system simulated in the previous chapter) perform thermally activated processes? The answer is that one should think of the transitions between inherent structures as involving only a few (local) degrees of freedom, while the remaining degrees of freedom provides the heath bath.
In his classic paper from 1969 Goldstein argued that there exists a transition temperature, which we will term $`T_x`$, below which the flow of viscous liquids is dominated by potential barriers high compared to thermal energies, while above $`T_x`$ this is no longer true. Or in other words; below $`T_x`$ the dynamics is governed by the “vibrations plus transitions” scenario given above. Goldstein gave as a (very) rough estimate, that the shear relaxation time at $`T=T_x`$ is on the order of $`10^9`$ seconds. Later it was noted by Angell , that experimentally it is often found that the shear relaxation time is on the order of $`10^9`$ seconds at the mode coupling temperature, $`T_c`$. This lead to the argument that Goldstein’s transition temperature, $`T_x`$, is identical to the mode coupling temperature, $`T_c`$. A similar argument was recently given by Sokolov .
Recently, Sastry, DeBenedetti and Stillinger demonstrated that the onset of non-exponential relaxation ($`\beta <1`$) in a simulated glass-forming Lennard-Jones liquid is associated with a change in the system’s exploration of its potential energy landscape. In figure 4.2a we present data similar to the data presented in . $`E_{quench}(T)`$ is here the mean value of the energy of inherent structures, found by quenching configurations from a normal MD run equilibrated at the temperature $`T`$. Like in we find that this approaches a constant value at high temperatures, where the relaxation becomes exponential ($`\beta `$ 1, compare fig. 3.12). Similar results was found in for the quenched enthalpy (letting the volume change during the quench).
When plotting the total energy per particle, $`E_{tot}(T)/N`$ (in figure 3.2), we found that the specific heat capacity at constant volume, $`C_v(T)`$, increases as the system is cooled. In figure 4.2b this is shown explicitly, with $`C_v(T)`$ calculated from eq. 3.6 (using central differences). Also shown in figure 4.2b is the “quenched heat capacity”, $`C_{quench}(T)`$, calculated in the same way, but with $`E_{quench}(T)`$ substituted for $`E_{tot}(T)`$. Clearly the increase in $`C_v(T)`$ upon cooling is related to the increase in $`C_{quench}(T)`$.
### 4.2 The Inherent Dynamics
The basic idea of the “inherent dynamics” approach is the following (see figure 4.3); After equilibration at a given temperature, a time series of configurations, $`𝐑(t)`$, is produced by a normal MD simulation. Each of the configurations in $`𝐑(t)`$ is now quenched<sup>2</sup><sup>2</sup>2 In the present work this minimization was done using the conjugate gradient method , which uses a succession of line minimizations in configuration space to minimize the potential energy., to produce the corresponding time series of inherent structures, $`𝐑^I(t)`$. We now have two “parallel” time series of configurations. The time series $`𝐑(t)`$ defines the “true dynamics”, which is simply the normal (Newtonian) MD dynamics as presented in the previous chapter, described by $`r^2(t)`$, $`F_s(q,t)`$, etc. In a completely analogous way, we define the “inherent dynamics” as the dynamics described by the time series $`𝐑^I(t)`$. In other words: If a function describing the true dynamics (eg. $`r^2(t)`$ or $`F_s(q,t)`$) is calculated by $`f(𝐑(t))`$, then the corresponding function in the inherent dynamics (eg. $`r^2(t)^I`$ or $`F_s^I(q,t)`$) is calculated in exactly the same way, except using the time series of inherent structures: $`f(𝐑^I(t))`$.
If the true dynamics can be separated into vibrations around inherent structures, and transitions between these, as stated in the “vibrations plus transitions” scenario given above, then the inherent dynamics can be thought of as what is left after the thermal vibrations is removed from the true dynamics.
In the bottom part of figure 4.3 the inherent dynamics approach is applied to the trajectory of the hopping particle, which was shown in figure 3.1; All the configurations that were used to plot the true trajectory were quenched, and the position of the particle in the resulting time series of inherent structures is plotted as the “inherent trajectory”. The quenching procedure is seen to remove the vibrations in the true trajectory. The motion that seems to be left from the vibrations (eg. a jump from $`x2.2`$ to $`x2.4`$) will be discussed in section 4.5.
Paper II and III explores the possibilities of comparing the true and inherent dynamics of the model glass-former described in the previous chapter. The concept of inherent dynamics was first introduced in paper II (without using that name). In that paper the true and inherent versions of the mean square displacement and the van Hove correlation function, was compared on a qualitatively level. Paper III compares the true and inherent version of the self part of the intermediate scattering function, which can be done on a quantitative level by means of stretched exponentials (see section 3.5). This turns out to have important consequences, which is why we here discuss this paper first.
### 4.3 Paper III
In this paper the concept of inherent dynamics is applied to the self part of the intermediate scattering function, i.e. we compare $`F_s(q_{max},t)`$, with its inherent counterpart $`F_s^I(q_{max},t)`$. As explained above $`F_s^I(q,t)`$ is defined in the same way as $`F_s(q,t)`$, except that the normal particle coordinates, $`𝐫_j(t)`$, are substituted by the corresponding coordinates in the inherent structures, $`𝐫_j^I(t)`$ (compare equation 3.14)<sup>3</sup><sup>3</sup>3 $`𝐫_j^I(t)`$ is here the 3-dimensional vector describing the position of the $`j`$’th particle in the inherent structure $`𝐑^I(t)`$.:
$`F_s^I(𝐪,t)`$ $``$ $`\mathrm{cos}𝐪(𝐫_j^I(t)𝐫_j^I(0))`$ (4.2)
In figure 4b in paper III $`F_s^I(q_{max},t)`$ is plotted for the A particles, at the 8 temperatures studied. Here we plot the similar data for the B particles in figure 4.4. In both cases we find that the long time relaxation of $`F_s^I(q_{max},t)`$ can be fitted to stretched exponentials, like is the case for $`F_s(q_{max},t)`$. This is an important finding, since it enables us to make an quantitative comparison between $`F_s(q_{max},t)`$ and $`F_s^I(q_{max},t)`$, by comparing the two sets of fitting parameters. In the following, the set of fitting parameters for $`F_s(q_{max},t)`$ is denoted $`\{\tau _\alpha ,\beta ,f_c\}`$, while the corresponding set for $`F_s^I(q_{max},t)`$ is denoted $`\{\tau _\alpha ^I,\beta ^I,f_c^I\}`$
The key question now is: If the true dynamics follows the “vibrations plus transitions” scenario, i.e. if the dynamics can be separated into vibrations around inherent structures, and transitions between these, how is $`\{\tau _\alpha ^I,\beta ^I,f_c^I\}`$ expected to be related to $`\{\tau _\alpha ,\beta ,f_c\}`$? To answer this question, we assume that the initial relaxation in $`F_s(q,t)`$ is due to vibrations, which is the widely accepted explanation (see section 3.5). If this is the case, then we expect the quenching procedure to remove the initial relaxation (since it removes the vibrations), which means that $`F_s^I(q,t)`$ can be thought of as $`F_s(q,t)`$ with the initial relaxation removed. This in turn means, that $`F_s^I(q,t)`$ should be identical to the long time relaxation of $`F_s(q,t)`$, but rescaled to start at unity (by definition): $`\{\tau _\alpha ^I,\beta ^I,f_c^I\}=\{\tau _\alpha ,\beta ,1\}`$.
The identity of the relaxation times, and the stretching parameters, $`\{\tau _\alpha ^I,\beta ^I\}=\{\tau _\alpha ,\beta \}`$, will probably be true for any “coarse-graining” of the dynamics we might apply; if we eg. decide to add a small random displacement to all the particles (instead of quenching), we would still expect the long time relaxation to have the same shape and characteristic time, i.e. $`\{\tau _\alpha ^I,\beta ^I\}=\{\tau _\alpha ,\beta \}`$. It is when we find $`f_c^I=1`$ we know that the procedure we have applied removes the vibrations (or in more general terms: removes the part of the dynamics responsible for the initial relaxation).
In figure 5 in paper III the fitting parameters for $`F_s(q_{max},t)`$ and $`F_s^I(q_{max},t)`$ are compared for the A particles. Here we show similar data for the B particles in figure 4.5. In both cases we find for all temperatures $`\tau _\alpha ^I\tau _\alpha `$. The stretching parameters are more difficult to compare at high temperatures, but they become identical at low temperatures, $`\beta ^I=\beta `$ (within the error-bars). Whereas $`f_c`$ is roughly constant as a function of temperature, $`f_c^I`$ is clearly seen to approach unity, as the system is cooled. We have thus found evidence, that the the vibrations plus transitions scenario holds below a transition temperature $`T_x`$, as argued by Goldstein. We estimate that $`T_x`$ is close to (or just below) the lowest of the temperatures simulated ($`T=0.591\pm 0.002`$).
At the transition we find $`\tau _\alpha 310^3`$, which in “Argon units” corresponds to $`\tau _\alpha 10^9`$ seconds (see section 3.1), i.e. Goldstein’s estimate of the shear relaxation time at $`T_x`$. This agreement is however probably “to perfect”; As mentioned above Goldstein’s estimate is very rough, and it regards the shear relaxation time and not $`\tau _\alpha `$. Consequently an agreement better than “orders of magnitude” should probably be considered a coincidence.
Having found evidence that there exists a transition temperature, $`T_x`$, and estimated the (approximate) position of this, it is natural to proceed to check Angell’s proposal, that $`T_xT_c`$. We find that both estimated values for $`T_c`$ ($`0.592\pm 0.005`$ from relaxation times and $`0.574\pm 0.005`$ from diffusion) are in the temperature range where $`f_c^I`$ is approaching unity. To check the proposition that $`T_xT_c`$ on the system used here, is of course somewhat problematic, since it doesn’t conform very well to the predictions of the mode coupling theory, as discussed in the previous chapter. It should be noted however, that the arguments given by Angell (and Sokolov), only relates to $`T_c`$ as the temperature where power-law fits to experimental data tends to break down, i.e. the “usage” of MCT in this argument is similar to the way we have estimated $`T_c`$ in the previous chapter, and does not require e.g. time-temperature super-position.
At the end of paper III we present results on the nature of transitions between inherent structures. These will be discussed in section 4.5.
### 4.4 Paper II
In paper II the concept of inherent dynamics is applied to the mean square displacement, $`\mathrm{\Delta }r^2(t)`$, and the distribution of particle displacements,
$`4\pi r^2G_s(r,t)`$.
Comparing $`\mathrm{\Delta }r^2(t)^I`$ to $`\mathrm{\Delta }r^2(t)`$ (figure 1 in paper II), we find that the quenching procedure removes the plateau seen in $`\mathrm{\Delta }r^2(t)`$. This is taken as (qualitative) evidence for the “cage”-explanation for the plateau (see section 3.3). This conclusion is consistent with the conclusion we drew from $`F_s^I(q,t)`$ in the previous section, since the “caging” is what we call vibrations in configuration space.
Comparing $`4\pi r^2G_s^I(r,t)`$ and $`4\pi r^2G_s(r,t)`$ (figure 2 in paper II), we find that the hopping peak is slightly sharper in $`4\pi r^2G_s^I(r,t)`$, and that the first peak is moved to the left. A point that is *not* made in paper II is the following; if the particles only moved by either “rattling” in their local cages, or hopping approximately an inter-particle distance, then one would expect the first peak seen in $`4\pi r^2G_s(r,t)`$ to be quenched to a delta-peak at $`r=0`$ in $`4\pi r^2G_s^I(r,t)`$. Figure 2b in paper II shows that this is clearly not the case. The reason for this will be discussed in the next section.
### 4.5 Transitions between inherent structures
Having established that our lowest temperature (T=0.59) is close to $`T_x`$ (i.e. the dynamics can be thought of as separated into vibrations around inherent structures and transitions between these) we are now interested in studying the nature of transitions between inherent structures at that temperature. We have identified 440 such transitions, by quenching the true MD configurations every $`0.1\tau `$ (i.e. every 10 MD-steps), and looking for signatures of the system making a transition from one inherent structure to another. In figure 4.6a is shown the energy of the inherent structures as a function of time, $`E_{quench}(t)`$. As expected $`E_{quench}(t)`$ is found to be constant for some time-intervals, and then jump to another level, i.e. the system has made a transition from one inherent structure to another. In figure 4.6b is plotted as a function of time, the distance in configuration space between two successive quenched configurations :
$$\mathrm{\Delta }R^I(t)|𝐑^I(t+0.1)𝐑^I(t)|=\sqrt{\underset{j=1}{\overset{N}{}}\left(𝐫_j^I(t+0.1)𝐫_j^I(t)\right)^2}$$
(4.3)
Each jump in $`E_{quench}(t)`$ is associated with a peak in $`\mathrm{\Delta }R^I(t)`$, indicating the system has moved to a different inherent structure. A transition might in principle occur between two inherent structures with exactly the same (quenched) energy. Such a transition would not be seen in $`E_{quench}(t)`$, and consequently we use $`\mathrm{\Delta }R^I(t)`$ to identify transitions. The distribution of $`\mathrm{log}_{10}(\mathrm{\Delta }R^I(t))`$ is shown in figure 4.7. The distribution is seen to have two separated peaks. The peak to the left (centered around $`\mathrm{log}_{10}(\mathrm{\Delta }R^I(t))3)`$ is due to numerical uncertainties; two configurations within the same basin of attraction is *not* quenched to the exactly the same configuration. The peak on the right is the one containing the physically interesting transitions, which is seen to have $`\mathrm{\Delta }R^I`$ on the order of unity. In the following we use the condition $`\mathrm{\Delta }R^I>0.1`$ to identify the (physically interesting) transitions.
Figure 4.7 shows that most of the quenches results in *not* finding a transition (left peak), and only a small fraction results in actually finding a transition (right peak). Or in actual numbers: Doing 7500 quenches we found 440 transitions. The quenching procedure takes a considerable amount of time (corresponding to approximately 1000 MD steps), which is why we haven’t simply continued this procedure to find a larger number of transitions. The data presented in the following is averaged over the 440 transitions found using the “brute force” method described above, and should be considered preliminary. (NOTE: In paper III results from 12000 transitions are reported. The results are similar to the results presented here, except of course with less noise).
For each transition, we monitor the displacements of the particles from one inherent structure to the other. The distribution of all particle displacements is shown in Fig. 4.8. While many particles move only a small distance ($`r<0.2`$) during the transition, a number of particles move farther, and in particular, we find that the distribution for $`r>0.2`$ is to a good approximation exponential. At present we have no explanation for this. The dotted curve is a fit to a power-law with exponent $`5/2`$, which is a prediction from linear elasticity theory , describing the displacements of particles in the surroundings of a local “event”. This power-law fit does not look very convincing by it self, but we note that the exponent was not treated as a fitting parameter (i.e. only the prefactor was fitted), and the power-law is *expected* to break down at small displacements, since these corresponds to distances far away from the local event, and is thus not seen in our finite sample. From the change in behavior of $`p(r)`$ at $`r0.2`$, it is reasonably to think of particles with displacements larger than $`0.2`$, as those taking part in the local event, and the rest of the particles as being in the surroundings, adjusting to the local event. Using this definition it is found that on average approximately 10 particles participate in an event.
Figure 4.8 has two important consequences with regards to points discussed earlier in this thesis; i) The distribution of particle displacements during transitions shows no preference for displacement of the average inter-particle distance ($`1`$ in the used units). This shows that the hopping indicated by the secondary peak in $`4\pi r^2G_s(r,t)`$ (figure 3.7 and 3.9) at low temperatures is not due to transitions over single energy barriers. A (correlated) sequence of the these is needed, to “build up” the secondary peak. This is consistent with the behavior seen in the inherent trajectory in figure 4.3; The jump does not happen in one step, but through a number of “intermediate” inherent structures. ii) Particles in the surroundings of a local event are displaced small distances, to adjust to the larger displacements occurring in the local region of the event itself. This kind of motion is very hard to detect in the true dynamics, since it is dominated by the thermal vibrations. Presumably this kind of motion is the reason why the inherent trajectory in figure 4.3 still retains some motion “within” the vibrations; as a consequence of an event in the surroundings, the particle starts vibrating around a position that is slightly displaced. This view of the dynamics is also consistent with the fact, that the first peak in $`4\pi r^2G_s^I(r,t)`$ is not a delta function in $`r=0`$ (see discussion in section 4.4).
From visual inspection of a number of the identified transitions it was found, that these are cooperative and string-like in nature. By visual inspection is here meant, that the position of particles moving more than $`0.2`$ during the transition, is plotted before and after the transition, see figure 4.9. String-like motion has been found to be an important part of the dynamics of glass-forming liquids. It is the natural consequence of particles hopping to positions previous occupied by other particles (as concluded from figure 3.10), and it was found (and quantified) in the Kob & Andersen system, when looking at the “mobile” particles . In paper II strings was also found when looking at how particles was displaced (in the inherent structures) during the time interval $`t^{}`$ (figure 3 in paper II). In paper II this was described as “vacancy hopping”; one particle jumps, leaving room for another particle to jump, etc. The finding that also transitions over energy barriers are associated with strings, indicates that the vacancy hopping interpretation might not be correct; it seems to indicate that (at least some of) the strings are really cooperative in nature, i.e the particles in the string move at the same time. Further (quantitative) investigations are obviously needed to answer this question.
### 4.6 Conclusion
In this chapter we have presented results from analyzing the dynamics of a model glass-forming liquid in terms of its potential energy landscape. We did so by introducing the new concept of “inherent dynamics”, which can be thought of as a course-graining of the true dynamics, where the part of the dynamics related to vibrations around single inherent structures is removed.
Comparing the self intermediate scattering function, $`F_s(q,t)`$, with its inherent counterpart, $`F_s^I(q,t)`$, we found direct numerical evidence for the existence of a transition temperature, $`T_x`$, below which the true dynamics is separated into vibrations around inherent structures and transitions between these (the “vibrations plus transitions scenario”). We thus confirm the “energy landscape” picture, which is (at least) 30 years old. Given the fact the energy landscape *does* exist (since its simply the potential energy as function of the particle coordinates) and it *does* have a number of local minima (inherent structures), it is not surprising that the dynamics becomes dominated by the energy barriers at sufficiently low temperatures. What we have done here using the concept of inherent dynamics, is to provide direct numerical evidence for this, *and* we have shown that this regime can be reached by (pseudo-) equilibrium molecular dynamics (for the particular system investigated here). To our knowledge this is the first time such evidence has been presented.
The fact that we have been able to cool the system, under equilibrium conditions, to temperatures where the separation between vibrations around inherent structures and transitions between these is (almost) complete, means that it now makes sense to study the individual transitions over energy barriers, since these in this regime are “significant”. There is a lot of interesting questions to investigate regarding these transitions, and we have here only investigated a few of these. Specifically, we have not determined the energy barriers, but only compared the two inherent structures involved in a particular transition. This is an obvious point for further investigations.
## Chapter 5 The Symmetric Hopping Model
The symmetric hopping model is introduced, and three analytical approximations for calculating the frequency dependent diffusion coefficient, $`D(s)`$, in the extreme disorder limit (low temperatures) of the model is described; the Effective Medium Approximation (EMA), the Percolation Path Approximation (PPA), and the Diffusion Cluster Approximation (DCA). DCA is a new approximation, developed by my supervisor Jeppe C. Dyre (See paper IV and V). Two numerical methods for calculating $`D(s)`$ in the extreme disorder limit is discussed. The first method is derived from the mean square displacement and is equivalent with the method of the ac Miller-Abrahams (ACMA) electrical equivalent circuit. The second method (VAC) is derived from the velocity auto correlation and is a new method. Numerical results using the VAC method are compared to the three analytical approximations, and previous results from the ACMA method.
The main results in this chapter are the development of the VAC method (section 5.5), and the numerical results it leads to (section 5.6 and 5.7). Results in this chapter are published in paper IV and V.
### 5.1 The Symmetric Hopping Model
The symmetric hopping model is defined in the following way: A particle ’lives’ on the sites of a $`𝒟`$-dimensional regular lattice (see figure 5.1 for a 1-dimensional illustration), where all the sites has the same energy (which we set to 0). The particle jumps over energy barriers connecting nearest neighbor sites, with jump rates (probability per unit time) given by $`\mathrm{\Gamma }(ki)=\mathrm{\Gamma }_0\mathrm{exp}(\beta E_{ki})`$, where $`\mathrm{\Gamma }_0`$ is the (constant) “attack-frequency”, $`\beta (k_BT)^1`$ and $`E_{ki}`$ is the energy barrier between the two sites. The energy barriers are chosen randomly, from a probability distribution, $`p(E)`$, (to be specified). Jump rates between sites that are not nearest neighbors are zero. It follows from the above, that $`\mathrm{\Gamma }(ik)=\mathrm{\Gamma }(ki)`$, i.e. the jump-rates are symmetric. In the following we will use $`\mathrm{\Gamma }(E)\mathrm{\Gamma }_0\mathrm{exp}(\beta E)`$, and denote the lattice constant $`a`$.
If all the energy barriers are identical, i.e. $`p(E)=\delta (EE_0)`$, the model describes diffusion in an ordered structure, and one finds normal diffusive behavior:
$`\mathrm{\Delta }X^2(t)=2Dt`$ (5.1)
where $`\mathrm{\Delta }X^2(t)`$ is the mean square displacement along the x-direction, and $`D=a^2\mathrm{\Gamma }(E_0)`$ is the diffusion coefficient . The average in equation 5.1 can either be a time-average or a ensemble average. We will in the following use the ensemble average, which is characterized by all the sites having the same probability (since they have the same energy). Instead of ensembles it is convenient to think of a (large) number of independent particles moving around in the sample.
In a sample where the energy barriers are not identical, equation 5.1 does not hold at all time scales. Picture for example a 1-dimensional sample with mostly small energy barriers, $`E_{small}`$, and a few much larger energy barriers, $`E_{large}`$, (see eg. ). At small time scales most of the particles will “think” they are living on an ordered sample with only the small energy barriers (since they haven’t yet encountered a large energy barrier), and the ensemble will follow equation 5.1 with a large diffusion coefficient ($`\mathrm{\Gamma }(E_{small})`$). However, at long time scales the particles will start to “feel” the effect of the large energy barriers, which will slow down the diffusion. At very long time and length scales the sample will appear homogeneous and the system again become diffusive, but now with a small diffusion coefficient ($`\mathrm{\Gamma }(E_{large})`$).
The deviations from equation 5.1 can be quantified using the time dependent diffusion coefficient , $`D(t)`$, or the frequency dependent diffusion coefficient , $`D(s)`$ (with $`s`$ being the Laplace frequency, $`s=i\omega `$):
$`D(t)`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{d}{dt}}\mathrm{\Delta }X^2(t)`$ (5.2)
$`D(s)`$ $`=`$ $`s{\displaystyle _0^{\mathrm{}}}e^{st}D(t)𝑑t={\displaystyle \frac{s^2}{2}}{\displaystyle _0^{\mathrm{}}}e^{st}\mathrm{\Delta }X^2(t)𝑑t`$ (5.3)
Note, that for diffusion in a sample with identical energy barriers, equation 5.1 leads to $`D(s)=D(t)=D`$.
In a disordered sample one in general finds that the system becomes diffusive at long time scales (corresponding to low frequencies) when particles are moving much longer than the appropriate correlation length (assuming that such a correlation length exists); $`D_0D(s0)=D(t\mathrm{})`$. In one dimension one finds : $`D_0=a^2\mathrm{\Gamma }^1^1`$, where the average is over the distribution of jump rates, $`\mathrm{\Gamma }`$. In higher dimensions no such simple expression exists, but for $`\beta \mathrm{}`$ one finds : $`D_0\mathrm{\Gamma }(E_c)`$, where $`E_c`$ is the so-called percolation energy (see section 5.3.2).
At infinitely short time particles never jumps more than once, and one finds (where $`\mathrm{\Gamma }_L`$ and $`\mathrm{\Gamma }_R`$ are respectively the jump rates to the left and to the right of a given site):
$`\mathrm{\Delta }X^2(\mathrm{\Delta }t)`$ $`=`$ $`a^2\mathrm{\Gamma }_L\mathrm{\Delta }t+a^2\mathrm{\Gamma }_R\mathrm{\Delta }t`$ (5.4)
$`=`$ $`2a^2\mathrm{\Gamma }\mathrm{\Delta }t`$ (5.5)
In the high-frequency limit, we thus find: $`D_{\mathrm{}}D(s\mathrm{})=D(t0)=a^2\mathrm{\Gamma }`$.
In and paper I, the symmetric hopping model was treated as a model for frequency dependent conduction in glasses. In that context the particles represent non-interacting charge carriers (ions or electrons), and the lattice sites represents the positions in the glass where the charge carriers can reside. The frequency dependent conductivity, $`\sigma (s)`$, is related to $`D(s)`$ by the generalized Einstein relation :
$$\sigma (s)=\frac{e^2n}{k_BT}D(s)$$
(5.6)
where $`n`$ is the density of particles, and $`e`$ is the charge of the particles.
In the present work we shift the attention slightly, and treat the symmetric hopping model more generally as a model for diffusion in disordered media. We focus on the “extreme disorder limit”, i.e. low temperatures, where the model is found to exhibit universal behavior; $`D(s)`$ becomes independent (suitably scaled) of the temperature *and* the chosen probability distribution for the energy barriers, $`p(E)`$. In the extreme disorder limit, it is *not* feasible to simulate the model using standard Monte Carlo (MC) techniques ; Due to the large difference in jump rates ($`e^{\beta E}`$), the particle will jump many times over small energy barriers, and only “sample” the rest of the lattice on much larger time scales. Although this reflects the physics of the model in the extreme disorder limit, it makes MC methods unsuitable in this limit.
In the following we set the attack frequency, $`\mathrm{\Gamma }_0=1`$, the lattice constant, $`a=1`$, and Boltzmann’s constant, $`k_B=1`$, thus defining the scales for time, length, and energy respectively.
### 5.2 The Master Equation
The starting point for both the analytical and the numerical methods, is the master equation for the model; Let $`P(i,t|j,0)`$ denote the probability of the particle being at site $`i`$ at time $`t`$, given that it was at site $`j`$ at $`t=0`$. The master equation for the system is then:
$$\frac{dP(i,t|j,0)}{dt}=\gamma _iP(i,t|j,0)+\underset{k}{}\mathrm{\Gamma }(ki)P(k,t|j,0)$$
(5.7)
where $`\gamma _i_k\mathrm{\Gamma }(ik)`$. The first term on the right hand side is the probability flow out of site $`i`$, and the second term is the flow into site $`i`$. Defining the ($`N^𝒟`$x$`N^𝒟`$) matrix $`𝐏(t)`$ with the components $`𝐏_{ij}(t)=P(i,t|j,0)`$, we write the master equation on matrix form:
$$\frac{d}{dt}𝐏(t)=\mathrm{𝐇𝐏}(t)$$
(5.8)
$`𝐇`$ is here a ($`N^𝒟`$x$`N^𝒟`$) matrix containing the jump rates; $`𝐇_{ii}=\gamma _i`$ and $`𝐇_{ik}=\mathrm{\Gamma }(ki)`$. Note that only $`2𝒟+1`$ elements in each row in $`𝐇`$ are different from zero (the diagonal and the $`2𝒟`$ elements corresponding to nearest neighbors); $`𝐇`$ is sparse. This is an essential feature in the numerical methods; without it we wouldn’t be able to treat large enough sample sizes.
Taking the Laplace transform of equation 5.8 we get:
$$s𝐆(s)𝐏(t=0)=\mathrm{𝐇𝐆}(s)$$
(5.9)
where the $`𝐆_{ij}(s)`$ is the Green’s function, i.e. the Laplace transform of $`𝐏_{ij}(t)`$: $`𝐆_{ij}(s)_0^{\mathrm{}}𝐏_{ij}(t)e^{st}𝑑t`$. The initial condition, $`𝐏(t=0)`$, is given by the identity matrix, $`𝐈`$, and we thus find:
$$(s𝐈𝐇)𝐆(s)=𝐈𝐆(s)=(s𝐈𝐇)^1$$
(5.10)
Before proceeding to discuss the analytical approximations and the numerical methods, we present a “preview” of the dynamics of the symmetric hopping model in figure 5.2. A 2 dimensional sample was set up using a box-distribution of barrier energies; $`p(E)=1,0E1`$. The time evolution of $`𝐏(t)`$ was determined by a discrete version of the master equation (eq. 5.8):
$$𝐏(t+\mathrm{\Delta }t)=𝐏(t)+\mathrm{\Delta }t\mathrm{𝐇𝐏}(t)=(𝐈+\mathrm{\Delta }t𝐇)𝐏(t)$$
(5.11)
A time step of $`\mathrm{\Delta }t=0.1`$ was used, and as initial condition a particle was placed at a particular site. In figure 5.2a $`𝐏(t=2)`$ is shown for $`\beta =0`$, i.e. infinitely high temperature. In this limit all the jump rates are identical (equal to unity), and as expected the probability distribution spreads in a symmetric manner. In figure 5.2b $`𝐏(t=10^5)`$ is shown for $`\beta =40`$, with the same starting position, and same energy landscape (i.e. same set of energy barriers) as in figure 5.2a. The time at which the two temperatures are compared was chosen so that the probability of finding the particle at the starting point was the same ($`0.04`$). The qualitative difference between the dynamics at high and low temperature is evident; At the low temperature (figure 5.2b) the spread of the probability distribution $`𝐏(t)`$ is highly irregular, as a consequence of the particles “preferring” low energy barriers and avoiding high energy-barriers. The particular structure of $`𝐏(t)`$ at the low temperature depends strongly on where the particle was started. If it was started in one of the two enclosed empty sites (white in figure 5.2b) it would be stuck there on the time scale used here, since these must be connected with high energy-barriers to the surroundings (otherwise they would not be empty in figure 5.2b).
### 5.3 Analytical approximations
In this section we’ll briefly describe the 3 analytical approximations which will be compared to the numerical results, in section 5.6 and 5.7.
#### 5.3.1 Effective Medium Approximation (EMA)
In EMA the disordered sample is replaced by an ordered sample (the “effective medium”), where all the jump rates are replaced by an “effective” jump-rate, $`\mathrm{\Gamma }_E(s)`$. The value of $`\mathrm{\Gamma }_E(s)`$ is determined by a self-consistency condition; A single jump rate in the effective medium is replaced by its value in the disordered sample, and an average is performed over the distribution of jump rates. The condition that this average should be equivalent with the effective medium leads to the following condition :
$`{\displaystyle \frac{D\mathrm{\Gamma }}{𝒟D(D\mathrm{\Gamma })(1s\stackrel{~}{G})}}_\mathrm{\Gamma }`$ $`=`$ $`0`$ (5.12)
where $`\stackrel{~}{G}`$ is the diagonal element of the Green’s function for the effective medium, which depends on the spatial dimension, $`𝒟`$. In the extreme disorder limit ($`\beta \mathrm{}`$) one finds for $`𝒟2`$:
$`\stackrel{~}{D}\mathrm{ln}(\stackrel{~}{D})=\stackrel{~}{s}`$ (5.13)
where $`\stackrel{~}{D}D(s)/D_0`$, and $`\stackrel{~}{s}s/s_c`$, where $`s_c`$ is a suitably defined characteristic frequency. The EMA thus predicts universality in the extreme disorder limit; Properly scaled $`D(s)`$ becomes independent of temperature *and* the distribution of energy barriers, $`p(E)`$.
In the predictions of EMA was compared with numerical results using the ACMA method (described in section 5.4) in 2 dimensions. The existence of universality was confirmed by the numerical results (for 4 different energy-distributions), but the shape of $`\stackrel{~}{D}(\stackrel{~}{s})`$ was found to deviate from the one predicted by EMA. Or to be more specific: EMA predicts a shape of the universal $`\stackrel{~}{D}(\stackrel{~}{s})`$ that is growing to rapidly at the onset of frequency dependence; it is to “sharp” (figure 5 in ). The EMA prediction for $`\mathrm{\Delta }X^2(t)`$ is discussed in .
#### 5.3.2 Percolation Path Approximation (PPA)
In the extreme disorder limit, the low-frequency diffusion is expected to be dominated by percolation ; when diffusing over long distances (i.e. in the low-frequency limit) the particles “chose” to do so by jumping over energy-barriers that are as small as possible. The largest energy barrier that a particle has to cross to move through the sample is given by the “percolation energy”, $`E_c`$, defined by :
$`{\displaystyle _0^{E_c}}p(E)𝑑E=p_c`$ (5.14)
where $`p_c`$ is the percolation threshold for bond percolation (for an introduction to percolation theory, see ). In a 2 dimensional regular lattice $`p_c`$ equals $`1/2`$ (exact) and in a 3 dimensional cubic sample $`p_c`$ equals $`0.2488`$ (approximated). Equation 5.14 can be interpreted as follows; In a given sample mark the energy-barriers (bonds) starting with the smallest energy-barrier, then the next smallest etc. At some point the marked energy-barriers will “percolate” i.e. make a cluster that stretches through the whole sample (the “percolation cluster”). The highest energy needed to get percolation is (for an infinite sample) $`E_c`$. For the box distribution of energy barriers used in the present work ($`p(E)=1,0E1`$), the percolation energy equals the (bond) percolation threshold: $`E_c=p_c`$. The percolation cluster is a fractal, with fractal dimension 1.9 and 2.5 in 2 and 3 dimensions respectively .
In paper I we argue, that the reason why EMA does not predict the shape of the universal curve $`\stackrel{~}{D}(\stackrel{~}{s})`$ in the extreme disorder limit correctly might be, that it replaces the disordered sample by an effective media of the *same* dimension, whereas the actual low frequency diffusion happens on a cluster of lower dimensionality. We will here term this cluster “the diffusion cluster”, and its fractal dimension $`D_f`$ (this is *not* the same as the percolation cluster, see below). PPA can be considered as a first attempt at trying to incorporate the fractal dimension of the diffusion cluster into an analytical approximation.
The percolation cluster contain “dead-ends”, which contributes little to the low frequency diffusion. Removing dead-ends from the percolation cluster leaves us with the “backbone”, with fractal dimension 1.6 and 1.7 in 2 and 3 dimensions respectively . In paper I we argue, that $`D_f`$ is even lower, since the backbone contains loops, where one of the branches usually will be preferred (since the jump-rates depend strongly on the energy barriers in the extreme disorder limit). In PPA the extreme view that $`D_f=1`$ is taken, i.e. that the (low frequency) diffusion takes place on 1-dimensional “percolation paths”. With this assumption we in paper I arrive at the PPA approximation, given by:
$`\sqrt{\stackrel{~}{s}\stackrel{~}{D}}\mathrm{ln}\left(1+\sqrt{\stackrel{~}{s}\stackrel{~}{D}}\right)=\stackrel{~}{s}`$ (5.15)
PPA thus predicts universality in the extreme disorder limit, like is the case for EMA. In paper I , we find that PPA agrees better with the numerical results than EMA. This is especially true in 3 dimensions. The agreement is however not perfect, and a number of problems remain unanswered (see last section in paper I). The PPA prediction for $`\mathrm{\Delta }X^2(t)`$ is discussed in .
#### 5.3.3 The Diffusion Cluster Approximation (DCA)
The Diffusion Cluster Approximation can be thought of as a refinement of PPA (paper IV and V); instead of setting $`D_f=1`$, this is in DCA left as a parameter in the model. Using the approach of EMA, but with an fractal effective medium with fractal dimension $`D_f`$, one finds:
$`\mathrm{ln}(\stackrel{~}{D})=\left({\displaystyle \frac{\stackrel{~}{s}}{\stackrel{~}{D}}}\right)^{D_f/2}`$ (5.16)
This expression is limited to $`1<D_f<2`$. For $`D_f2`$ DCA reduces to the EMA prediction (eq. 5.13), and for $`D_f1`$ it is undefined. Since we do not have an independent estimate of $`D_f`$ we will in the present work treat it as a fitting parameter. This must obviously be taken into consideration, when comparing with EMA and PPA, since these have no fitting parameters.
From the arguments given above, we expect $`D_f`$ to be limited above by the fractal dimension of the backbone. Furthermore we expect $`D_f`$ to limited below by the fractal dimension of the so-called “red bonds”, which are the singly connected bonds on the backbone, i.e. if a red bond is removed the backbone is broken into 2 parts. The fractal dimension of the red bonds are $`3/4`$ and $`1.14`$ in 2 and 3 dimensions respectively.
### 5.4 The ACMA Method
In this section we briefly describe the ACMA method, which was used to obtain the numerical results reported in and paper I.
Defining $`X_i`$ as the x-coordinate of site $`i`$ we can write the frequency dependent diffusion coefficient (equation 5.3):
$`D_x(s)`$ $`=`$ $`{\displaystyle \frac{s^2}{2}}{\displaystyle _0^{\mathrm{}}}e^{st}{\displaystyle \frac{1}{N^d}}{\displaystyle \underset{i,j}{}}(X_iX_j)^2P(i,t|j,0)dt`$ (5.17)
$`=`$ $`{\displaystyle \frac{s^2}{2N^d}}{\displaystyle \underset{i,j}{}}(X_iX_j)^2G(i,s|j)`$ (5.18)
The subscript $`x`$ in $`D_x(s)`$ is here used to emphasize that the diffusion coefficient here is calculated from the mean square displacement in the x-direction. In more than 1 dimension similar expressions obviously hold for other directions.
Equation 5.10 and equation 5.18 together constitutes a method for computing D(s); calculate $`G(i,s|j)`$ by inverting $`(s𝐈𝐇)`$ (equation 5.10) and insert the result in equation 5.18 .
In my master thesis it was demonstrated, that the method described above is equivalent to the method of the ac Miller-Abrahams (ACMA) electrical equivalent circuit . In this method a large network of linear electric components are set up (see figure 5.3 for an illustration of the 1-dimensional case). After eliminating all the ’internal sites’ (those without numbers in figure 5.3) by the so-called generalized Star-Mesh transformation, $`\sigma (s)`$ ($`D(s)`$, see equation 5.6) is calculated as a weighted sum over the effective conductances between the ’external sites’ (those with numbers in figure 5.3). The elimination process described above is equivalent with Gauss-Jordan elimination on the matrix $`(s𝐈𝐇)`$ , and the Fogelholm algorithm for the order in which the internal sites are eliminated is equivalent with the minimum degree pivoting algorithm : At each step eliminate the site/row with the smallest number of connections/elements.
A number of ’tricks’ is used in the ACMA method, to improve efficiency over a straight forward computation of first equation 5.10 and then equation 5.18. The most important trick is that the (dense) matrix $`𝐆(s)`$ is not explicitly calculated, but instead a ’condensed’ version where columns corresponding to the sites with the same x-coordinate are added together. Thus instead of solving for $`N^𝒟`$ right-hand sides (inverting $`s𝐈𝐇`$) only $`N`$ right-hand sides are used.
In the simulations using the ACMA method presented in and paper I ’perfect electrodes’ were used in the x-direction (see figure 5.3). This ensures that $`D(s)`$ goes to a constant value as $`s0`$, corresponding to $`\mathrm{\Delta }X^2(t)`$ being proportional to $`t`$, as $`t\mathrm{}`$, i.e. the system being diffusive. However, with this modification the model does *not* correspond directly to the solution of a master equation. This is evident from the fact that the mean square displacement as calculated in equation 5.17 is finite for any finite sample, and thus it can not be proportional to $`t`$ at long times.
To what extent the ’perfect electrodes’ mimics the effect of real electrodes when experimentally measuring the frequency dependent conductivity, $`\sigma (s)`$, will not be discussed here. What we are after here, is to calculate the *bulk* value of the frequency dependent diffusion coefficient $`D(s)`$ in the extreme disorder limit.
### 5.5 The Velocity Auto Correlation (VAC) method
In this section the Velocity Auto Correlation (VAC) method is developed. Its main advantages over the ACMA method is, that it can be (and is) used with periodic boundary conditions, and still give a diffusive regime for $`s0`$ (i.e. $`t\mathrm{}`$). Another advantage over the ACMA method (as it was implemented in the previous work), is that the VAC method is formulated in terms of sparse matrices, which means that standard methods for solving these can be used<sup>1</sup><sup>1</sup>1The numerical results presented here was done using Matlab version 5.3.0. . The differences between the two methods will be discussed further in section 5.8.
To derive the VAC method we express the diffusion coefficient in terms of the velocity auto correlation function :
$$D_x(s)=_0^{\mathrm{}}v(t)v(0)_xe^{st}𝑑t$$
(5.19)
$`v(t)`$ is here the velocity (in the x-direction) of the particle at time t. The motion of the particle in the symmetric hopping model is instantaneous; when it jumps from one site to another, this happens in a infinitesimally small time-interval, $`\mathrm{\Delta }t`$. We may thus assign the constant velocity $`\pm a/\mathrm{\Delta }t`$ (where we briefly reintroduce the lattice constant, $`a`$) to the particle in the time interval, $`\mathrm{\Delta }t`$, ensuring that it moves one lattice constant, $`a`$, either to the left or to the right. $`v(t)`$ thus takes on the value $`a/\mathrm{\Delta }t`$ in a time interval $`\mathrm{\Delta }t`$ when the particle jumps to the right, the value $`a/\mathrm{\Delta }t`$ in a time interval $`\mathrm{\Delta }t`$ when the particle jumps to the left, and zero the rest of the time (It is *not* important whether or not this is a physically reasonable choice for $`v(t)`$. The only requirement for eq. 5.19 to hold is that $`x(t)=_o^tv(t)𝑑t+x(o)`$). With this choice of $`v(t)`$, the function $`v(t^{})v(0)`$ has the value $`a^2/\mathrm{\Delta }t^2`$ if the particle jumps in the same direction at $`t=t^{}`$ and $`t=0`$, the value $`a^2/\mathrm{\Delta }t^2`$ if the particle jumps in opposite directions at $`t=t^{}`$ and $`t=0`$, and zero otherwise (i.e if the particle does not jump *both* at $`t=t^{}`$ and $`t=0`$).
The probability, $`P_{RR}`$, that the particle jumps to the right both at $`t=t^{}`$ and $`t=0`$ is (for now we assume that $`t^{}0`$):
$`P_{RR}`$ $`=`$ $`{\displaystyle \frac{1}{N^d}}{\displaystyle \underset{i,j}{}}\mathrm{\Gamma }_R(j)\mathrm{\Delta }tP(i,t|j+1,0)\mathrm{\Gamma }_R(i)\mathrm{\Delta }t`$ (5.20)
$`=`$ $`{\displaystyle \frac{\mathrm{\Delta }t^2}{N^d}}{\displaystyle \underset{i,j}{}}\mathrm{\Gamma }_L(j)P(i,t|j,0)\mathrm{\Gamma }_R(i)`$ (5.21)
Equation 5.20 can be read from the left as follows; If the particle starts at site $`j`$ the probability of jumping to the right at $`t=0`$ is $`\mathrm{\Gamma }_R(j)\mathrm{\Delta }t`$. This means that it is now at site $`j+1`$, and the probability that it moves to site $`i`$ in the time $`t`$ is $`P(i,t|j+1,0)`$, and finally the probability that it jumps to the right from site $`i`$ is $`\mathrm{\Gamma }_R(i)\mathrm{\Delta }t`$. In equation 5.21 $`j`$ is substituted with $`j1`$, and $`\mathrm{\Gamma }_R(j1)=\mathrm{\Gamma }_L(j)`$ is used. Calculating in the same way the probability of the other events that contributes to the velocity auto correlation function, $`v(t)v(0)_x`$ ($`P_{RL}`$, $`P_{LR}`$, and $`P_{LL}`$), we can now write $`v(t)v(0)_x`$, in terms of the Green’s function (where again use $`a=1`$, and the delta function takes care of the special case $`t=0`$, see below):
$`v(t)v(0)_x`$ $`=`$ $`C\delta (t)+{\displaystyle \frac{1}{\mathrm{\Delta }t^2}}{\displaystyle \underset{i,j}{}}(P_{RR}+P_{LL}P_{RL}P_{LR})`$ (5.22)
$`=`$ $`C\delta (t)+{\displaystyle \frac{1}{N^d}}{\displaystyle \underset{i,j}{}}P(i,t|j,0)`$
$`\left[\mathrm{\Gamma }_L(j)\mathrm{\Gamma }_R(i)+\mathrm{\Gamma }_R(j)\mathrm{\Gamma }_L(i)\mathrm{\Gamma }_L(j)\mathrm{\Gamma }_L(i)\mathrm{\Gamma }_R(j)\mathrm{\Gamma }_R(i)\right]_{}`$
$`=`$ $`C\delta (t)^{}+`$
$`{\displaystyle \frac{1}{N^d}}{\displaystyle \underset{i,j}{}}(\mathrm{\Gamma }_R(i)\mathrm{\Gamma }_L(i))P(i,t|j,0)(\mathrm{\Gamma }_L(j)\mathrm{\Gamma }_R(j))`$
Taking the Laplace transform of equation 5.5 we get:
$$D_x(s)=C+\frac{1}{N^d}\underset{i,j}{}(\mathrm{\Gamma }_R(i)\mathrm{\Gamma }_L(i))G(i,s|j)(\mathrm{\Gamma }_L(j)\mathrm{\Gamma }_R(j))$$
(5.25)
In the high- frequency limit, $`s\mathrm{}`$, we find from equation 5.10
$`G(i,s|j)0`$. We thus get $`C=D_{\mathrm{}}=\mathrm{\Gamma }`$ (see equation 5.5):
$$D_x(s)=\mathrm{\Gamma }\frac{1}{N^d}\underset{i,j}{}(\mathrm{\Gamma }_R(i)\mathrm{\Gamma }_L(i))G(i,s|j)(\mathrm{\Gamma }_R(j)\mathrm{\Gamma }_L(j))$$
(5.26)
In the case where all the jump-rates are identical, $`\mathrm{\Gamma }_R(i)=\mathrm{\Gamma }_L(i)`$, equation 5.26 leads to $`D_x(s)=\mathrm{\Gamma }=D_{\mathrm{}}`$, as expected (see section 5.1).
Equation 5.26 can be used together with equation 5.10 to calculate $`D_x(s)`$. However this method poses a serious numerical problem at low temperatures where $`D_0`$ is very small compared to $`D_{\mathrm{}}`$. Calculating $`D_x(s)`$ at low frequencies from equation 5.26 in this limit amounts to calculating a small difference between two large numbers, which leads to large uncertainties in the result.
In the following we will derive a version of the VAC method, which does not have the numerical problems described above. We do this first in 1 dimension (section 5.5.1), followed by a general derivation in $`𝒟`$ dimensions (section 5.5.2).
#### 5.5.1 The VAC method in 1 Dimension
The idea behind the derivation of the VAC method in 1 dimension is to rewrite the problem in terms of the variable:
$`J_R(i,t|j,0)\mathrm{\Gamma }_R(i)\left[P(i,t|j,0)P(i+1,t|j,0)\right]`$ (5.27)
which can be interpreted as the particle current to the right of site $`i`$ at time $`t`$ given that the particle started at site $`j`$ at $`t=0`$. The master equation (5.7) can then be written as:
$`{\displaystyle \frac{dP(i,t|j,0)}{dt}}=J_R(i1,t|j,0)J_R(i,t|j,0)`$ (5.28)
We now define two new matrices; $`𝚪_R`$ is a diagonal matrix containing the jump rates to the right of each site: $`(𝚪_R)_{ij}=\delta (i,j)\mathrm{\Gamma }_R(i)`$. $`𝐀`$ is a matrix with $`+1`$ on the diagonal, $`1`$ on the sub-diagonal and zero everywhere else: $`(𝐀)_{ij}=\delta (i,j)\delta (i1,j)`$. In more general terms; $`𝐀_{ij}`$ equals $`1`$ if site $`j`$ is the left neighbor to site $`i`$. Here (as in the rest of this chapter) periodic boundary conditions are implicit, i.e. $`(𝐀)_{1N}=1`$. We can now write the 1 dimensional master equation as:
$`{\displaystyle \frac{d}{dt}}𝐏(t)`$ $`=`$ $`\mathrm{𝐀𝐉}_R(t),𝐉_R(t)=𝚪_R𝐀^T𝐏(t)`$ (5.29)
These equations provides us with an an explicit expression for the matrix $`𝐇`$ in the master equation (5.8): $`𝐇=𝐀𝚪_R𝐀^T`$. What we are after here is however an equation for $`𝐉_R`$. To get this we take the Laplace transform of equation 5.29:
$`s𝐆(s)+\mathrm{𝐀𝐉}_R(s)=𝐏(t=0)=𝐈`$ (5.30)
where $`𝐉_R(s)`$ is the Laplace transform of $`𝐉_R(t)`$. By multiplying from the left by $`𝐀^T`$ and using $`𝐉_R(s)=𝚪_R𝐀^T𝐆(s)`$ we get:
$`(s𝚪_R^1+𝐀^T𝐀)𝐉_R(s)=𝐀^T`$ (5.31)
$`𝐉_R(s)=(s𝚪_R^1+𝐀^T𝐀)^1𝐀^T`$ (5.32)
Here we assume that the diagonal matrix $`𝚪_R`$ is invertible, i.e. that all the jump rates are different from zero. We are now ready to derive an equation for D(s) in 1 dimension. To do this we rewrite equation 5.26 using $`\mathrm{\Gamma }_L(i)=\mathrm{\Gamma }_R(i1)`$:
$`D(s)=\mathrm{\Gamma }_R{\displaystyle \frac{1}{N}}{\displaystyle \underset{i,j}{}}(\mathrm{\Gamma }_R(i)\mathrm{\Gamma }_R(i1))G(i,s|j)(\mathrm{\Gamma }_R(j)\mathrm{\Gamma }_R(j1))`$
This can be written as (where $`\mathrm{𝟏}`$ is a column vector containing all 1’s):
$`D(s)N`$ $`=`$ $`\mathrm{𝟏}^T𝚪_R\mathrm{𝟏}\mathrm{𝟏}^T𝚪_R𝐀^T𝐆(s)𝐀𝚪_R\mathrm{𝟏}`$ (5.33)
$`=`$ $`\mathrm{𝟏}^T𝚪_R\mathrm{𝟏}\mathrm{𝟏}^T𝐉_R(s)𝐀𝚪_R\mathrm{𝟏}`$ (5.34)
$`=`$ $`\mathrm{𝟏}^T𝚪_R\mathrm{𝟏}\mathrm{𝟏}^T(s𝚪_R^1+𝐀^T𝐀)^1𝐀^T𝐀𝚪_R\mathrm{𝟏}`$
$`=`$ $`\mathrm{𝟏}^T(s𝚪_R^1+𝐀^T𝐀)^1\left[(s𝚪_R^1+𝐀^T𝐀)𝐀^T𝐀\right]𝚪_R\mathrm{𝟏}`$
$`=`$ $`s\mathrm{𝟏}^T(s𝚪_R^1+𝐀^T𝐀)^1\mathrm{𝟏}`$ (5.36)
$`=`$ $`s\mathrm{𝟏}^T𝐱,(s𝚪_R^1+𝐀^T𝐀)𝐱=\mathrm{𝟏}`$ (5.37)
We have now avoided the problematic subtraction (in equation 5.26 and 5.33), and reduced the problem to finding the vector $`𝐱`$ by solving a sparse system of linear equations.
#### 5.5.2 The VAC method in $`𝒟`$ Dimensions
In this section we generalize the VAC method to $`𝒟`$ dimensions. In analogy with $`𝐉_R(t)`$ used in the 1-dimensional version, we define $`𝐉_k(t)`$ as the particle current in the $`k`$’th direction:
$`𝐉_k(t)=𝚪_k𝐇_k^T𝐏(t)`$ (5.38)
where $`𝚪_k`$ and $`𝐇_k`$ are generalization of the matrices $`𝚪_R`$ and $`𝐀`$ used in the previous section. $`𝚪_k`$ is a diagonal matrix, which for each site contains the jump rate to the “right” in the $`k`$’th direction. $`(𝐇_k)_{ij}`$ equals $`1`$ if site $`j`$ is the “left” neighbor to site $`i`$ along direction $`k`$, and it has $`1`$ on the diagonal and zero everywhere else. The structure of $`𝐇_k`$ depends on the numbering scheme chosen for the sites, but the resulting physics is obviously independent of this. An explicit expression for $`𝐇_k`$ (and the corresponding numbering of the sites) is given in appendix A.
The goal is now to set up a master-equation for the currents, and from that calculate $`D(s)`$, like it was done in the 1-dimensional case. Obviously it is not possible to solve for $`𝐉_k(s)`$ without at the same time solving for the currents in the other directions. We thus generalize equation 5.38 to
$`𝐉_{}(t)=𝚪_{}𝐇_{}^T𝐏(t)𝐉_{}(s)=𝚪_{}𝐇_{}^T𝐆(s)`$ (5.39)
where we have defined the following block-matrices:
$`𝐉_{}(t)\left(\begin{array}{c}𝐉_1(t)\\ 𝐉_2(t)\\ \mathrm{}\\ 𝐉_𝒟(t)\end{array}\right),𝚪_{}`$ $``$ $`\left(\begin{array}{cccc}𝚪_1& \mathrm{𝟎}& & \mathrm{𝟎}\\ \mathrm{𝟎}& 𝚪_2& & \mathrm{𝟎}\\ & & \mathrm{}& \\ \mathrm{𝟎}& \mathrm{𝟎}& & 𝚪_𝒟\end{array}\right),𝐇_{}(𝐇_1,\mathrm{},𝐇_𝒟)`$ (5.48)
and $`𝐉_{}(s)`$ is the Laplace transform of $`𝐉_{}(t)`$.
The master equation for $`𝐏(t)`$ is similar to the one in 1-dimension (eq. 5.29), but the change in the probability at a given site now has contributions from all $`𝒟`$ directions:
$`{\displaystyle \frac{d}{dt}}𝐏(t)={\displaystyle \underset{k=1}{\overset{𝒟}{}}}𝐇_k𝐉_k(t)=𝐇_{}𝐉_{}(t)=𝐇_{}𝚪_{}𝐇_{}^T𝐏(t)`$ (5.50)
$`s𝐆(s)+𝐇_{}𝐉_{}(s)=𝐏(t=0)`$ (5.51)
By multiplying equation 5.51 from the left with $`𝚪_{}𝐇_{}^T`$, and substituting for $`𝐉_{}(s)`$ (equation 5.39), we arrive at the (Laplace transformed) master equation for $`𝐉_{}(s)`$:
$`\left(s𝐈+𝚪_{}𝐇_{}^T𝐇_{}\right)𝐉_{}(s)=𝚪_{}𝐇_{}^T𝐏(t=0)`$ (5.52)
$`𝐉_{}(s)=\left(s𝐈+𝚪_{}𝐇_{}^T𝐇_{}\right)^1𝚪_{}𝐇_{}^T𝐏(t=0)`$ (5.53)
To derive the final result for $`D_k(s)`$ (i.e. the diffusion coefficient as calculated from velocity correlations in direction $`k`$) we define $`\mathrm{𝟏}_k`$ as a column vector with $`𝒟N^𝒟`$ elements, where the $`N^𝒟`$ elements corresponding to currents in the $`k`$’th direction is $`1`$ and the rest is zero (so that eg. $`\mathrm{𝟏}_k^T𝚪_{}\mathrm{𝟏}_k=\mathrm{𝟏}^T𝚪_k\mathrm{𝟏}`$):
$`\mathrm{𝟏}_k\left(\begin{array}{c}\mathrm{𝟎}\\ \mathrm{}\\ \mathrm{𝟏}\\ \mathrm{}\\ \mathrm{𝟎}\end{array}\right)`$ (5.59)
From equation 5.26 we finally get (using $`𝐏(t=0)=𝐈`$):
$`D_k(s)N^𝒟`$ $`=`$ $`\mathrm{𝟏}^T\left(𝚪_k𝐉_k(s)𝐇_k𝚪_k\right)\mathrm{𝟏}`$ (5.60)
$`=`$ $`\mathrm{𝟏}_k^T\left(𝐈𝐉_{}(s)𝐇_{}\right)𝚪_{}\mathrm{𝟏}_k`$ (5.61)
$`=`$ $`\mathrm{𝟏}_k^T\left(𝐈\left(s𝐈+𝚪_{}𝐇_{}^T𝐇_{}\right)^1𝚪_{}𝐇_{}^T𝐇_{}\right)𝚪_{}\mathrm{𝟏}_k`$
$`=`$ $`\mathrm{𝟏}_k^T\left(s𝐈+𝚪_{}𝐇_{}^T𝐇_{}\right)^1\left(\left(s𝐈+𝚪_{}𝐇_{}^T𝐇_{}\right)𝚪_{}𝐇_{}^T𝐇_{}\right)𝚪_{}\mathrm{𝟏}_k`$
$`=`$ $`s\mathrm{𝟏}_k^T\left(s𝐈+𝚪_{}𝐇_{}^T𝐇_{}\right)^1𝚪_{}\mathrm{𝟏}_k`$ (5.63)
$`=`$ $`s\mathrm{𝟏}_k^T\left(s𝚪_{}^1+𝐇_{}^T𝐇_{}\right)^1\mathrm{𝟏}_k`$ (5.64)
$`=`$ $`s\mathrm{𝟏}_k^T𝐱,\left(s𝚪_{}^1+𝐇_{}^T𝐇_{}\right)𝐱=\mathrm{𝟏}_k`$ (5.65)
Like in the 1-dimensional case in the previous section, we have now reduced the problem to finding the vector $`𝐱`$ by solving a sparse system of linear equations. In $`𝒟`$ dimensions the matrix involved in this is a ($`𝒟N^𝒟`$x$`𝒟N^𝒟`$) matrix.
For real $`s>0`$ the matrix $`\left(s𝚪_{}^1+𝐇_{}^T𝐇_{}\right)`$ is positive definite<sup>2</sup><sup>2</sup>2The matrix $`𝐀`$ is symmetric and positive definite if $`𝐲^T\mathrm{𝐀𝐲}>0`$ for any $`𝐲`$., which ensures the numerical stability when doing Gaussian elimination :
$`𝐲^T\left(s𝚪_{}^1+𝐇_{}^T𝐇_{}\right)𝐲`$ $`=`$ $`s𝐲^T𝚪_{}^1𝐲+𝐲^T𝐇_{}^T𝐇_{}𝐲`$ (5.66)
$`=`$ $`s𝐲^T𝚪_{}^1𝐲+(𝐇_{}𝐲)^T(𝐇_{}𝐲)`$ (5.67)
The second term is greater than or equal to zero, and the first term is greater than zero, since $`𝚪_{}^1`$ is a diagonal matrix with only positive values on the diagonal. When representing the problem in finite precision in the computer this last statement is violated for very low $`s`$ values; The elements corresponding to small energy barriers (large jump-rates) are effectively set to zero when the two terms are added (all elements in $`𝐇_{}`$ are O(1)). This means that the matrix in practice is *not* positive definite and the attempt to solve it fails. To avoid this numerical problem a small constant $`\delta `$ ($`=10^{14}`$) is added to the diagonal ensuring that the matrix is positive definite and thus can be solved. In physical terms this corresponds to applying a minimum value for the energy barriers, ensuring that all jump-rates are finite ($`<s/\delta `$). It was checked numerically (by varying $`\delta `$) that the effect of this procedure is negligible (relative error $`10^4`$), as expected from our physical understanding of the model (see section 5.3.2).
Like in paper I the jump rates corresponding to energy barriers larger than $`E_c+K/\beta `$ ($`K=6.4`$) was set to zero, to speed up the calculations (and reduce the amount of memory needed). By varying $`K`$ the relative error introduced by this procedure was estimated to be less than 1%.
Equation 5.65 was solved using Cholesky factorization for real $`s`$-values and LU factorization for imaginary $`s`$-values. In both cases pivoting was done using the minimum degree algorithm, and all computations was done using Matlab version 5.3.0.
### 5.6 Numerical results, $`D(s)`$
In this section we report the results of numeric calculations of $`D(s)`$ in 3 dimensions using the VAC method. Like in and paper I the calculations are done with for real Laplace frequencies, $`s`$, i.e. corresponding to imaginary values of $`\omega `$ ($`s=i\omega `$). When comparing the numeric results and analytical approximations for real values of $`s`$, we are relying on the principle of analytical continuation ; If two complex functions coincide on a line in the complex plane, they coincide everywhere in the complex plane (where they are both well-defined).
In section 5.7 we present the corresponding results for real values of $`\omega `$. Besides providing results directly comparable with experimental results for $`\sigma (\omega )`$, it will be clearly demonstrated that we can indeed trust the principle of analytical continuation when comparing analytical approximations and numerical results.
It has already been demonstrated in and paper I that the symmetric hopping model becomes universal in the extreme disorder limit, i.e. that $`\stackrel{~}{D}(\stackrel{~}{s})`$ (or equivalently $`\stackrel{~}{\sigma }(\stackrel{~}{s}))`$ becomes independent of temperature and the distribution of energy-barriers. Here we will focus on the shape of $`\stackrel{~}{D}(\stackrel{~}{s})`$, and only present results for the box-distribution of energy barriers: $`p(E)=1,0E1`$.
Figure 5.4 show the frequency dependent diffusion coefficient for real Laplace frequencies, $`D(s)`$, in a log-log plot, for 4 $`\beta `$-values. Both axis in figure 5.4 are scaled by the DC level, $`D_0D(s0)`$. In the inset of figure 5.4 is shown $`D_0`$ scaled by $`\mathrm{\Gamma }(E_c)`$ vs. $`\beta `$. For $`\beta 80`$ $`D_0`$ is seen to be well approximated by:
$$D_0\beta ^\gamma \mathrm{\Gamma }(E_c),\gamma =0.89\pm 0.01$$
(5.68)
Note that the dominant $`\beta `$-dependence of $`D_0`$ is given by $`\mathrm{\Gamma }(E_c)`$, which changes by 30 orders of magnitude for the $`\beta `$-values used (see table 5.1). In the simulations reported in paper I , we found: $`\gamma =0.81\pm 0.04`$ .
To illustrate how $`D(s)`$ approaches a universal curve (suitable scaled) at high $`\beta `$-values, the data presented in figure 5.4 was scaled by the characteristic frequency, $`s_c`$, defined here by $`D(s_c)/D_0=\sqrt{10}`$. In figure 5.5 the result of scaling the data in this way is shown. It is clearly seen that $`\stackrel{~}{D}(\stackrel{~}{s})`$ approaches a universal curve, as the $`\beta `$-values increases. Or in other words: the way the system approaches (long time) diffusion becomes universal as the temperature is decreased. The universal curve is estimated to be close to the data for $`\beta =320`$, and the frequency regime for which the data follows the universal curve is seen to increase as $`\beta `$ increases.
In the inset of figure 5.5 is shown $`s_c`$ scaled by $`D_0`$ vs. $`\beta `$. For $`\beta 80`$ $`s_c`$ is seen to be well approximated by:
$$s_c\beta ^\eta D_0,\gamma =1.37\pm 0.01$$
(5.69)
The dominant $`\beta `$-dependence of $`s_c`$ is given by $`D_0`$, which changes by 31 orders of magnitude for the $`\beta `$-values used (see table 5.1). In the simulations reported in paper I, the power law was less well defined, and $`\eta `$ was estimated to be $`1.48\pm 0.06`$ (depending on the $`\beta `$-range used in the power-law fit) .
For both of the scaling parameters, $`D_0`$ and $`s_c`$, the $`\beta `$-dependence given by the power-laws reported are very small compared to the $`\beta `$-dependence given by the other involved quantities. In the simulations presented here, the main focus has not been on determining the scaling exponents, $`\gamma `$ and $`\eta `$, precisely, but on the *shape* of the universal curve seen in figure 5.5.
Before proceeding to compare the numeric results with the analytical approximations, we will briefly discuss the possibilities of computing $`\mathrm{\Delta }X^2(t)`$ from $`D(s)`$. As discussed above we will compare numerical results and analytical approximations for $`D(s)`$, and consequently we are calculating $`\mathrm{\Delta }X^2(t)`$
only for illustrative purposes. From equation 5.3 this is in principle straight forward; $`\mathrm{\Delta }X^2(t)`$ is the inverse Laplace transform of $`2D(s)/s^2`$. The general inverse Laplace transform is however in practice problematic , and here we furthermore have the problem of the many decades of frequency/time involved. Instead we use the following trick; we fit $`D(s)`$ in such a way that $`2D(s)/s^2`$ has a known inverse Laplace transform. This is done for $`\beta =320`$ in figure 5.6a, where $`D(s)`$ is fitted to a sum of power-laws. Note that this is a purely empirical fit, and it does not necessarily have anything to do with the true functional form of $`D(s)`$. The inset in figure 5.6a shows $`2D(s)/s^2`$ and the corresponding fit, i.e. the quantity to which the inverse Laplace transform is applied. The resulting $`\mathrm{\Delta }X^2(t)`$ is shown in figure 5.6b, without using the scaling parameters, so that the actual size of the quantities is seen ($`t_c1/s_c`$). Note that the system becomes diffusive at very long time scales; $`t10^{40}`$. This illustrates the impossibility of using traditional MC techniques, since this would require at least $`10^{40}`$ time steps. The inset in figure 5.6b shows $`D(t)`$ and $`\mathrm{\Delta }X^2(t)/2t`$ as resulting from the procedure described above.
In figure 5.7 the universal curve for $`D(s)`$ (represented by the data for $`\beta =320`$) is compared with the 3 analytical approximations described in section 5.3. The fractal dimension $`D_f`$ used in the Diffusion Cluster Approximation (DCA) was treated as a fitting parameter. The fitting procedure will be explained when discussing figure 5.9. The fit of DCA is seen to be almost perfect, and clearly better than both PPA and EMA. This is perhaps not surprising given that DCA has a fitting parameter which PPA and EMA has not, and the value of fit achieved by DCA thus greatly depends on whether an independent argument can be given for the value of $`D_f`$. For now we only note, that the value $`D_f=1.35`$ is in the expected range: $`1.14<D_f<1.7`$ (see section 5.3).
Like in paper I, we now proceed to focus on the shape of $`\stackrel{~}{D}(\stackrel{~}{s})`$, by calculating the apparent power-law exponent, i.e. the logarithmic derivative of $`\stackrel{~}{D}(\stackrel{~}{s})`$:
$$\mu \frac{d\mathrm{log}_{10}(D)}{d\mathrm{log}_{10}(s)}$$
(5.70)
When plotting $`\mu `$ as a function of $`\mathrm{log}_{10}(D(s)/D_0)`$ we can compare the shapes of the data presented above, without relying on an empirical scaling of the frequency axis. This is done in figure 5.8, for the same data as shown in 5.5. Plotting the data in this way is seen to be more sensitive; The small difference between the data for $`\beta =160`$, and $`\beta =320`$ in figure 5.5 is more pronounced in figure 5.8.
In figure 5.9 we compare the universal curve for $`\mu `$ (approximated by the data for $`\beta =320`$) with the 3 analytical approximations. As in figure 5.7 the data from the simulations is seen to lie between EMA and PPA, and DCA is seen to give a good fit. The fitting of DCA was done by hand using this figure. The inset in figure 5.9 shows the same data, but focusing on the low values of $`\stackrel{~}{D}`$, corresponding to low values of $`\stackrel{~}{s}`$ and $`\mu `$. In this limit, the data is seen to deviate from DCA, and seems to approach EMA.
The behavior of $`\stackrel{~}{D}(\stackrel{~}{s})`$ at low frequencies, i.e. the departure from the DC level ($`D_0`$) is better investigated by plotting $`\stackrel{~}{D}(\stackrel{~}{s})D_0`$ as its done in figure 5.10. The universality for the numerical data is seen to hold even at the very low frequencies in this plot. This demonstrates that the universality seen at low frequencies in $`\stackrel{~}{D}(\stackrel{~}{s})`$ (figure 5.5), is *not* just a consequence of the DC level being dominant. On the other hand the apparent reasonable agreement between $`\stackrel{~}{D}(\stackrel{~}{s})`$ and DCA seen in figure 5.7 breaks down when the DC level is subtracted.
The low-frequency expansion of $`D(s)`$ is known to be :
$`\stackrel{~}{D}(\stackrel{~}{s})=1+A\stackrel{~}{s}+B\stackrel{~}{s}^{3/2}+\mathrm{}`$ (5.71)
In figure 5.10 the numerical data and EMA is seen to agree with the first two terms of this expansion.
### 5.7 Numerical Results, $`D(\omega )`$
In this section we present numerical results for $`D(\omega )`$, i.e. for imaginary Laplace-frequency, $`s=i\omega `$. In this case the diffusion coefficient is a complex quantity, and we can compare both the real and imaginary part with the analytical approximations.
In figure 5.11 is shown the (scaled) real part of $`D(\omega )`$ vs. the frequency (scaled). The scaling parameters, $`D_0`$ and $`\omega _c`$ is shown in figure 5.13. Like we found in the previous section for real Laplace frequencies the data is here seen to approach a universal curve which agrees well with DCA, with $`D_f=1.35`$. This is the value found in the previous section from $`D(s)`$; it was not found to be necessary to make a new fit for $`D(\omega )`$.
Figure 5.12 shows the imaginary part of $`D(\omega )`$ corresponding to the data in figure 5.11, and using the same scaling parameters (and $`D_f`$). Universality is evident at low frequencies (agreeing with EMA, apart from scaling), and is approached at higher frequencies.
As an analogue to the (apparent) exponent $`\mu `$ (eq. 5.70) used for $`D(s)`$ in the previous section, we can define an exponent for the real part of $`D(\omega )`$:
$$\mu _{real}\frac{d\mathrm{log}_{10}(D^{}(\omega ))}{d\mathrm{log}_{10}(\omega )}$$
(5.72)
$`\mu _{real}`$ is plotted in figure 5.14 as a function of $`\omega `$. Note that the convergence toward universality is more “abrupt” than it was found for $`D(s)`$ (fig. 5.8); Only the data for $`\beta =40`$ deviates significantly from DCA.
In figure 5.15 we focus on the low frequency behavior of $`D^{}(\omega )`$ by subtracting the DC level, which gives us the chance to check the agreement with third term in the low-frequency expansion (Eq. 5.71), i.e. the exponent $`3/2`$. Neither of the approximations agrees with this, which in particular means that EMA only agrees with the first 2 terms in the low-frequency expansion. The numerical data seems to agree better with the exponent $`3/2`$, although it is difficult to judge the numerical significance of this.
We define the apparent exponent for the imaginary part of $`D(\omega )`$ in the following way:
$$\mu _{imag}\frac{d\mathrm{log}_{10}(D^{\prime \prime }(\omega ))}{d\mathrm{log}_{10}(\omega )}$$
(5.73)
This is plotted as a function of $`D^{\prime \prime }(\omega )`$ in figure 5.16. EMA has the right value at the very low frequencies, as seen in fig. 5.12. DCA has the right behavior from $`D^{\prime \prime }(\omega )D_0`$, with small deviations at the highest values (notice the y-axis being different from the other exponent-plots).
In the data presented so far we have found much the same behavior in $`D(\omega )`$ as for $`D(s)`$ with regards to how the numerical data agrees with the analytical approximations; at high frequencies the numerical data falls between EMA and PPA and is well fitted by DCA with $`D_f=1.35`$, whereas the very low frequencies is governed by the low frequency expansion (eq. 5.71), which (for the first 2 terms) agrees with EMA. In contrast the approach to universality seems to different for especially $`D(s)`$ and $`D^{}(\omega )`$. The quantity that approaches universality in a similar manner as $`D(s)`$ is the absolute value of $`D(\omega )`$. This is illustrated in figure 5.17 where we plot the apparent exponent for the absolute value of $`D(\omega )`$:
$$\mu _{abs}\frac{d\mathrm{log}_{10}(|D(\omega )|)}{d\mathrm{log}_{10}(\omega )}$$
(5.74)
### 5.8 VAC vs. ACMA
In figure 5.18 the universal curve for $`D(s)`$ as calculated from the VAC method (fig. 5.7) and the ACMA method (fig. 1 in paper I) is compared. There is a small but significant difference between the results from the two methods, as can also be seen by comparing the apparent exponents from the two methods (fig. 5.8 and fig. 2 in paper I). The main difference between the two methods is the boundary conditions, so the differences in these are the “main suspect” for the (slightly) different results. At infinitely large samples we would expect the results of both methods to converge to the bulk-limit. The way the two methods converge to the bulk limit might however be different; In the ACMA method there is some fraction of the sites that are close to an electrode, and thus give a wrong contribution to $`D(s)`$. As the sample size increases this fraction decreases, and the results converges to the true bulk-limit. In the VAC method there are no sites that are “worse” than the other, and the finite size effects are more subtle; they arise from particles traveling trough the sample to where they started, and thus experiencing false correlations.
In figure 5.19 we compare the apparent exponent $`\mu `$ for $`\beta =320`$ as calculated with $`N=96`$ (i.e. the data in fig. 5.8) and $`N=64`$ (averaged over 600 samples). The agreement is seen to be excellent, with the largest error at the very low frequencies. The inset in figure 5.19 shows the N-dependence of $`D_0`$, with error-bars estimated from the standard deviation. The tendency to converge to a constant as $`N`$ increases is evident.
In figure 5.19 is shown $`\mu `$ calculated independently for each of the 100 (96x96x96) samples at $`\beta =320`$, i.e. without averaging. The noise shows that the samples are not self-averaging, i.e. it *is* necesarry to average over a number of samples. The data points are distributed evenly around DCA, showing that the averaging over (non self-averaging) samples do not introduce systematic errors (compare fig. 5.9)
### 5.9 Conclusions
The main result in this chapter is the development of the Velocity Auto Correlation (VAC) method. At first it might seem strange to work in terms of velocities in a hopping model, as it is evident from the following quote from Scher & Lax :
> The relation as it stands \[ i.e. eq. 5.19 in the present work, relating $`D(s)`$ to the velocity auto correlation function \] is inconvenient to use in a hopping model since it refers to velocities rather than positions.
It should be evident by now, that it is worth to suffer the (initial) inconvenience of working with velocities; The VAC method is clearly to prefer to the ACMA method, since it can be used with periodic boundary conditions, and still give a diffusive regime. The physical reason for this being possible is that the diffusive regime is characterised by loss of correlations, which is possible in a finite sample (as opposed to the mean square displacement being proportional to time).
The numerical results was found to be in excellent agreement with the Diffusion Cluster Approximation, except for the very low frequencies where the low frequency expansion holds. The agreement with DCA was achieved by a (single) fit to the fractal dimension of the diffusion cluster, $`D_f=1.35`$, which is within the limits expected from percolation arguments: $`1.14<D_f<1.7`$ (see section 5.3). Two questions needs to be answered, to decide whether the agreement between the numerical data and DCA signals that the picture behind the approximation is correct, or if it is simply a result of fitting:
* Is DCA a good approximation of diffusion on a diffusion cluster with fractal dimension $`D_f`$? This corresponds to the check of PPA in 1 dimension shown in figure 2a in paper I.
* Is the diffusion cluster (in 3 dimensions) a fractal with fractal dimension (close to) 1.35?
If either of the answers to these questions are negative, then the agreement found between numerical data and DCA is a coincidence, and DCA is merely a convenient way to describe the data. It is left as future work to answer these questions.
## Appendix A Numbering scheme for VAC method
The sites in the $`𝒟`$-dimensional regular lattice are numbered by:
$`\text{SiteIndex}={\displaystyle \underset{j=1}{\overset{𝒟}{}}}(C_j1)N^{j1}+1`$ (A.1)
where $`C_j`$ is the coordinate (counted from 1 to N) of the site in the $`j`$’th direction.
The ($`N^𝒟\times N^𝒟`$) matrix $`𝐇_k`$ used in section 5.5.2 (eq. 5.38) describes the connectiviy in the lattice; $`(𝐇_k)_{ij}`$ equals $`1`$ if site $`j`$ is the “left” neighbor to site $`i`$ along direction $`k`$, and it has $`1`$ on the diagonal and zero everywhere else. With the numbering scheme given above, $`𝐇_k`$ is given by $`𝐈𝐈\mathrm{}𝐀\mathrm{}𝐈𝐈`$, where $`𝐈`$ is the ($`N\times N`$) identity matrix, $``$ is the direct (Kronecker) multiplication, and the matrix $`𝐀`$ is at the $`j`$’th position from the right<sup>1</sup><sup>1</sup>1 The ($`N\times N`$) matrix $`𝐀`$ is here defined as in section 5.5.1: $`(𝐀)_{ij}=\delta (i,j)\delta (i1,j)`$ (with periodic boundary conditions).
In 2 dimensions with N=3 we eg. have:
$`𝐇_1`$ $``$ $`𝐈𝐀=\left(\begin{array}{ccc}𝐀& \mathrm{𝟎}& \mathrm{𝟎}\\ \mathrm{𝟎}& 𝐀& \mathrm{𝟎}\\ \mathrm{𝟎}& \mathrm{𝟎}& 𝐀\end{array}\right)`$ (A.5)
$`=`$ $`\left(\begin{array}{ccccccccc}\hfill 1& \hfill 0& \hfill 1& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0\\ \hfill 1& \hfill 1& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0\\ \hfill 0& \hfill 1& \hfill 1& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0\\ \hfill 0& \hfill 0& \hfill 0& \hfill 1& \hfill 0& \hfill 1& \hfill 0& \hfill 0& \hfill 0\\ \hfill 0& \hfill 0& \hfill 0& \hfill 1& \hfill 1& \hfill 0& \hfill 0& \hfill 0& \hfill 0\\ \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 1& \hfill 1& \hfill 0& \hfill 0& \hfill 0\\ \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 1& \hfill 0& \hfill 1\\ \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 1& \hfill 1& \hfill 0\\ \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 1& \hfill 1\end{array}\right)`$ (A.15)
$`𝐇_2`$ $``$ $`𝐀𝐈=\left(\begin{array}{ccc}\hfill 𝐈& \hfill \mathrm{𝟎}& \hfill 𝐈\\ \hfill 𝐈& \hfill 𝐈& \hfill \mathrm{𝟎}\\ \hfill \mathrm{𝟎}& \hfill 𝐈& \hfill 𝐈\end{array}\right)`$ (A.19)
$`=`$ $`\left(\begin{array}{ccccccccc}\hfill 1& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 1& \hfill 0& \hfill 0\\ \hfill 0& \hfill 1& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 1& \hfill 0\\ \hfill 0& \hfill 0& \hfill 1& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 1\\ \hfill 1& \hfill 0& \hfill 0& \hfill 1& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0\\ \hfill 0& \hfill 1& \hfill 0& \hfill 0& \hfill 1& \hfill 0& \hfill 0& \hfill 0& \hfill 0\\ \hfill 0& \hfill 0& \hfill 1& \hfill 0& \hfill 0& \hfill 1& \hfill 0& \hfill 0& \hfill 0\\ \hfill 0& \hfill 0& \hfill 0& \hfill 1& \hfill 0& \hfill 0& \hfill 1& \hfill 0& \hfill 0\\ \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 1& \hfill 0& \hfill 0& \hfill 1& \hfill 0\\ \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 1& \hfill 0& \hfill 0& \hfill 1\end{array}\right)`$ (A.29)
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# MZ-TH/00-22 Some remarks on the 𝜀-expansion of dimensionally regulated Feynman diagramsBased on the talk given by A. D. at the Zeuthen Workshop “Loops and Legs in Gauge Theories” (Bastei, Germany, April 2000).
## Abstract
Some problems related to construction of the $`\epsilon `$-expansion of dimensionally regulated Feynman integrals are discussed. For certain classes of diagrams, an arbitrary term of the $`\epsilon `$-expansion can be expressed in terms of log-sine integrals related to the polylogarithms. It is shown how the analytic continuation of these functions can be constructed in terms of the generalized Nielsen polylogarithms.
1. Dimensional regularization is one of the most powerful tools used in loop calculations. In some cases, one can derive results valid for an arbitrary space-time dimension $`n`$, usually in terms of various hypergeometric functions. However, for practical purposes the coefficients of the expansion in $`\epsilon `$ are important, where the regulator $`\epsilon `$ corresponds to the difference between $`n`$ and the (integer) number of dimensions of interest. Below we shall usually imply that $`n=42\epsilon `$. In multi-loop calculations higher terms of the $`\epsilon `$-expansion of one- and two-loop functions are needed, since they may get multiplied by poles in $`\epsilon `$, not only due to factorizable loops, but also as a result of applying the well-known reduction techniques .
In refs. , it was shown that the log-sine integral functions (see, e.g., in , chapter 7.9),
$$\text{Ls}_j(\theta )\underset{0}{\overset{\theta }{}}\text{d}\theta ^{}\mathrm{ln}^{j1}\left|2\mathrm{sin}\frac{\theta ^{}}{2}\right|,$$
(1)
happen to be very useful to represent results for higher terms of the $`\epsilon `$-expansion.
For instance, for the one-loop two-point function $`J^{(2)}(n;\nu _1,\nu _2)`$ with external momentum $`k`$, masses $`m_1`$ and $`m_2`$ and unit powers of propagators, the following result for an arbitrary term of the $`\epsilon `$-expansion has been obtained in :
$`J^{(2)}(42\epsilon ;1,1)=\text{i}\pi ^{2\epsilon }{\displaystyle \frac{\mathrm{\Gamma }(1+\epsilon )}{2(12\epsilon )}}`$
$`\times \{{\displaystyle \frac{m_1^{2\epsilon }+m_2^{2\epsilon }}{\epsilon }}+{\displaystyle \frac{m_1^2m_2^2}{\epsilon k^2}}(m_1^{2\epsilon }m_2^{2\epsilon })`$
$`+{\displaystyle \frac{\left[\mathrm{\Delta }(m_1^2,m_2^2,k^2)\right]^{1/2\epsilon }}{(k^2)^{1\epsilon }}}{\displaystyle \underset{j=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(2\epsilon )^j}{j!}}`$
$`\times {\displaystyle \underset{i=1}{\overset{2}{}}}[\text{Ls}_{j+1}(\pi )\text{Ls}_{j+1}(2\tau _{0i}^{})]\},`$ (2)
where
$`\mathrm{cos}\tau _{01}^{}=(m_1^2m_2^2+k^2)/(2m_1\sqrt{k^2}),`$
$`\mathrm{cos}\tau _{02}^{}=(m_2^2m_1^2+k^2)/(2m_2\sqrt{k^2}),`$
$`\mathrm{cos}\tau _{12}=(m_1^2+m_2^2k^2)/(2m_1m_2),`$ (3)
whereas the “triangle” function $`\mathrm{\Delta }`$ is defined as
$$\mathrm{\Delta }(x,y,z)=2xy+2yz+2zxx^2y^2z^2.$$
(4)
One can see that $`\tau _{12}+\tau _{01}^{}+\tau _{02}^{}=\pi `$. In fact, these angles can be associated with a triangle whose sides are $`m_1`$, $`m_2`$ and $`\sqrt{k^2}`$. Moreover, the area of this triangle is $`\frac{1}{4}\sqrt{\mathrm{\Delta }(m_1^2,m_2^2,k_{12}^2)}`$. For details of geometrical description, see in .
Note that the values of $`\text{Ls}_j(\pi )`$ can be expressed in terms of Riemann’s $`\zeta `$ function, see Eqs. (7.112)–(7.113) of . The infinite sum with $`\text{Ls}_j(\pi )`$ in (MZ-TH/00-22 Some remarks on the $`\epsilon `$-expansion of dimensionally regulated Feynman diagramsthanks: Based on the talk given by A. D. at the Zeuthen Workshop “Loops and Legs in Gauge Theories” (Bastei, Germany, April 2000).) can be converted into $`\mathrm{\Gamma }`$ functions,
$$\underset{j=0}{\overset{\mathrm{}}{}}\frac{(2\epsilon )^j}{j!}\text{Ls}_{j+1}(\pi )=\pi \frac{\mathrm{\Gamma }(1+2\epsilon )}{\mathrm{\Gamma }^2(1+\epsilon )}.$$
(5)
2. The $`\epsilon `$-expansion (MZ-TH/00-22 Some remarks on the $`\epsilon `$-expansion of dimensionally regulated Feynman diagramsthanks: Based on the talk given by A. D. at the Zeuthen Workshop “Loops and Legs in Gauge Theories” (Bastei, Germany, April 2000).) is directly applicable in the region where $`\mathrm{\Delta }(m_1^2,m_2^2,k^2)0`$, i.e. when $`(m_1m_2)^2k^2(m_1+m_2)^2`$. In other regions, the proper analytic continuation of the occurring $`\text{Ls}_j\left(\theta \right)`$ should be constructed. To do this, it is convenient to introduce the variable (cf. in ref. )
$$ze^{\mathrm{i}\sigma \theta },\mathrm{ln}(z)=\mathrm{ln}(z)\mathrm{i}\sigma \pi ,$$
(6)
where the choice of the sign $`\sigma =\pm 1`$ is related to the causal “+i0” prescription for the propagators.
Since $`\text{Ls}_1(\theta )=\theta `$, we get
$$\mathrm{i}\sigma \left[\text{Ls}_1\left(\pi \right)\text{Ls}_1\left(\theta \right)\right]=\mathrm{ln}(z).$$
(7)
For the next order, we can use the fact that $`\text{Ls}_2(\theta )=\text{Cl}_2(\theta )`$, where (see in )
$`\text{Cl}_j(\theta )=\{\begin{array}{c}\frac{1}{2\mathrm{i}}\left[\text{Li}_j\left(e^{\mathrm{i}\theta }\right)\text{Li}_j\left(e^{\mathrm{i}\theta }\right)\right],j\text{even}\hfill \\ \frac{1}{2}\left[\text{Li}_j\left(e^{\mathrm{i}\theta }\right)+\text{Li}_j\left(e^{\mathrm{i}\theta }\right)\right],j\text{odd}\hfill \end{array}`$ (10)
is the Clausen function, whereas $`\text{Li}_j`$ is the polylogarithm. In other words, $`\text{Cl}_j(\theta )`$ corresponds either to the imaginary part or to the real part of $`\text{Li}_j\left(e^{\mathrm{i}\theta }\right)`$, depending on whether $`j`$ is even or odd. Therefore, the analytic continuation reads
$$\mathrm{i}\sigma \left[\text{Ls}_2\left(\pi \right)\text{Ls}_2\left(\theta \right)\right]=\frac{1}{2}\left[\text{Li}_2\left(z\right)\text{Li}_2\left(1/z\right)\right],$$
(11)
where $`\text{Ls}_2\left(\pi \right)=0`$. The result for the $`\epsilon `$-term of the two-point function was obtained in .
To proceed further, we need similar relations between higher $`\text{Ls}_j\left(\theta \right)`$ and the imaginary (or real) parts of the polylogarithms. For $`j=3`$, $`\text{Ls}_3(\theta )`$ can be expressed in terms of the imaginary part of $`\text{Li}_3\left(1e^{\mathrm{i}\theta }\right)`$. Then, the imaginary part of $`\text{Li}_4\left(1e^{\mathrm{i}\theta }\right)`$ is already a mixture of $`\text{Ls}_4(\theta )`$ and $`\text{Cl}_4(\theta )`$, whereas its real part involves the generalized log-sine integral $`\text{Ls}_4^{(1)}(\theta )`$. All these relations can be found in (some misprints are mentioned in ). However, attempts to generalize these results to higher functions show that the relations get more and more cumbersome.
Instead of going that way, we suggest to consider how the higher $`\text{Ls}_j`$ functions are generated by the imaginary (and real) parts of the generalized Nielsen polylogarithms (see, e.g., in ),
$$S_{a,b}(z)=\frac{(1)^{a+b1}}{(a1)!b!}\underset{0}{\overset{1}{}}\text{d}\xi \frac{\mathrm{ln}^{a1}\xi \mathrm{ln}^b(1z\xi )}{\xi },$$
(12)
where $`S_{a,1}(z)=\text{Li}_{a+1}\left(z\right).`$ We obtain
$`\text{Re}S_{1,2}(e^{\mathrm{i}\theta })=\frac{1}{2}\text{Cl}_3\left(\theta \right)+\frac{1}{2}\zeta _3\frac{1}{2}\left(\pi \theta \right)\text{Ls}_2\left(\theta \right),`$
$`\text{Im}S_{1,2}(e^{\mathrm{i}\theta })=\frac{1}{2}\text{Ls}_3\left(\theta \right)\frac{1}{24}\theta \left(\theta ^23\pi \theta +3\pi ^2\right),`$
$`\text{Re}S_{1,3}(e^{\mathrm{i}\theta })=\frac{1}{4}\text{Ls}_4^{(1)}\left(\theta \right)+\frac{1}{4}\pi \text{Ls}_3\left(\theta \right)`$
$`+\frac{1}{90}\pi ^4+\frac{1}{48}\pi ^3\theta \frac{1}{32}\pi ^2\theta ^2+\frac{1}{48}\pi \theta ^3\frac{1}{192}\theta ^4,`$
$`\text{Im}S_{1,3}(e^{\mathrm{i}\theta })=\frac{1}{6}\text{Ls}_4\left(\theta \right)+\frac{1}{4}\text{Cl}_4\left(\theta \right)\frac{1}{4}\pi \zeta _3`$
$`+\frac{1}{4}\left(\pi \theta \right)\text{Cl}_3\left(\theta \right)\frac{1}{8}\left(\pi \theta \right)^2\text{Ls}_2\left(\theta \right),`$
where the generalized log-sine integral (see in ) is defined as
$$\text{Ls}_j^{(k)}\left(\theta \right)=\underset{0}{\overset{\theta }{}}\text{d}\theta ^{}\theta ^k\mathrm{ln}^{jk1}\left|2\mathrm{sin}\frac{\theta ^{}}{2}\right|.$$
(13)
In particular, $`\text{Ls}_j^{(0)}\left(\theta \right)=\text{Ls}_j\left(\theta \right)`$. Using these relations, we can express the $`\text{Ls}_3`$ and $`\text{Ls}_4`$ functions,
$`\mathrm{i}\sigma \left[\text{Ls}_3\left(\pi \right)\text{Ls}_3\left(\theta \right)\right]=`$
$`S_{1,2}(z)S_{1,2}(1/z)\frac{1}{12}\mathrm{ln}^3(z),`$ (14)
$`\mathrm{i}\sigma \left[\text{Ls}_4\left(\pi \right)\text{Ls}_4\left(\theta \right)\right]=3\left[S_{1,3}(z)S_{1,3}(1/z)\right]`$
$`+\frac{3}{4}\left[\text{Li}_4\left(z\right)\text{Li}_4\left(1/z\right)\right]`$
$`\frac{3}{4}\left[\text{Li}_3\left(z\right)+\text{Li}_3\left(1/z\right)\right]\mathrm{ln}(z)`$
$`+\frac{3}{8}\left[\text{Li}_2\left(z\right)\text{Li}_2\left(1/z\right)\right]\mathrm{ln}^2(z),`$ (15)
where $`\text{Ls}_3\left(\pi \right)=\frac{1}{12}\pi ^3`$, $`\text{Ls}_4\left(\pi \right)=\frac{3}{2}\pi \zeta _3`$.
We have also constructed further expressions,
$`\mathrm{i}\sigma \left[\text{Ls}_5\left(\pi \right)\text{Ls}_5\left(\theta \right)\right]=12\left[S_{1,4}(z)S_{1,4}(1/z)\right]`$
$`3\left[S_{2,3}(z)S_{2,3}(1/z)\right]`$
$`+3\left[S_{1,3}(z)+S_{1,3}(1/z)\right]\mathrm{ln}(z)`$
$`+\frac{1}{80}\mathrm{ln}^5(z),`$ (16)
$`\mathrm{i}\sigma \left[\text{Ls}_6\left(\pi \right)\text{Ls}_6\left(\theta \right)\right]=60\left[S_{1,5}(z)S_{1,5}(1/z)\right]`$
$`+15\left[S_{2,4}(z)S_{2,4}(1/z)\right]`$
$`\frac{15}{4}\left[\text{Li}_6\left(z\right)\text{Li}_6\left(1/z\right)\right]`$
$`15\left[S_{1,4}(z)+S_{1,4}(1/z)\right]\mathrm{ln}(z)`$
$`+\frac{15}{4}\left[\text{Li}_5\left(z\right)+\text{Li}_5\left(1/z\right)\right]\mathrm{ln}(z)`$
$`\frac{15}{8}\left[\text{Li}_4\left(z\right)\text{Li}_4\left(1/z\right)\right]\mathrm{ln}^2(z)`$
$`+\frac{5}{8}\left[\text{Li}_3\left(z\right)+\text{Li}_3\left(1/z\right)\right]\mathrm{ln}^3(z)`$
$`\frac{5}{32}\left[\text{Li}_2\left(z\right)\text{Li}_2\left(1/z\right)\right]\mathrm{ln}^4(z),`$ (17)
with $`\text{Ls}_5\left(\pi \right)=\frac{19}{240}\pi ^5`$, $`\text{Ls}_6\left(\pi \right)=\frac{45}{2}\pi \zeta _5+\frac{5}{4}\pi ^3\zeta _3`$.
Substituting (7), (11), (14)–(17) into Eq. (MZ-TH/00-22 Some remarks on the $`\epsilon `$-expansion of dimensionally regulated Feynman diagramsthanks: Based on the talk given by A. D. at the Zeuthen Workshop “Loops and Legs in Gauge Theories” (Bastei, Germany, April 2000).) (taking $`\theta =\tau _{01}^{}`$ or $`\theta =\tau _{02}^{}`$, denoting the $`z`$’s from Eq. (6) as $`z_1`$ and $`z_2`$, and setting $`\sigma =1`$), we arrive at the analytic continuation of the terms of the $`\epsilon `$-expansion, up to order $`\epsilon ^5`$. In fact, we have also obtained results for higher $`\text{Ls}_j`$ functions, up to $`j=10`$, which allowed us to reach the order $`\epsilon ^9`$.
It is instructive to consider the limit $`m_20`$. Introducing the variables $`x=m_1^2/k^2`$ and $`y=m_2^2/k^2`$ (and remembering that $`\sigma =1`$), we get
$`z_1{\displaystyle \frac{y}{(1x)^2}}+𝒪(y^2),z_2x+{\displaystyle \frac{2xy}{1x}}+𝒪(y^2),`$
$`\mathrm{ln}\left[\mathrm{\Delta }(m_1^2,m_2^2,k^2)/(k^2)^2\right]\text{i}\pi +2\mathrm{ln}(1x)`$
$`{\displaystyle \frac{2y(1+x)}{(1x)^2}}+𝒪(y^2).`$
Then, $`S_{a,b}(1/z_1)`$ can be converted into $`S_{a,b}(z_1)`$ by means of known relations given in Ref. . After this, the limit $`y0`$ ($`m_20`$) can be taken, since all $`\mathrm{ln}y`$ terms are cancelled. The obtained expressions can be simplified by transforming $`S_{a,b}(z_2)`$ into $`S_{a,b}(1/z_2)`$. In such a way, we also avoid appearance of terms like $`S_{a,b}(1)`$. Note that for this limit the terms up to order $`\epsilon ^3`$ can be extracted from Eq. (A.3) of ref. . Our expressions are in agreement with their results.
In fact, using hypergeometric representation
$`J^{(2)}(42\epsilon ;1,1)|_{\begin{array}{c}_{m_1=0}\hfill \\ ^{m_2m}\hfill \end{array}}=\text{i}\pi ^{2\epsilon }m^{2\epsilon }{\displaystyle \frac{\mathrm{\Gamma }(1+\epsilon )}{\epsilon (1\epsilon )}}`$ (20)
$`\times _2F_1\left(\begin{array}{c}1,\epsilon \\ 2\epsilon \end{array}|{\displaystyle \frac{k^2}{m^2}}\right)`$ (23)
(see, e.g., Eq. (10) of ), an arbitrary term of the $`\epsilon `$-expansion can be obtained. Employing Kummer’s relations for contiguous functions, one can transform the $`{}_{2}{}^{}F_{1}^{}`$ function from Eq. (20) into
$$\frac{1\epsilon }{12\epsilon }\left\{\frac{1+u}{2u}\frac{(1u)^2}{2u}{}_{2}{}^{}F_{1}^{}(\begin{array}{c}1,\mathrm{\hspace{0.33em}1}+\epsilon \\ 1\epsilon \end{array}|u)\right\},$$
with $`uk^2/m^2`$. This (transformed) $`{}_{2}{}^{}F_{1}^{}`$ function can be expressed in terms of a simple one-fold parametric integral,
$$(1u)^{12\epsilon }\left\{1\epsilon \underset{0}{\overset{1}{}}\frac{\text{d}t}{t}t^\epsilon \left[(1ut)^{2\epsilon }1\right]\right\}.$$
Expanding the integrand in $`\epsilon `$, we arrive at
$`J^{(2)}(42\epsilon ;1,1)|_{\begin{array}{c}_{m_1=0}\hfill \\ ^{m_2m}\hfill \end{array}}=\text{i}\pi ^{2\epsilon }m^{2\epsilon }{\displaystyle \frac{\mathrm{\Gamma }(1+\epsilon )}{(12\epsilon )}}`$ (26)
$`\times \{{\displaystyle \frac{1}{\epsilon }}{\displaystyle \frac{1u}{2u\epsilon }}[(1u)^{2\epsilon }1]`$
$`{\displaystyle \frac{(1u)^{12\epsilon }}{u}}{\displaystyle \underset{j=1}{\overset{\mathrm{}}{}}}\epsilon ^j{\displaystyle \underset{k=1}{\overset{j}{}}}(2)^{jk}S_{k,jk+1}(u)\},`$ (27)
which agrees with the results discussed earlier.
3. In ref. it was shown that similar explicit results can be constructed for the off-shell massless one-loop three-point function with external momenta $`p_1`$, $`p_2`$ and $`p_3`$ ($`p_1+p_2+p_3=0`$),
$`J(n;\nu _1,\nu _2,\nu _3|p_1^2,p_2^2,p_3^2)`$
$`{\displaystyle \frac{\text{d}^nr}{\left[(p_2r)^2\right]^{\nu _1}\left[(p_1+r)^2\right]^{\nu _2}(r^2)^{\nu _3}}},`$ (28)
as well as for the two-loop vacuum diagram with arbitrary masses $`m_1`$, $`m_2`$ and $`m_3`$,
$`I(n;\nu _1,\nu _2,\nu _3|m_1^2,m_2^2,m_3^2)`$
$`{\displaystyle \frac{\text{d}^np\text{d}^nq}{\left(p^2m_1^2\right)^{\nu _1}\left(q^2m_2^2\right)^{\nu _2}\left[(pq)^2m_3^2\right]^{\nu _3}}}.`$ (29)
Results for general $`n`$ and $`\nu _i`$ (in terms of hypergeometric functions of two variables) are available in Refs. . According to the magic connection , these integrals are closely related to each other. For example, in the case $`\nu _1=\nu _2=\nu _3=1`$ this connection (see Eq. (16) of ) yields
$`J(42\epsilon ;1,1,1)=\pi ^{3\epsilon }\text{i}^{1+2\epsilon }\left(p_1^2p_2^2p_3^2\right)^\epsilon `$
$`\times {\displaystyle \frac{\mathrm{\Gamma }(1+\epsilon )}{\mathrm{\Gamma }(12\epsilon )}}I(2+2\epsilon ;1,1,1),`$ (30)
where we assume that $`p_i^2m_i^2`$. Below we shall omit the arguments $`p_i^2`$ and $`m_i^2`$ in the integrals $`J`$ and $`I`$, respectively.
Then, using exact results in terms of $`{}_{2}{}^{}F_{1}^{}`$ functions , in combination with the formula
$`{\displaystyle \underset{j=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(2\epsilon )^j}{j!}}\text{Ls}_{j+1}(2\varphi )=2\pi {\displaystyle \frac{\mathrm{\Gamma }(12\epsilon )}{\mathrm{\Gamma }^2(1\epsilon )}}\theta (\mathrm{cos}\varphi )`$
$`{\displaystyle \frac{2^{12\epsilon }\mathrm{tan}\varphi }{(12\epsilon )\mathrm{sin}^{2\epsilon }\varphi }}{}_{2}{}^{}F_{1}^{}(\begin{array}{c}1,1/2\\ 3/2\epsilon \end{array}|\mathrm{tan}^2\varphi ),`$ (33)
the following results have been obtained in :
$`J(42\epsilon ;1,1,1)=2\pi ^{2\epsilon }\text{i}^{1+2\epsilon }{\displaystyle \frac{\mathrm{\Gamma }(1+\epsilon )\mathrm{\Gamma }^2(1\epsilon )}{\mathrm{\Gamma }(12\epsilon )}}`$
$`\times {\displaystyle \frac{\left[\mathrm{\Delta }(p_1^2,p_2^2,p_3^2)\right]^{1/2+\epsilon }}{(p_1^2p_2^2p_3^2)^\epsilon }}{\displaystyle \underset{j=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(2\epsilon )^j}{(j+1)!}}`$
$`\times \left[\text{Ls}_{j+2}(\pi ){\displaystyle \underset{i=1}{\overset{3}{}}}\left[\text{Ls}_{j+2}(\pi )\text{Ls}_{j+2}(2\varphi _i)\right]\right],`$ (34)
$`I(42\epsilon ;1,1,1)=\pi ^{42\epsilon }{\displaystyle \frac{\mathrm{\Gamma }^2(1+\epsilon )}{(1\epsilon )(12\epsilon )}}`$
$`\times \{{\displaystyle \frac{1}{2\epsilon ^2}}[{\displaystyle \frac{m_1^2+m_2^2m_3^2}{(m_1^2m_2^2)^\epsilon }}`$
$`+{\displaystyle \frac{m_2^2+m_3^2m_1^2}{(m_2^2m_3^2)^\epsilon }}+{\displaystyle \frac{m_3^2+m_1^2m_2^2}{(m_3^2m_1^2)^\epsilon }}]`$
$`+\left[\mathrm{\Delta }(m_1^2,m_2^2,m_3^2)\right]^{1/2\epsilon }{\displaystyle \underset{j=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(2\epsilon )^j}{(j+1)!}}`$
$`\times [\text{Ls}_{j+2}(\pi ){\displaystyle \underset{i=1}{\overset{3}{}}}[\text{Ls}_{j+2}(\pi )\text{Ls}_{j+2}(2\varphi _i)]]\},`$ (35)
where the angles $`\varphi _i`$ ($`i=1,2,3`$) are defined via
$`\mathrm{cos}\varphi _1=(p_2^2+p_3^2p_1^2)/(2\sqrt{p_2^2p_3^2}),`$
$`\mathrm{cos}\varphi _2=(p_3^2+p_1^2p_2^2)/(2\sqrt{p_3^2p_1^2}),`$
$`\mathrm{cos}\varphi _3=(p_1^2+p_2^2p_3^2)/(2\sqrt{p_1^2p_2^2})`$ (36)
(remember that $`p_i^2m_i^2`$ for the integrals $`I`$), so that $`\varphi _1+\varphi _2+\varphi _3=\pi `$. Note that the angles $`\theta _i`$ from are related to $`\varphi _i`$ as $`\theta _i=2\varphi _i`$. By analogy with the two-point case (MZ-TH/00-22 Some remarks on the $`\epsilon `$-expansion of dimensionally regulated Feynman diagramsthanks: Based on the talk given by A. D. at the Zeuthen Workshop “Loops and Legs in Gauge Theories” (Bastei, Germany, April 2000).), the angles $`\varphi _i`$ can be understood as the angles of a triangle whose sides are $`\sqrt{p_1^2}`$, $`\sqrt{p_2^2}`$ and $`\sqrt{p_3^2}`$, whereas its area is $`\frac{1}{4}\sqrt{\mathrm{\Delta }(p_1^2,p_2^2,p_3^2)}`$.
For the two lowest orders ($`\epsilon ^0`$ and $`\epsilon ^1`$), we reproduce eqs. (9)–(10) from . Useful representations for the $`\epsilon ^0`$ terms of both types of diagrams can also be found in . Moreover, in Eq. (26) of a one-fold integral representation for $`J(42\epsilon ;1,1,1)`$ is presented (for its generalization, see Eq. (7) of ). Expanding the integrand in $`\epsilon `$, we were able to confirm the $`\epsilon `$-expansion (MZ-TH/00-22 Some remarks on the $`\epsilon `$-expansion of dimensionally regulated Feynman diagramsthanks: Based on the talk given by A. D. at the Zeuthen Workshop “Loops and Legs in Gauge Theories” (Bastei, Germany, April 2000).) numerically.
To construct the analytic continuation of the terms of the $`\epsilon `$-expansion (MZ-TH/00-22 Some remarks on the $`\epsilon `$-expansion of dimensionally regulated Feynman diagramsthanks: Based on the talk given by A. D. at the Zeuthen Workshop “Loops and Legs in Gauge Theories” (Bastei, Germany, April 2000).) and (MZ-TH/00-22 Some remarks on the $`\epsilon `$-expansion of dimensionally regulated Feynman diagramsthanks: Based on the talk given by A. D. at the Zeuthen Workshop “Loops and Legs in Gauge Theories” (Bastei, Germany, April 2000).), we need just to apply substitutions (7), (11), (14)–(17), with $`\theta =2\varphi _i`$ ($`i=1,2,3`$). The remaining terms $`\text{Ls}_{j+2}\left(\pi \right)`$ can actually be treated in the same way, if we substitute $`\theta =0`$. In any case, their values are known (in terms of $`\zeta `$ function) and can be summed into a combination of $`\mathrm{\Gamma }`$ functions, Eq. (5). We get three variables $`z_i=e^{2\mathrm{i}\sigma \varphi _i}`$, see Eq. (6), such that $`z_1z_2z_3=1`$. The causal prescription requires to take $`\sigma =+1`$ for the $`J`$-integrals and $`\sigma =1`$ for the $`I`$-integrals. Using the substitutions presented above, we obtain the analytic continuation of the results (MZ-TH/00-22 Some remarks on the $`\epsilon `$-expansion of dimensionally regulated Feynman diagramsthanks: Based on the talk given by A. D. at the Zeuthen Workshop “Loops and Legs in Gauge Theories” (Bastei, Germany, April 2000).) and (MZ-TH/00-22 Some remarks on the $`\epsilon `$-expansion of dimensionally regulated Feynman diagramsthanks: Based on the talk given by A. D. at the Zeuthen Workshop “Loops and Legs in Gauge Theories” (Bastei, Germany, April 2000).) up to $`\epsilon ^4`$. We note that the result for the $`\epsilon `$-term was known (in terms of $`\text{Li}_3`$).
When the masses are equal, $`m_1=m_2=m_3m`$ (this also applies to the symmetric case $`p_1^2=p_2^2=p_3^2p^2`$), the three angles $`\varphi _i`$ are all equal to $`\pi /3`$, whereas $`\mathrm{\Delta }(m^2,m^2,m^2)=3m^4`$. Therefore, in this case the r.h.s. of Eq. (MZ-TH/00-22 Some remarks on the $`\epsilon `$-expansion of dimensionally regulated Feynman diagramsthanks: Based on the talk given by A. D. at the Zeuthen Workshop “Loops and Legs in Gauge Theories” (Bastei, Germany, April 2000).) becomes
$`\pi ^{42\epsilon }{\displaystyle \frac{\mathrm{\Gamma }^2(1+\epsilon )m^{24\epsilon }}{(1\epsilon )(12\epsilon )}}\{{\displaystyle \frac{3}{2\epsilon ^2}}+{\displaystyle \frac{\sqrt{3}}{3^\epsilon }}{\displaystyle \underset{j=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(2\epsilon )^j}{(j+1)!}}`$
$`\times [3\text{Ls}_{j+2}\left(\frac{2\pi }{3}\right)2\text{Ls}_{j+2}(\pi )]\}.`$ (37)
For instance, in the contribution of order $`\epsilon `$ the transcendental constant $`\text{Ls}_3(2\pi /3)`$ appears. This constant was discussed in detail in . The fact that $`\text{Ls}_3(2\pi /3)`$ occurs in certain two-loop on-shell integrals and three-loop vacuum integrals has been noticed in . Moreover, in it was observed that the higher-$`j`$ terms from (MZ-TH/00-22 Some remarks on the $`\epsilon `$-expansion of dimensionally regulated Feynman diagramsthanks: Based on the talk given by A. D. at the Zeuthen Workshop “Loops and Legs in Gauge Theories” (Bastei, Germany, April 2000).) form a basis for certain on-shell integrals with a single mass parameter. Connection of $`\text{Ls}_3(2\pi /3)`$ with multiple binomial sums is discussed in . We also note that in the constant $`\text{Ls}_3(\pi /2)`$ appeared.
4. One of the interesting problems is to construct terms of the $`\epsilon `$-expansion for the one-loop three-point function with general masses. In this sense, the geometrical description seems to be rather instructive. The geometrical approach to the three-point function is discussed in section V of (see also in ). This function can be represented as an integral over a spherical (or hyperbolic) triangle, as shown in Fig. 6 of , with a weight factor $`1/\mathrm{cos}^{12\epsilon }\theta `$ (see eqs. (3.38)–(3.39) of ). This triangle 123 is split into three triangles 012, 023 and 031. Then, each of them is split into two rectangular triangles, according to Fig. 9 of . We consider the contribution of one of the six resulting triangles, namely the left rectangular triangle in Fig. 9. Its angle at the vertex 0 is denoted as $`\frac{1}{2}\phi _{12}^+`$, whereas the height dropped from the vertex 0 is denoted $`\eta _{12}`$.
The remaining angular integration is (see eq. (5.16) of )
$`{\displaystyle \frac{1}{2\epsilon }}{\displaystyle \underset{0}{\overset{\phi _{12}^+/2}{}}}\text{d}\phi \left[1\left(1+{\displaystyle \frac{\mathrm{tan}^2\eta _{12}}{\mathrm{cos}^2\phi }}\right)^\epsilon \right]`$
$`=\frac{1}{2}{\displaystyle \underset{j=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(\epsilon )^j}{(j+1)!}}{\displaystyle \underset{0}{\overset{\phi _{12}^+/2}{}}}\text{d}\phi \mathrm{ln}^{j+1}\left(1+{\displaystyle \frac{\mathrm{tan}^2\eta _{12}}{\mathrm{cos}^2\phi }}\right).`$ (38)
First of all, we note that the l.h.s. of Eq. (MZ-TH/00-22 Some remarks on the $`\epsilon `$-expansion of dimensionally regulated Feynman diagramsthanks: Based on the talk given by A. D. at the Zeuthen Workshop “Loops and Legs in Gauge Theories” (Bastei, Germany, April 2000).) yields a representation valid for an arbitrary $`\epsilon `$ (i.e., in any dimension). To get the result for the general three-point function, we need to consider a sum of six such integrals. The resulting representation is closely related to the representation in terms of hypergeometric functions of two arguments (see also in for some special cases).
In the limit $`\epsilon 0`$ we get a combination of $`\text{Cl}_2`$ functions, eq. (5.17) of . Collecting the results for all six triangles, we get the result for the three-point function with arbitrary masses and external momenta, corresponding (at $`\epsilon =0`$) to the analytic continuation of the well-known formula presented in . The higher terms of the $`\epsilon `$-expansion correspond to the angular integrals on the r.h.s. of Eq. (MZ-TH/00-22 Some remarks on the $`\epsilon `$-expansion of dimensionally regulated Feynman diagramsthanks: Based on the talk given by A. D. at the Zeuthen Workshop “Loops and Legs in Gauge Theories” (Bastei, Germany, April 2000).). The problem of constructing closed representations for these terms, as well as their analytic continuation, is very important. We note that the $`\epsilon `$-term of the three-point function with general masses has been calculated in in terms of $`\text{Li}_3`$.
5. We have shown that the compact structure of the coefficients of the $`\epsilon `$-expansion of the two-point function (MZ-TH/00-22 Some remarks on the $`\epsilon `$-expansion of dimensionally regulated Feynman diagramsthanks: Based on the talk given by A. D. at the Zeuthen Workshop “Loops and Legs in Gauge Theories” (Bastei, Germany, April 2000).), the massless off-shell three-point function (MZ-TH/00-22 Some remarks on the $`\epsilon `$-expansion of dimensionally regulated Feynman diagramsthanks: Based on the talk given by A. D. at the Zeuthen Workshop “Loops and Legs in Gauge Theories” (Bastei, Germany, April 2000).) and two-loop massive vacuum diagrams (MZ-TH/00-22 Some remarks on the $`\epsilon `$-expansion of dimensionally regulated Feynman diagramsthanks: Based on the talk given by A. D. at the Zeuthen Workshop “Loops and Legs in Gauge Theories” (Bastei, Germany, April 2000).), in terms of log-sine integrals, allows to perform analytic continuation in terms of generalized Nielsen polylogarithms (12), in some cases (26) even for an arbitrary order of the $`\epsilon `$-expansion. It is likely that a further generalization of these results is possible, e.g. for the three-point function with different masses, two-point integrals with two (and more) loops and three-loop vacuum integrals. In particular, numerical analysis of the coefficients of the expansion of certain two-point on-shell integrals and three-loop vacuum integrals shows that in some cases the values of generalized log-sine integrals $`\text{Ls}_j^{(k)}`$, Eq. (13), may be involved. For instance, in ref. it was shown that $`\text{Ls}_4^{(1)}\left(2\pi /3\right)`$ is connected with $`V_{3,1}`$ from .
The fact that the generalization of $`\text{Ls}_2=\text{Cl}_2`$ goes in the $`\text{Ls}_j`$ direction, rather than in $`\text{Cl}_j`$ direction (see Eq. (10)), is very interesting. There is another example , the off-shell massless ladder three- and four-point diagrams with an arbitrary number of loops, when such a generalization went in the $`\text{Cl}_j`$ direction (for details, see ). It could be also noted that the the two-loop non-planar (crossed) three-point diagram gives in this case the square of the one-loop function, $`(\text{Cl}_2(\theta ))^2`$ (cf. Eq. (23) of ), leading to the structure $`(\text{Cl}_2(\pi /3))^2`$ in the symmetric ($`p_i^2=p^2`$) case. Recently, these constants have been also found in massive three-loop calculations .
The construction of analytic continuation of the generalized log-sine functions should be investigated in more detail. In fact, it may require including some other generalizations of polylogarithms (see, e.g., in Ref. ).
Acknowledgements. A. D. would like to thank the organizers of ‘Loops and Legs 2000’, it was really a very useful conference. A. D.’s research and participation in the conference were supported by DFG. M. K. is grateful to the THEP group (University of Mainz) for their hospitality during his research stay, which was supported by BMBF under contract 05 HT9UMB 4. At an earlier stage (before November 1999), A. D.’s research was supported by the Alexander von Humboldt Foundation; also partial support from the grants RFBR No. 98–02–16981 and Volkswagen No. I/73611 is acknowledged.
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# From Stars to Super-planets: the Low-Mass IMF in the Young Cluster IC348 1footnote 11footnote 1Based on observations with the NASA/ESA Hubble Space Telescope, obtained at the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc., under NASA contract No. NAS5-26555.
## 1 Introduction
The low-mass end of the stellar initial mass function (IMF) is of interest for our understanding of both baryonic dark matter in the Galaxy and, perhaps more importantly, the formation processes governing stars, brown dwarfs, and planets. In the stellar mass regime, the complex interplay between a wide array of physical processes is believed to determine the eventual outcome of the star formation process, the masses of stars. These diverse processes include those that govern molecular cloud structure and evolution, subsequent gravitational collapse, disk accretion, stellar winds, multiplicity, and stellar mergers. What is the distribution of object masses that results from the interaction between these processes? Do the same processes that form stars also produce less massive objects extending into the brown dwarf and planetary regimes? While such questions can be answered directly by constructing inventories of stellar and substellar objects, there is also the hope that the same set of data can shed light on the nature of the interaction between the physical processes and, thereby, bring us closer to a predictive theory of star and brown dwarf formation.
While the stellar IMF has long been studied (e.g., Salpeter (1955)), the very low-mass and substellar IMF is much less well known since the very existence of substellar objects has only recently been demonstrated, and reliable inventories of substellar objects are only now becoming available. The Pleiades has proven to be one of the most popular sites for low-mass IMF studies both due to its proximity ($`d125`$ pc) and because it is at an age ($`100`$ Myr) at which our understanding of stellar evolution is fairly robust. The large area subtended by the Pleiades poses several challenges: studies of the low mass IMF must survey large areas and distinguish low mass cluster members from the growing Galactic interloper population at faint magnitudes. For example, recent deep imaging surveys of the Pleiades carried out over several square degrees have used broad band color selection criteria to probe the cluster IMF to masses below the hydrogen burning limit (e.g., to $`0.04M_{}`$; Bouvier et al. 1999), where the fraction of objects that are cluster members is much less than 1%.
In a complementary development, new large area surveys (e.g., 2MASS, DENIS, and SDSS) are now probing the low mass IMF of the field population in the solar neighborhood, extending into the substellar regime. In an account of the progress to date, Reid et al. (1999) model the spectral type distribution of the low mass population drawn from 2MASS and DENIS samples obtained over several hundred square degrees in order to constrain the low mass IMF. Since substellar objects cool as they age, the observed spectral type distribution depends on both the mass and age distributions of the local field population. As a result, the lack of strong constraints on the age distribution poses a challenge for the determination of the field IMF at low masses. For example, assuming a flat age distribution over $`010`$ Gyr, Reid et al. find an IMF that is fairly flat, $`dN/d\mathrm{log}MM^\alpha `$ where $`\alpha `$ -1 to 0, where the uncertainty in the slope does not include the uncertainty in the age distribution of the population.
In comparison with the solar neighborhood and older open clusters such as the Pleiades, young stellar clusters ($`10`$ Myr) are a complementary and advantageous environment in which to carry out low-mass IMF studies. As in the situation for the Pleiades, stars in young clusters share a common distance and metallicity and, at low masses, are much brighter due to their youth. As a well recognized consequence, it is possible to readily detect and study even objects much below the hydrogen-burning limit. In addition, young clusters also offer some significant advantages over the older open clusters. For example, since young clusters are less dynamically evolved than older open clusters, the effects of mass segregation and the evaporation of low mass cluster members are less severe. Since young clusters are less dynamically evolved, they also subtend a more compact region on the sky. As a result, the fractional foreground and background contamination is much reduced and reasonable stellar population statistics can be obtained by surveying small regions of the sky. These advantages are (of course) accompanied by challenges associated with the study of young environments. These include the need to correct for both differential reddening toward individual stars and infrared excess, the excess continuum emission that is believed to arise from circumstellar disks. Pre-main sequence evolutionary tracks pose the greatest challenge to the interpretation of the observations because the tracks have little observational verification, especially at low masses and young ages. The temperature calibration for low-mass pre-main-sequence stars is an additional uncertainty.
While thus far the luminosity advantage of young clusters has been used with great success to detect some very low mass cluster members (e.g., $`0.02M_{}`$ objects in IC348 \[Luhman 1999\] and the $`\sigma `$ Ori cluster \[Zapatero Osorio et al. 2000\]), attempts to study the low mass IMF in young clusters have stalled at much higher masses, in the vicinity of the hydrogen burning limit (e.g., the Orion Nebula Cluster—Hillenbrand 1997), due to the need for complete sampling to low masses and potentially large extinctions. Since reddening and IR excesses can greatly complicate the determination of stellar masses from broad band photometry alone (e.g., Meyer et al. 1997), stellar spectral classification to faint magnitudes, an often time-consuming task, is typically required.
Stellar spectral classification in young clusters has been carried out using a variety of spectroscopic methods. These include the use of narrow atomic and molecular features in the $`K`$-band (e.g., Ali et al. (1995); Greene & Meyer (1995); Luhman et al. 1998, hereinafter LRLL), the $`H`$-band (e.g., Meyer (1996)), and the $`I`$-band (e.g., Hillenbrand (1997)), each of which have their advantages. While spectral classification at the longer wavelengths is better able to penetrate higher extinctions, spectral classification at the shorter wavelengths is less affected by infrared excess. With the use of high spectral resolution and the availability of multiple stellar spectral features, it is possible to diagnose and correct for infrared excess. This technique has been used with great success at optical wavelengths in the study of T Tauri star photospheres (e.g., Hartigan et al. (1989)). Alternatively, the difficulty of correcting for infrared excess can be avoided to a large extent by studying somewhat older (5-20 Myr old) clusters, in which infrared excesses are largely absent but significant dynamical evolution has not yet occurred.
In this paper, we develop an alternative, efficient method of spectral classification: filter photometric measures of water absorption band strength as an indicator of stellar spectral type. Water bands dominate the infrared spectra of M stars and are highly temperature sensitive, increasing in strength with decreasing effective temperature down to the coolest M dwarfs known ($`2000`$K; e.g., Jones et al. (1994)). The strength of the water bands and their rapid variation with effective temperature, in principle, allows the precise measurement of spectral type from moderate signal-to-noise photometry. At the same time, water bands are relatively insensitive to gravity (e.g., Jones et al. (1995)), particularly above 3000K, becoming more sensitive at lower temperatures where dust formation is an added complication (e.g., the Ames-Dusty models; Allard et al. (2000); Allard 1998b ). Synthetic atmospheres (e.g., NextGen: Hauschildt et al. (1999); Allard et al. (1997)) also indicate a modest dependence of water band strength on metallicity (e.g., Jones et al. (1995)).
Because strong absorption by water in the Earth’s atmosphere can complicate the ground-based measurement of the depth of water bands, we used HST NICMOS filter photometry to carry out the measurements. The breadth of the water absorption bands requires that any measure of band strength adequately account for the effects of reddening. Consequently, we used a 3 filter system to construct a reddening independent index that measures the band strength. Of the filters available with NICMOS, only the narrow band F166N, F190N, and F215N filters which sample the depth of the 1.9 $`\mu `$m water band proved suitable. On the one hand, the narrow filter widths had the advantages of excluding possible stellar or nebular line emission and limiting the differential reddening across the bandpass. On the other hand, similar filters with broader band passes would have made it feasible to study much fainter sources, e.g., in richer clusters at much larger distances. Despite the latter difficulty, there were suitable nearby clusters such as IC348 to which this technique could be profitably applied.
IC348 is a compact, young cluster located near an edge of the Perseus molecular cloud. It has a significant history of optical study (see, e.g., Herbig (1998) for a review), and because of its proximity ($`d300`$pc), youth ($`<10`$ Myr), and rich, compact nature, both the star formation history and the mass function that characterizes the cluster have been the subject of several recent studies.
Ground-based $`J`$, $`H`$, $`K`$ imaging of the cluster complete to $`K`$=14 (Lada & Lada (1995)) revealed signficant spatial structure, in which the richest stellar grouping is the “a” subcluster ($`r=3.5^{}`$; hereinafter IC348a) with approximately half of the cluster members. The near-IR colors indicate that IC348 is an advantageous environment in which to study the stellar properties of a young cluster since only a moderate fraction of cluster members possess near-IR excesses ($``$20% for the cluster overall; $``$ 12% for IC348a) and most cluster members suffer moderate extinction ($`A_V5`$ with a spread to $`A_V`$ $`>`$ 20). Lada & Lada (1995) showed that the $`K`$-band luminosity function of IC348 is consistent with a history of continuous star formation over the last $`57`$ Myr and a time-independent Miller-Scalo IMF in the mass range $`0.120M_{}`$. The inferred mean age of a few Myr is generally consistent with the lack of a significant population of excess sources since disks are believed to disperse on a comparable timescale (Meyer et al. (2000)).
Herbig (1998) subsequently confirmed a significant age spread to the cluster ($`0.712`$ Myr) based on BVRI imaging of a $`7\mathrm{}\times 12\mathrm{}`$ region, which included much of IC348a, and $`R`$-band spectroscopy of a subset of sources in the field. In the mass range in which the study is complete ($`M_{}>0.3M_{}`$), the mass function slope was found to be consistent with that of Scalo (1986). A more detailed study of a $`5\mathrm{}\times 5\mathrm{}`$ region centered on IC348a was carried out by LRLL using IR and optical spectroscopy complete to $`K=12.5`$. They also found an age spread to the subcluster ($`510`$ Myr), a mean age of $`3`$ Myr, and evidence for a substellar population. The mass function of the subcluster was found to be consistent with Miller & Scalo (1979) in the mass range $`0.253M_{}`$ (i.e., flatter in slope than deduced by Herbig) and flatter than Miller-Scalo at masses below $`0.25M_{}`$; however, completeness corrections were significant below $`0.1M_{}`$. Luhman (1999) has further probed the substellar population of IC348 using optical spectral classification of additional sources ($`I19.5`$) both in and beyond the $`5^{}\times 5^{}`$ core.
In this paper, we extend previous studies of IC348 by probing 4 magnitudes below the $`K`$ spectral completeness limit of LRLL , enabling a more detailed look at the population in the low-mass stellar and substellar regimes. We find that, with our spectral classification technique, our measurement of the IMF in IC348 is complete to the deuterium burning limit ($`0.015M_{}`$), a fiducial boundary between brown dwarf and planetary mass objects (e.g., Saumon et al. 2000; Zapatero Osorio et al. 2000). To avoid potential misunderstanding, we note that this boundary is only very approximate. A precise division between the brown dwarf and planetary regimes is unavailable and perhaps unattainable in the near future given the current disagreement over fundamental issues regarding the definition of the term “planet”. These include whether the distinction between brown dwarfs and planets should be made in terms of mass or formation history (e.g., gravitational collapse vs. accumulation) and whether planetary mass objects that are not companions can even be considered to be “planets”. Here, we hope to side-step such a discussion at the outset and, instead, explore how the IMF of isolated objects over the range from $`1M_{}`$ to $`0.015M_{}`$, once measured, can advance the discussion, i.e., provide clues to the formation and evolutionary histories of stellar and substellar objects. The HST observations are presented in section 2. The resulting astrometry and near-infrared luminosity functions are discussed in sections 3 and 4. In section 5, we discuss the calibration of the water index and the determination of stellar spectral types. The reddening corrections are discussed in section 6, and the resulting observational HR diagram in section 7. In section 8, we identify the interloper population and compare the cluster population with the predictions of pre-main sequence evolutionary tracks. Given these results, in section 9, we identify possible cluster binaries and derive a mass function for the cluster. Finally, in section 10, we present our conclusions.
## 2 Observations, Data Reduction, and Calibration
### 2.1 Photometry
We obtained HST NIC3 narrow band photometry for 50 $`(51^{\prime \prime }\times 51^{\prime \prime })`$ fields in the IC348a subcluster, nominally centered at $`\alpha =3^\mathrm{h}44^\mathrm{m}31\stackrel{\mathrm{s}}{\mathrm{.}}9`$, $`\delta =32^{}09\mathrm{}54\stackrel{}{\mathrm{.}}2`$ (J2000). The NICMOS instrument and its on-orbit performance have been described by Thompson et al. (1998) and Calzetti & Noll (1998). Figure 1 shows the relative positions of the fields with respect to the $`5\mathrm{}\times 5\mathrm{}`$ core of the subcluster. The NIC3 field positions were chosen to avoid bright stars much above the saturation limit ($`K9`$) and to maximize area coverage. As a result, the fields are largely non-overlapping, covering most of the $`5\mathrm{}\times 5\mathrm{}`$ core and a total area of $`34.76`$ sq. arcmin. Each field was imaged in the narrow band F166N, F190N, and F215N filters, centered at 1.66 $`\mu `$m, 1.90 $`\mu `$m, 2.15 $`\mu `$m respectively, at two dither positions separated by $`5.1\mathrm{}`$. The exposure time at each dither position was 128 seconds, obtained through four reads of the NIC3 array in the SPARS64 MULTIACCUM sequence, for a total exposure time in each field of 256 seconds.
To calibrate the non-standard NIC3 colors, we observed a set of 23 standard stars chosen to cover spectral types K2 through M9 that have the kinematics and/or colors typical of solar neighborhood disk stars (e.g., Leggett (1992); see Table 1) and, therefore, are likely to have metallicities similar to that of the cluster stars. Although most of the standard stars were main-sequence dwarfs, we also observed a few pre-main sequence stars in order to explore the effect of lower gravity. We chose for this purpose pre-main sequence stars known to have low infrared excesses (weak lined T Tauri stars; WTTS) so that the observed flux would be dominated by the stellar photosphere. The standard stars were observed in each of the F166N, F190N, and F215N filters and with the G141 and G206 grisms. The stars were observed with each spectral element at two or three dither positions separated by $`5.1\mathrm{}`$ in MULTIACCUM mode.
Since NICMOS does not have a shutter, the bright standard stars could potentially saturate the array as the NIC3 filter wheel rotates through the broad or intermediate band filters located between the narrow band filters and grisms used in the program. To avoid the resulting persistence image that would compromise the photometric accuracy, dummy exposures, taken at a position offset from where the science exposure would be made, were inserted between the science exposures in order to position the filter wheel at the desired spectral element before actually taking the science exposure.
Much of the data for IC348 (45 of the 50 fields) and all of the data for the standard stars were obtained during the first (January 12 – February 1, 1998) and second (June 4 – 28, 1998) NIC3 campaigns in which the HST secondary was moved to bring NIC3 into focus. A log of our observations is provided in Table 2. The data were processed through the usual NICMOS calnic pipeline (version 3.2) with the addition of one step. After the cosmic ray identification, column bias offsets were removed from the final readout in order to eliminate the “banding” (constant, incremental offsets of $`30`$ counts about 40 columns wide) present in the raw data.
No residual reflection nebulosity is noticeable in the reduced (dither-subtracted) images. Consequently, removal of nebular emission was not a concern for the stellar photometry. To perform the stellar photometry, we first identified sources in each of the images using the IRAF routine daofind. Due to the strongly varying noise characteristics of the NIC3 array, daofind erroneously identified numerous noise peaks as point sources, and so the detections were inspected frame by frame to eliminate spurious detections. A detection was considered to be real if the source was detected in both the F215N and F190N frames. With these identification criteria, we were likely to obtain robust detections of heavily extincted objects (in F215N) as well as spectral types for all identified sources, F190N typically having the lowest flux level at late spectral types.
Since the frames are sparsely populated, we used the aperture photometry routine phot to measure the flux of each identified source. To optimize the signal-to-noise of the photometry on faint objects ($`K16`$), we adopted a 4-pixel radius photometric aperture that included the core of the PSF and $`91\%`$ of the total point source flux (the exact value varied by about $`1\%`$ from filter to filter) with an uncertainty in the aperture correction of $`<1\%`$ in all filters. The aperture correction was derived from observations of calibration standards and/or bright, unsaturated objects in the IC348 fields. Despite the difference in focus conditions between the data taken in and out of the NIC3 campaigns, the aperture corrections were statistically identical. As a result, the same aperture and procedures were used for both data sets. The conversion from ADU/s to both Janskys and magnitudes was made using the photometric constants kindly provided by M. Rieke (1999, personal communication). These constants are tabulated in Table 3.
### 2.2 Spectroscopy
In order to confirm the calibration of the filter photometric water index against stellar spectral type, we also obtained NIC3 G141 and G206 grism spectra for 17 of our 23 standard stars. The spectral images were processed identically to the photometric images, including the removal of the bias jumps. The spectra were extracted using NICMOSlook (version 2.6.5; Pirzkal & Freudling 1998a ), the interactive version of the standard pipeline tool (CalnicC; Pirzkal & Freudling 1998b ) for the extraction of NIC3 grism spectra. The details of the extraction process and subsequent analysis are presented in Tiede et al. (2000). The 1.9 $`\mu `$m H<sub>2</sub>O band strengths obtained from a preliminary analysis of the spectra were found to be consistent with the filter photometric results reported in section 5.
### 2.3 Intrapixel Sensitivity and Photometric Accuracy
Because infrared arrays may have sensitivity variations at the sub-pixel scale, the detected flux from an object, when measured with an undersampled PSF, may depend sensitively on the precise position of the object within in a pixel. As shown by Lauer (1999), such intrapixel sensitivity effects can be significant when working with undersampled NIC3 data ($`0.2\mathrm{}`$ pixels). To help us quantify the impact of this effect on our data set, Lauer kindly calculated for us the expected intrapixel dependence of the detected flux from a point source as a function of intrapixel position, using TinyTim PSFs appropriate for the filters in our study and the NIC3 intrapixel response function deduced in Lauer (1999). As expected, the intrapixel sensitivity effect is more severe at shorter wavelengths where the undersampling is more extreme. In the F215N filter, the effect is negligible: the variation in the detected flux as a function of intrapixel position is within $`\pm 0.3`$% of the flux that would be detected with a well sampled PSF. For the F190N and F166N filters, the same quantity varies within $`\pm 3.5`$% and $`\pm 8.5`$%, respectively.
Although intrapixel sensitivity can be severe at the shorter wavelengths, the effect on photometric colors is mitigated if the intrapixel response is similar for the three filters (the assumption made here) and the sub-pixel positional offsets between the observations in each filter are small. For example, with no positional offset between the 3 filters, the error in the reddening independent water index, $`Q_{\mathrm{H2O}}`$, discussed in section 5, is $`<`$1% which impacts negligibly on our conclusions. Since pointing with HST is expected to be accurate to better than a few milliarcseconds for the $`17`$ minute duration of the observations on a given cluster field (M. Lallo 1999, personal communication), pointing drifts are unlikely to introduce significant positional offsets. The HST jitter data for our observations confirm the expected pointing accuracy. Over the $`5`$ minute duration of the exposure in a single filter, the RMS pointing error is on average $`4`$ milliarcseconds (0.02 NIC3 pixels).
Systematic positional offsets between filters could also arise from differing geometric transformations between the filters. To test this, we examined the centroid position of the bright cluster sources and standard stars for individual dither positions in each filter. No systematic differences in centroid positions between filters were found. The 1–$`\sigma `$ scatter about the mean was 0.05 pixels which represents the combination of our centroiding accuracy and any true positional variations. To quantify the impact of the latter possibility on our results, random positional variations of 0.05 pixels in each filter translate into a maximal error in $`Q_{\mathrm{H2O}}`$ of less than $`\pm `$4%.
## 3 Astrometry
Because three of the recent studies of IC348 (Herbig (1998), LRLL , and Luhman 1999) have examined regions surrounding and including IC348a, we can directly compare the previous results with ours via the overlaps in the stellar samples. Figure 2 shows the spatial distribution of the samples from the previous and present studies. The present study covers a more compact region than the previous studies, but is complete to much greater depth.
Table From Stars to Super-planets: the Low-Mass IMF in the Young Cluster IC348 <sup>1</sup><sup>1</sup>1Based on observations with the NASA/ESA Hubble Space Telescope, obtained at the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc., under NASA contract No. NAS5-26555. presents the source designations for all of the stars in our sample, the corresponding designations from previous studies, and the J2000 celestial coordinates of each star. Our designations are comprised of the 3-digit field number followed by the 2-digit number of the star in that field. For example, 021-05 is from field 021 and is star number 5 in that field. The celestial coordinates in Table From Stars to Super-planets: the Low-Mass IMF in the Young Cluster IC348 <sup>1</sup><sup>1</sup>1Based on observations with the NASA/ESA Hubble Space Telescope, obtained at the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc., under NASA contract No. NAS5-26555. are based on the NICMOS header values associated with the central pixel in each field. The total error in the relative accuracy of the coordinates due to photometric centroiding, geometric field distortion, and repeat pointing errors, are estimated to be $`0\stackrel{}{\mathrm{.}}2`$ per star. This error is a function of the stellar position in the NIC3 field of view: stars located toward the corners of a frame have larger errors primarily due to field distortion which we have not attempted to correct. While absolute astrometry is not required for the present study, we can obtain an estimate of the absolute astrometric error by comparing our coordinates to those obtained in previous investigations. Comparison with the celestial coordinates reported in LRLL typically resulted in disagreements of less than $`1\mathrm{}`$.
## 4 Completeness and Luminosity Functions
### 4.1 Completeness and Photometric Accuracy
At the bright end, our sample is limited by saturation. Inspection of the error flags output by CALNICA implied that our saturation limits are $`10.96\pm 0.49`$ magnitudes in F166N, $`10.89\pm 0.44`$ in F190N; and $`10.62\pm 0.35`$ in F215N. The flux range over which saturation occured reflects the sensitivity variation across the array and the variation in the intrapixel position of individual stars.
Given the noise characteristics of, and significant quantum efficiency variations across, the NIC3 array, we used simulated data to evaluate the efficiency of our detection algorithm at the faint end and the accuracy of our photometric measurements. We first added to a representative frame for each filter a known number of point sources, positioned randomly within the frame, with known magnitudes and zero color, then performed detection and stellar photometry on the frames in a method identical to those used for the real data. Since crowding was not an issue in the real frames, care was taken to ensure that none of the artificial stars where lost to superposition. While we did not explore the full color range of the actual data set, the adopted simulation was sufficient to obtain a robust estimate of our detection efficiency in the individual filters.
Artificial PSFs were generated using the program TinyTim version 4.4 (Krist & Hook (1997)). Each artificial PSF was created with a factor of 10 oversampling, i.e, in a 240 $`\times `$ 240 grid with each element of the grid representing $`0.02\mathrm{}`$ on the sky, to facilitate sub-pixel interpolation in positioning the artificial stars. The extent of the artificial PSF ($`2.4\mathrm{}`$) was chosen to equal the radius at which the flux level for even the brightest stars in the data set is less than the noise fluctuations in the background.
Inspection of the empirical luminosity functions, the theoretical photometric errors, and signal-to-noise values indicated that our sample was likely complete to $`17.5`$ mag (0.1 mJy in F215N). To derive the completeness limit quantitatively for each band, we created two sets of artificial stars to be added and recovered from a representative frame in each band. The first set of 50 stars was linearly distributed over the magnitude range in which photometric errors become significant (15.0 to 19.5). The second set of 50 stars was linearly distributed between 17.0 and 18.5 magnitude in order to “zero-in” on the $`100\%`$ completeness limit. After the addition of the artificial stars with the appropriate noise, each of the images was photometrically processed in a manner identical to the real data frames.
The completeness as a function of F215N magnitude is displayed in Fig. 3. The results are essentially identical for F190N. The Figure shows the number of stars input into (solid line) and the number detected in (dotted line) each 0.5 magnitude bin. Our photometry is $`100\%`$ complete through the bin centered at 17.25 magnitudes, beyond which the detection efficiency drops rapidly. It is 80% at 17.75, 11% at 18.25, and finally no detections beyond 18.5. When the results are tabulated in 0.1 magnitude bins, we find that we are $`100\%`$ complete to 17.6 magnitudes. Since the last $`100\%`$ complete bin only contains 5 stars and because the rest of analysis is done in 0.5 magnitude increments, we adopt 17.5 magnitudes as a conservative estimate of our $`100\%`$ completeness limit.
In addition to calculating the completeness limit, the artificial stars also allowed us to gauge the accuracy of our photometry and photometric error estimates. Since we knew the magnitudes of the artificial stars that we added to the frame, we could calculate the “True Error” of each photometric measurement (True Error $``$ measured magnitude $``$ input magnitude). The top panels of Fig. 4 show the absolute value of the resulting true errors as a function of input magnitude. For each photometric measurement, we calculated the photometric uncertainty due to photon statistics. The bottom panels of Fig. 4 show this estimated error versus input magnitude. Although the scatter in the absolute value of the true errors is much larger than the scatter in the estimated errors, the estimated errors provide a good approximation to the true errors in an average sense. This remains true down to the completeness limit. In all three bands, the estimated errors fall along the curves fit to the true errors with significant deviation only below $`17.5`$ magnitudes.
### 4.2 Empirical and Combined Luminosity Functions
The luminosity functions (LFs) for each of the narrow band filters are shown in Fig. 5. No corrections for reddening or completeness have been made. The range in magnitude over which saturation occurs is indicated by the grey band in each panel. The vertical dotted lines indicate the mean saturation limit and the completeness limit of $`17.5`$ magnitudes. The F215N luminosity function is relatively flat between the saturation and completeness limits, with a dip between 14 and 15.5 magnitudes. The structure in the F166N and F190N luminosity functions is similar.
In order to compare our LF with previously determined LFs for IC348, we converted our measured F215N magnitudes to standard $`K`$ magnitudes. The F215N filter measures a relatively feature-free region of the standard CTIO/CIT $`K`$ filter. Therefore, the F215N magnitude should correlate well with $`K`$, requiring a zero-point offset and possibly, due to increasing water band strengths in the coolest M stars, a color term. To determine the offset, we compared our F215N magnitudes with published $`K`$ magnitudes for the 61 stars in our sample that are in common with Lada & Lada (1995; see tabulation in LRLL ) and/or Luhman (1999) and are below the saturation limit ($`K>11`$). The fit had a slope statistically identical to unity ($`1.003\pm 0.012`$), so we derived the mean offset between the two magnitude systems, $`(KF215N)=0.115\pm 0.011`$, where the error is the error in the mean. The 1–$`\sigma `$ residual to the fit was $`0.085`$. This residual is comparable to the typical combined photometric accuracy of the Luhman and our data. We investigated a possible color term in the transformation, but found that if any is present it is smaller than this scatter about the mean.
The accuracy and completeness of the bright end of our luminosity function ($`K11`$) is compromised by both saturation and our deliberate avoidance of bright cluster stars. To correct for this deficiency, we combined our derived $`K`$ photometry at $`K11`$ with $`K`$ photometry of the $`5\mathrm{}\times 5\mathrm{}`$ core from LRLL for $`K<11`$. This combination is reasonable since, as shown in Fig. 1, the region of our survey largely overlaps the $`5\mathrm{}\times 5\mathrm{}`$ core. To correct for the different areas covered by two surveys, we multiplied the counts in each bin of the LRLL luminosity function by the ratio of the survey areas, $`34.76/25.00=1.39`$.
The combined $`K`$ luminosity function for our 34.76 sq. arcmin region, complete to $`K17.5`$, is shown in Fig. 6 as the solid line histogram. To estimate the background contribution to the $`K`$ luminosity function, we used the prediction of the star count model of Cohen (1994). The predicted background $`K`$ counts, reddened by the mean reddening of the background population ($`A_K=0.71`$; see section 6), is shown as the dotted line histogram in Fig. 6. In section 8, we compare in greater detail the results for our data set with the predictions of the model. Here we simply note a few points. The contamination of the cluster by background stars is insignificant to $`K13`$ and the number of cluster stars is larger than the number of background stars until the $`K=14.25`$ bin. While the background rises steadily, we appear to have detected a few cluster stars to our completeness limit.
## 5 Spectral Classification
To derive spectral types for the stars in the sample, we combined the measured narrow band fluxes into a reddening independent index,
$$Q_{\mathrm{H2O}}2.5\mathrm{log}\left(\frac{\mathrm{F166}}{\mathrm{F190}}\right)+1.37\times 2.5\mathrm{log}\left(\frac{\mathrm{F190}}{\mathrm{F215}}\right),$$
that measures the strength of the $`1.9\mu `$m H<sub>2</sub>O absorption band. In this expression, F166, F190, and F215 are the fluxes in the F166N, F190N, and F215N filters, respectively. The value 1.37 is the ratio of the reddening color excesses:
$$\frac{E(\mathrm{F166}/\mathrm{F190})}{E(\mathrm{F190}/\mathrm{F215})}=1.37,$$
which is derived from the infrared extinction law $`A_\lambda /A_V=0.412(\lambda /\mu \mathrm{m})^{1.75}`$ (Tokunaga (1999)).
To explore the utility of the water index as an indicator of spectral type, we examined the relation between $`Q_{\mathrm{H2O}}`$ and spectral type for both the standard stars and a subset of IC348 stars that have optically determined spectral types from LRLL and Luhman (1999). For the standard stars, we adopted spectral types from the literature that are derived consistently from the classification scheme of Kirkpatrick et al. (1995). As shown in the top panel of Figure 7, $`Q_{\mathrm{H2O}}`$ is strongly correlated and varies rapidly with spectral type among the standard stars, confirming the expected sensitivity of the water band strength to stellar effective temperature. As is evident, there is real scatter among the standard stars that cannot be explained by errors in $`Q_{\mathrm{H2O}}`$ and spectral type. The scatter may reflect the inherent diversity in the standard star sample, a property that is evident from their $`JHK`$ colors. The spread in broad band color for a given spectral type is usually interpreted as the result of varying metallicity (e.g., Fig. 1 from Leggett et al. (1996)).
To compare these results with those for a population that has a more homogeneous metallicity distribution and the same mean metallicity and gravity to the IC348 sample, we also examined the $`Q_{\mathrm{H2O}}`$ vs. spectral type relation for the subset of IC348 stars that have optical spectral types determined by LRLL and Luhman (1999) (middle panel of Fig. 7). Although LRLL found no systematic difference between their IR and optical spectral types, there is significant dispersion between the two systems (their IR spectral types differ from the optical spectral types by as much as 3 subclasses). We find that the water band strengths are better correlated with the optical spectral types, with a smaller dispersion, than the IR spectral types, suggesting that their optical spectral types are more precise.
With the use of optical spectral types, we were also able to compare directly the results for the dwarf standards and the IC348 population, since both sets of objects are classified on the same system. The two samples exhibit a similar relation between spectral type and $`Q_{\mathrm{H2O}}`$ despite the difference in gravity between the two samples, with some evidence for a shallower slope for the pre-main sequence stars compared to the dwarfs. However, with the present data alone, we cannot claim such a difference with much certainty because the sample sizes are not large enough, the IC348 stars are not distributed evenly enough in spectral type, and there could be small systematic differences in the spectral typing of the IC348 and standard stars. The possibility of a difference between the two relations could be explored with more extensive optical spectral typing of the IC348 population.
The horizontal and vertical error bars in the lower left corner of the middle panel of Fig. 7 represent the typical errors in $`Q_{\mathrm{H2O}}`$ and spectral type for the cluster stars. Some of the scatter may arise from infrared excesses (which would uniquely affect the young star sample, compared to the standard star sample), although this effect is expected to be limited given the relatively small fraction of cluster sources that have IR excesses. For example, based on their $`JHK`$ photometry, Lada & Lada (1995) determined that $`<`$12% of sources brighter than $`K=14`$ in IC348a have substantial IR excesses. The LRLL study spectroscopically inferred $`K`$ continuum excesses in a similar fraction (15%) of sources in the subcluster.
To examine the possible impact of IR excess on our derived $`Q_{\mathrm{H2O}}`$ values, we considered excesses of the form $`\mathrm{\Delta }F_\nu (\lambda )\lambda ^\beta `$ and explored the effect of the excess on the $`Q_{\mathrm{H2O}}`$ values for two of our standards, the M3 dwarf Gl388, and the M6 dwarf Gl406. Since classical T Tauri stars have excesses at $`K`$ of $`r_K0.6`$ (Meyer et al. (1997); where $`r_K`$ is the ratio of the excess emission to the stellar flux), the IC348 sources, being more evolved, are likely to have much weaker excesses, typically $`r_K<0.2`$. With a spectral index of $`\beta =1/3,`$ appropriate for both disks undergoing active accretion and those experiencing passive reprocessing of stellar radiation, an IR excess produces an increase in $`Q_{\mathrm{H2O}}`$. Since the spectral slope is shallow and the maximum excess is small, only modest excursions are possible. For example, the $`Q_{\mathrm{H2O}}`$ index for Gl388 varies from its observed value, -0.28, at 0% excess to -0.24 at 20% excess in F215N. Over the same range of 0 to 20% excess in F215N, the $`Q_{\mathrm{H2O}}`$ index for Gl406 ranges from -0.52 to -0.43. This range of variation is sufficiently large that IR excess could account for most of the scatter of IC348 stars away from the mean trend to larger values of $`Q_{\mathrm{H2O}}`$. Explaining the scatter to smaller values of $`Q_{\mathrm{H2O}}`$ as the result of IR excesses requires more extreme values of $`\beta `$. For Gl388, values of $`\beta <3`$ are needed to decrease $`Q_{\mathrm{H2O}}`$ from its value at 0% excess. Such extreme spectral indices are unlikely as they would produce unusual broad band colors. For these reasons, it appears unlikely that IR excess is responsible for the majority of scatter about the mean relation between $`Q_{\mathrm{H2O}}`$ and spectral type. Other processes are implied, possibly including those that produce true differences in stellar water band strengths among stars with equivalent $`I`$-band spectral types.
Since we were not able to distinguish a systematic difference between the mean trends for the standard star sample and the IC348 sample, we used the combined samples to calibrate the relation between $`Q_{\mathrm{H2O}}`$ and spectral type (lower panel of Fig. 7). In order to use the error information in both $`Q_{\mathrm{H2O}}`$ and spectral type, we performed a linear fit in both senses (i.e., spectral type vs. $`Q_{\mathrm{H2O}}`$ and $`Q_{\mathrm{H2O}}`$ vs. spectral type; dotted lines in Fig. 7) and used the bisector of the two fits as the calibration relation (solid line in Fig. 7). Due to the non-uniform distribution of stars along the fit, the slope of the fit is sensitive to the inclusion or exclusion of stars near the sigma-clipping limit and at the extremes of either $`Q_{\mathrm{H2O}}`$ or spectral type. Doing a fit in both senses, and including the error information in both quantities, allowed us to better identify and exclude outliers. In the lower panel of Fig. 7, solid symbols indicate the stars that were included in the fit while open symbols indicate excluded stars.
The equation of the bisector, the relation we subsequently used to estimate spectral class for the entire cluster sample, is:
$$\mathrm{M}\mathrm{subtype}=1.09(\pm 0.39)13.01(\pm 0.50)\times Q_{\mathrm{H2O}}.$$
(5-1)
For a typical value of $`Q_{\mathrm{H2O}}`$, the formal spectral type uncertainty in the fit is $`\pm 0.46`$, while the scatter about the fit is $`0.81`$, just a little under one subtype. It is noteworthy that the discrete nature of spectral type versus the continuous nature of $`Q_{\mathrm{H2O}}`$ is responsible for a mean scatter of $`0.77`$ in $`Q_{\mathrm{H2O}}`$ in each subtype bin, which is a significant contribution to the total scatter.
Finally, we note that stars earlier than M2 have less certain spectral types due to the combination of the inherent scatter in the $`Q_{\mathrm{H2O}}`$ vs. spectral type relation and the decreasing sensitivity of the $`1.9\mu `$m H<sub>2</sub>O absorption band to spectral type as the K spectral types are approached. As a result, stars with spectral types of K and earlier can be misclassified by our method as later-type objects. For example, a comparison of the spectral types obtained by LRLL and Luhman (1999) with those obtained by our method shows that stars earlier than $``$K5 are classified by us as late K or M0 stars and late-K stars are classified as late-K and M0-M1 stars.
## 6 Extinction
Although the stellar spectral typing could be carried out without determining the reddening to each object, extinction corrections are required in order to investigate the masses and ages of cluster objects. We estimated the extinction toward each star by dereddening the observed F166/F190 and F190/F215 colors to a fiducial zero-reddening line in the color-color plane. Since extinction estimates for the dwarf standard stars were not available in the literature, we adopted the usual assumption that they suffer zero extinction. Figure 8 diagrams the process. First, we fit a line to the positions of the standard stars in the color-color plane (top panel), which is defined to be a locus of zero reddening. The WTTS were excluded from the fit. Gl569A was regarded as an outlier and also excluded from the fit. The resulting linear relation is:
$$2.5\mathrm{log}\left(\frac{\mathrm{F166}}{\mathrm{F190}}\right)=0.277(\pm 0.009)0.358(\pm 0.083)\times 2.5\mathrm{log}\left(\frac{\mathrm{F190}}{\mathrm{F215}}\right)$$
with a mean deviation about the fit of 1–$`\sigma =0.036`$. The extinction toward each star in the cluster fields was determined from the shift in each color required to deredden the star to the zero-reddening line. The resulting extinction estimates and errors are given in column 11 of Table From Stars to Super-planets: the Low-Mass IMF in the Young Cluster IC348 <sup>1</sup><sup>1</sup>1Based on observations with the NASA/ESA Hubble Space Telescope, obtained at the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc., under NASA contract No. NAS5-26555.. Note that the reddening vector (shown for $`A_V=10`$ in the bottom panel of Fig. 8), is nearly perpendicular to the standard star locus in the color-color plane. Consequently, reddening and spectral type are readily separable with moderate signal-to-noise photometry even given modest uncertainties in the slope of the reddening vector.
The subset of our standards used for the reddening calibration span the spectral class range K2V to M9V. This range is indicated by the dotted lines in the lower panel of Fig. 8. The few stars in the field with spectral types outside this range have extinction estimates based on the extrapolation of the fiducial line. As we show in Section 8, most of the early type stars are likely background objects. Finally, while the formal uncertainty in the fit of the fiducial line to the standards is small, 0.04 magnitudes, the scatter about the line for the latest standards is significantly larger than the scatter for the earlier standards (top panel of Fig. 8). Part of this scatter is due to the larger photometric errors; the late type standards are also the dimmest. However, four of the five late type standards fall above the fiducial line. In order from upper left to lower right these standards are LHS3003(M7V), Gl569B(M8.5V), VB10(M8V), LHS2924(M9V), and VB8(M7V). With the exception of VB8, these stars are aligned in the expected order in both colors but seem to be systematically shifted about 0.1 magnitudes to the red in $`2.5\mathrm{log}(\mathrm{F166}/\mathrm{F190})`$. While we cannot exclude the possibility that the relationship is non-linear for dwarfs later than M6, some of the scatter about the fit may be due to inherent variation in the photometric properties of the standard stars.
We can compare our extinction estimates to those of LRLL for the M dwarfs common to both samples. In Figure 9, the horizontal error bars indicate the formal (1–$`\sigma `$) uncertainty in our $`A_K`$ estimate (typically $`<0.1`$ mag). LRLL used various extinction estimators, citing their internal errors rather than values for individual stars. Their errors in $`A_K`$ range from 0.07 to 0.19 mag for the stars shown with a nearly equal systematic uncertainty in the zero point. For the M-dwarfs common to both samples, the mean difference in $`A_K`$, in the sense $`\mathrm{Ours}\mathrm{LRLL}=0.01\pm 0.03`$ with a scatter about the mean of 0.24 magnitudes. Considering the uncertainties, the agreement is good.
The resulting $`A_K`$ distribution (Fig. 10; solid-line histogram), has a pronounced tail to large values of $`A_K`$ and a peak at $`A_K=0.1`$. The extinction distribution for the subset of objects identified as the background population (as determined in section 8; dashed-line histogram) is also shown. Note that our extinction estimates include a few negative values (Figs. 8 and 9). While these values might suggest that our fiducial line needs to be lowered to bluer colors, that would imply a bias toward larger extinctions given the distribution of standard stars in the color-color plane. Therefore, we retain our original fit and, for all subsequent analysis, stars with negative extinction estimates are assigned an extinction of 0.0 with an error equal to the greater of the absolute value of the original extinction estimate or the formal uncertainty in the estimate.
With this revision, the mean extinction is $`A_K=0.44`$ with an error in the mean of 0.04 and a median of $`(A_K)_{\frac{1}{2}}=0.31`$. Our adjustment of the negative values impacts negligibly on the statistics. (If the negative extinction values were retained, the mean would be $`A_K=0.43`$ with the error and median unchanged.) When our sample is restricted to those stars in common with LRLL , we find approximately the same mean reddening ($`A_K=0.30`$) that they quote for their sample ($`A_K=0.34`$). The larger mean reddening in the present study indicates that, on average, we have sampled a more extincted population of the cluster than has been investigated previously. Using the position of the main sequence at the distance of the cluster (see section 8) to divide the sample into cluster and background objects, we find that the cluster objects have $`A_K=0.31\pm 0.04`$ with a scatter about the mean of 0.36. The background stars, which include most of the stars in the extended high extinction tail, have $`A_K=0.71\pm 0.07`$, with a scatter about the mean of 0.49.
The more heavily reddened stars in our sample are spatially intermixed with stars experiencing lower extinction. Figure 11 shows the same area plotted in Fig. 1. The gray symbols denote stars in our sample that were observed by other investigators (LRLL ; Herbig (1998); Luhman 1999), whereas the black symbols denote stars that were not observed by these investigators. The point size is scaled to our estimate of the extinction to the object (larger points corresponding to larger reddening), which ranges from $`A_K=0.0`$ to $`A_K=2.33`$. The higher average extinction among the black points is apparent. The extinction distribution is characterized by an overall gradient from NE (larger values) to SW (smaller) with significant small scale variation. Given the broad extinction distributions for both cluster and background objects, and the patchy distribution of extinction on the sky, it is evident that cluster membership cannot be determined on the basis of extinction alone. Membership based on extinction would erroneously assign low extinction background members to the cluster and highly extincted cluster members to the background.
## 7 Observational HR Diagram
With the spectral types determined in section 5 and the extinction extimates from section 6, we can construct an observational HR diagram of the cluster fields. In Figure 12, the vertical axes are apparent $`K`$ magnitude (left panel) and dereddened $`K`$ magnitude, $`K_0`$ (right panel). For comparison, the solid curve in the right panel is the fiducial main sequence at the distance of the cluster (see section 8.2). Examination of both panels reveals a well defined cluster sequence at $`K14`$. This locus is marginally tighter after being dereddened which supports the accuracy of our reddening estimates.
Spectral type errors are not shown, both to limit confusion and because some stars have systematic as well as random error. For example, although the typical random error is $`\pm 1`$ spectral subtype, stars earlier than M2 have systematically later $`Q_{\mathrm{H2O}}`$ spectral types than optical spectral types (section 5). Given the possible inaccuracy of our spectral typing scheme at spectral types earlier than M2, we adopted the optical spectral types of Luhman (1999) or LRLL for these objects where available. The original $`Q_{\mathrm{H2O}}`$ spectral types of these stars are shown as open circles in Figure 12. When optical spectral types of these stars are adopted instead (see subsequent figures), the photometric width of the distribution at M2 and earlier is reduced. In general, the random error in spectral type increases with increasing magnitude (see column 13 of Table From Stars to Super-planets: the Low-Mass IMF in the Young Cluster IC348 <sup>1</sup><sup>1</sup>1Based on observations with the NASA/ESA Hubble Space Telescope, obtained at the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc., under NASA contract No. NAS5-26555.). All stars with $`K<15.5`$ have spectral type errors $`1`$ subtype. Since our spectral type errors grow rapidly below $`K=16`$, with stars fainter than $`K=16.5`$ having spectral type errors $`2.5`$ subtypes, we identify $`K=16.5`$ as our effective magnitude limit for accurate spectral typing.
While some objects have spectral types as late as “M13”, this should be interpreted simply as an indication of strong water absorption rather than an advocacy of M spectral types beyond M9. The existence of objects with stronger water absorption than that of M9 dwarfs is in general agreement with the predictions of atmospheric models (e.g., the Ames-Dusty and Ames-MT-Dusty models of Allard et al. (2000)). These suggest that even in the presence of dust, the $`1.9\mu `$m H<sub>2</sub>O absorption band continues to increase in strength down to $`2000`$ K at pre-main sequence gravities. In the Ames-Dusty models, $`Q_{\mathrm{H2O}}`$ increases in strength by 45% between 2450K (equivalent to M8 in the dwarf temperature scale; see section 8.2) and 2000K. The $`Q_{\mathrm{H2O}}`$ vs. spectral type relation in eq. 5-1 implies that $`Q_{\mathrm{H2O}}`$ is 54% stronger at M13 than at M8, in general agreement with the predictions.
The dearth of stars at $`K15.5`$ in the $`K`$ luminosity function is also evident in the left panel of Fig. 12. Part of the deficit is due to the higher average reddening of the background stars. Stars with $`K>15.5`$ have an average extinction greater than stars with $`K<15.5`$ and when they are dereddened, they fill in the deficit somewhat. Our photometric completeness limit of $`K=17.5`$ is shown in the left panel as a horizontal dotted line. To quantify our detection limit as a function of extinction, we also show the completeness limit dereddened by $`A_K=0.31`$, the mean extinction among the cluster stars (lower horizontal dotted line in the right panel of Fig. 12) and by $`A_K=2.33`$, the greatest extinction detected in the cluster fields (upper horizontal dotted line). Both limits, $`K=17.19`$ and $`K=15.17`$, are considerably dimmer than the typical cluster M star.
These results imply that we have fully sampled the cluster population over a significant range in extinction. The extinction range that we probe is, of course, a function of spectral type. As examples, of the two cluster stars in the tail of the reddening distribution shown in Fig. 10, one is an M2 star with $`K=12.16`$ ($`A_K=1.97`$) and the other is an M9 star with $`K=16.73`$ ($`A_K=1.52`$). We would have been able to detect and spectral type the first star through another $`4.4`$ mag of extinction (to $`A_K6.4`$). The second star, observed through almost 5 times the average cluster extinction, is close to our spectral typing limit.
## 8 Comparison with Evolutionary Tracks
### 8.1 Evolutionary Models
Evolutionary models for low mass objects have developed greatly in recent years, with several different models now available over a large range in mass. D’Antona & Mazzitelli (1997) have recently updated their pre-main sequence calculations, retaining the use of the Full Spectrum Turbulence model of Canuto & Mazzitelli (1991) and making improvements in opacities and the equation of state. For the purpose of this paper, we use their 1998 models <sup>2</sup><sup>2</sup>2These models are available at: http://www.mporzio.astro.it/ dantona (hereinafter DM98) which cover the mass range $`0.0170.3M_{}`$ and include further improvements, e.g., in the treatment of deuterium burning, that affect the very low mass tracks.
Other groups (e.g., Baraffe et al. (1998); Burrows et al. 1997) have also presented new evolutionary models that include improvements in the treatment of the stellar interior and use non-gray atmospheres as an outer boundary condition. The corrections associated with the latter are particularly significant at low masses since the presence of molecules in low temperature atmospheres results in spectra that are significantly non-blackbody. Models by Baraffe et al. (1998; hereinafter B98) explore the mass range $`0.0251.0M_{}`$ using the Allard et al. (1997) NextGen synthetic atmospheres. Although there are known inconsistencies in the NextGen models (e.g., they overpredict the strength of the IR water bands; TiO opacities are suspected to be incomplete; grain formation is not included), the B98 models nevertheless reproduce well the main sequence properties of low metallicity populations, e.g., the optical color-magnitude diagram of globular clusters and halo field subdwarfs. There is also good agreement with the optical and IR properties of nearby disk populations, although some discrepancies remain at low masses ($`<0.15M_{}`$).
Non-gray models have been developed independently by Burrows et al. (1997) who focus on the properties of objects at lower mass ($`0.370M_J`$, where $`M_J`$ is the mass of Jupiter). The Burrows et al. evolutionary tracks differ qualitatively from those of B98 in the upper mass range, but are more qualititatively similar at masses $`60M_J.`$ The qualitative difference between these models, which appear to have similar input physics, may indicate the current level of uncertainty in the evolutionary tracks at low masses. Quantitatively, an effective temperature of 3340 K and luminosity of $`0.076L_{}`$ corresponds to a mass and age of $`0.090M_{}`$ and 1.8 Myr with the Burrows et al. tracks and $`0.3M_{}`$ and 8 Myr with the B98 tracks. The tracks agree better in mass in the lower mass range: at 2890 K and $`0.022L_{}`$, Burrows et al. predict $`0.05M_{}`$ at 1.2 Myr, and the B98 tracks predict $`0.06M_{}`$ at 3.2 Myr.
### 8.2 Interloper Population
As reviewed by Herbig (1998), the distance to IC348 has been previously estimated on the basis of both nearby stars in the Per OB2 association and stars in the IC348 cluster itself. For the purpose of comparing our results with evolutionary tracks, we adopt a distance to IC348 of $`d=300`$ pc, $`(mM)_0=7.4`$. This value is in good agreement with current estimates of the distances to the Per OB2 cluster (318$`\pm 27`$ pc; de Zeeuw et al. 1999) and to IC348 itself (261$`\pm 25`$ pc; Scholz et al. 1999) inferred from Hipparcos data. The adopted distance is also in agreement with the value adopted by both Herbig (1998) and LRLL and thereby allows ready comparison of our results with those obtained in previous studies.
To delineate the background population, the position of the main sequence at the cluster distance is indicated by the solid curve in the right panel of Fig. 12, where we have used the 12 Gyr isochrone from the B98 evolutionary tracks and a temperature scale that places the isochrone in good agreement with the main sequence locus of nearby field stars (e.g., Kirkpatrick & McCarthy (1994)). The temperature scale used,
$$\mathrm{M}\mathrm{subtype}=(4000T_{\mathrm{eff}})/180,$$
is generally consistent with the Leggett et al. (1996) dwarf temperature scale.
The magnitude and spectral type distributions of the background population, located to the lower left of the main sequence, are in very good agreement with the total interloper population predicted by models of the point source infrared sky (Wainscoat et al. (1992); Cohen 1994) at the Galactic latitude and longitude of IC348. Table 4 compares the observed and model counts as a function of $`K`$ magnitude and spectral type. To $`K_0=17`$, significant departures between the model and observed counts are apparent only for spectral types earlier than M3 at $`K_0>16`$. Given the large spread in the reddening distribution of the background population (to $`A_K>2`$; Fig. 10), this discrepancy in the counts probably arises from photometric incompleteness below $`K=17.5`$. This result (the good agreement between the model prediction for the total interloper population and the observed background population), implies a negligible foreground contamination (at most $`12`$ stars) of the cluster population at late spectral types. The large reddening of many of the faint late-type stars also statistically argues against a foreground origin for these objects. Note, however, that the errors on some of the fainter objects identified as older cluster members (e.g., objects in the range $`K_0=15.516.5`$, M6$``$M8) allow for the possibility that they are background objects even if they are not predicted to be so by the Galactic structure model.
### 8.3 Temperature Scale and Bolometric Correction
A generic difficulty in comparing measured stellar fluxes and spectral types with evolutionary tracks is the need to adopt relations between spectral type, effective temperature, and bolometric correction. In principle, such relations could be avoided by using synthetic spectra from model atmospheres to go directly from observed spectra and colors to temperature and gravity, and hence to mass and age using the theoretical evolutionary tracks. For example, we might hope to compare directly the water band strengths of the Allard & Hauschildt atmospheres used in the B98 models with the water band strengths that we measured. However, since there remain significant quantitative differences between the predicted and observed water band strengths of M stars (e.g., the models consistently overpredict water band strengths; see also Tiede et al. (2000)), this approach cannot be used in the present case. In other words, although current synthetic atmospheres may be sufficiently accurate for the purpose of evolutionary calculations and the prediction of broad band colors, they are insufficiently accurate as templates for spectral typing. Hence, we adopted the less direct method of first calibrating our water index versus spectra type (section 5), and then selecting an appropriate spectral type to temperature conversion.
Ideally, we would want to use a relation between spectral type and effective temperature that is appropriate to the gravity and metallicity of the IC348 population. Unfortunately, an empirical calibration of spectral type and effective temperature appropriate for pre-main-sequence conditions has yet to be made. In the meantime, since pre-main-sequence gravities are similar to dwarf gravities, temperature scales close to the dwarf scale (e.g., Leggett et al. (1996)) are often used in the study of young populations (e.g., LRLL; Wilking et al. (1999)). Because the temperature scale may differ from that of dwarfs at PMS gravities, other choices have also been investigated, including temperature scales intermediate between those of dwarfs and giants (e.g., White et al. (1999); Luhman 1999).
The validity of the various evolutionary tracks can be evaluated by a number of criteria including whether stellar masses predicted by evolutionary tracks agree with dynamical estimates, and whether populations believed to be coeval appear so when compared with evolutionary tracks (e.g., Stauffer et al. 1995). Dynamical mass constraints are becoming available in the $`1M_{}`$ range (see, e.g., Mathieu et al. (2000)) but are thus far unavailable at the masses of interest in the present study. In contrast, coeval population constraints are more readily available at these lower masses. For example, in the GG Tau hierarchical quadruple system (White et al. (1999)), the four components of the system, arguably coeval, span a wide range in spectral type (K7 to M7; open squares in Fig. 13, upper left), thereby outlining, in rough form, an isochrone spanning a large mass range. When plotted at a common distance, the IC348 cluster locus identified in the present study overlaps the locus defined by the GG Tau components over the same range of spectral types (Fig. 13). This both reinforces the validity of the GG Tau system as a coeval population constraint and argues that the mean age of the IC348 cluster is approximately independent of mass. Similar results have been found previously at spectral types earlier than M6 (Luhman 1999).
The uncertainty in the pre-main-sequence temperature scale complicates our understanding of the validity of the tracks. As discussed by Luhman (1999), combinations of evolutionary tracks and temperature scales that are consistent with a coeval nature for the GG Tau system and the IC348 cluster locus include (1) DM98 tracks and a dwarf temperature scale (2) B98 tracks and an otherwise arbitrary temperature scale intermediate between that of dwarfs and giants. Our results are compared in Figure 13 with these combinations of temperature scales and tracks. For comparison, the two alternative combinations of temperature scales and tracks are shown. In comparing the B98 models with the observations, we have used the model $`K`$ magnitudes and a linear fit to either the dwarf temperature scale
$$\mathrm{M}\mathrm{subtype}=(3914T_{\mathrm{eff}})/183.3$$
(8-1)
or the Luhman (1999) intermediate temperature scale
$$\mathrm{M}\mathrm{subtype}=(3850T_{\mathrm{eff}})/141.0$$
(8-2)
In approximating the dwarf temperature scale, particular weight was given to the dwarf temperature determinations by Tsuji et al. (1996) who used the IR flux measurement technique. As they show, this technique is relatively insensitive to the details of synthetic atmospheres (e.g., dust formation). The fit thus obtained is in good agreement with the temperature determinations of Leggett et al. (1996) which are based on a comparison of synthetic atmospheres with measured IR colors and spectra. In comparing the DM98 models with the observations, we have used, in addition to these temperature scales, a bolometric correction
$$\mathrm{BC}_\mathrm{K}=M_{\mathrm{bol}}M_K=4.19T_{\mathrm{eff}}/2240$$
that extrapolates the values obtained by Leggett et al. (1996) and Tinney et al. (1993) to low temperatures.
The combination of the B98 models and the Luhman intermediate temperature scale (eq. 8-2; Fig. 13 upper left) implies that the mean age of the cluster is approximately independent of mass over the range $`0.70.04M_{}`$. The comparison implies a mean age $``$3 Myr with a age spread from $`<1`$ to $`20`$ Myr. The faint cluster population between spectral types M5 and M8 appears to constitute an old cluster population ($`5`$ to $`>20`$ Myr) with masses $`0.130.05M_{}`$. If the dwarf temperature scale (eq. 8-1; Fig. 13 upper right) is used instead, the cluster is, on average, significantly younger at late spectral types.
The combination of the DM98 models and the dwarf temperature scale (eq. 8-1; Fig. 13 lower right) implies that the mean cluster age is approximately independent of mass at spectral types earlier than M7 but younger at late types. The comparison implies a mean age $`1`$ Myr with a age spread from $`<1`$ to $`10`$ Myr. With these models, the faint cluster population between spectral types M5 and M8 is spread over a larger range in mass $`0.160.025M_{}`$. If the Luhman intermediate temperature scale (eq. 8-2; Fig. 13 lower left) is used instead, the cluster is older at late types with a larger spread in age. With all combinations of models and temperature scales, the brighter cluster population beyond M8 is systematically younger, $`<1`$ Myr old. If this is an artifact, it may indicate the likely inadequacy of the assumed linear relation between effective temperature and spectral type over the entire range of spectral types in the sample. Deficiencies in the evolutionary tracks are another possibility.
It is interesting to examine the motivation for the intermediate temperature scale adopted by White et al. (1999) and Luhman (1999). These authors have argued that since the M giant temperature scale is warmer than the dwarf scale, PMS stars, which are intermediate in gravity, may be characterized by a temperature scale intermediate between that of giants and dwarfs. Luhman (1999) has further shown that the spectra of pre-main sequence stars in IC348 are better fit by an average of dwarf and giant spectra of the same spectral type.
There are several caveats to this argument. Firstly, the giant temperature scale considered by Luhman (1999) is derived from the direct measurement of stellar angular diameters (e.g., Perrin et al. (1998); Richichi et al. (1998); van Belle et al. 1999 ), whereas the dwarf temperature scale is typically determined with the use of model spectra (e.g., Leggett et al. (1996); Jones et al. (1994); Jones et al. (1996)). The different methods by which the two temperature scales are derived may introduce systematic differences that do not reflect a true temperature difference.
Secondly, we can turn to synthetic atmospheres for insight into the gravity-dependent behavior of the temperature scale. In the current generation of the Allard & Hauschildt atmospheres (e.g., Ames-Dusty, Ames-MT-Dusty), the 1.9$`\mu `$m water band strength is relatively insensitive to gravity above 3000K ($``$M5 in the dwarf scale). At effective temperatures below 3000K, dust formation is significant, introducing added complexity to the gravity dependence of the atmosphere in the 1.9$`\mu `$m region. In this temperature range, the water index first increases in strength ($`Q_{\mathrm{H2O}}`$ decreases) at fixed temperature from $`\mathrm{log}g3.5`$ to $`\mathrm{log}g5.05.5`$ (due to increased water abundance) then decreases in strength with higher gravity (due to increased dust formation and consequent backwarming and dissociation of water). The net result is a cooler temperature scale for pre-main-sequence gravities below 3000K. For example, at $`2700`$ K pre-main-sequence objects ($`\mathrm{log}g=3.54.0`$) are $`200`$K cooler than dwarfs ($`\mathrm{log}g=5.05.5`$) with an equivalent water strength.
On the basis of these models, there is little physical motivation for an intermediate temperature scale beyond M4 for the interpretation of water band strengths. Of course, these considerations apply to the interpretation of 1.9$`\mu `$m water band strengths rather than the 6500-9000Å region studied by Luhman (1999). A detailed examination of current synthetic atmospheres for the latter spectral region may provide better motivation for a hotter temperature scale at lower gravities.
Note that the gravity dependence of $`Q_{\mathrm{H2O}}`$ in the synthetic atmospheres is modest over the range of gravities relevant to low mass pre-main sequence stars in the age range of the cluster (1–10 Myr). For example, in the B98 model, an $`0.06M_{}`$ object follows a vertical evolutionary track at $`T_{\mathrm{eff}}2860`$K with $`\mathrm{log}g=3.64.2`$ in the age interval 1–10 Myr which corresponds to a fractional change in $`Q_{\mathrm{H2O}}`$ of $`15`$% or $`1`$ subtype, given the relation between $`Q_{\mathrm{H2O}}`$ and spectral type discussed in section 5.
In summary, while we can find little physical motivation for an intermediate temperature scale with which to interpret our results, we interpret the better fit to the IC348 cluster locus that we obtain with the combination of this temperature scale and the B98 models as an indication of the direction in which the evolutionary model calculations might themselves evolve in order to better reproduce observations of young clusters. With these caveats in mind, we discuss, in the next section, the cluster mass function implied by 2 combinations of tracks and temperature scales. However, it is already clear that there will be reasonable uncertainty associated with such results.
## 9 Discussion
### 9.1 Binarity
The area and depth that we have covered at relatively high angular resolution, combined with our ability to discriminate cluster members from background objects, allows us to place some useful constraints on the binary star population of the cluster. At the pixel scale of NIC3, pairs of stars with separations $`0.8\mathrm{}`$ are easily identified over the entire magnitude range of our sample; for fainter primaries, companions could be similarly detected at smaller separations. A significant obstacle to the detection of faint companions at separations $`0.8\mathrm{}`$ is the complex, extended structure in the NICMOS PSF which also makes it difficult to quantify our detection completeness. More refined techniques, such as PSF subtraction or deconvolution, when applied to the data, are likely to reveal close binary systems that we have missed.
Table 5 tabulates all of the stars in our sample that were found to have a nearest neighbor within $`8\mathrm{}`$. The stars have been designated primary and secondary based on their $`K_0`$ magnitudes. The spectral types for the G dwarfs are from LRLL and the other spectral types are our spectral types as determined in Section 5. Figure 14 shows the positions of the close pairs in the observational HR diagram. To identify the pairs, the components are connected by lines. Although we were sensitive to separations $`0.8\mathrm{}`$, only pairs with separations $`>1.5\mathrm{}`$ were detected. Based on their locations in the observational HR diagram, seven of the close pairs are chance projections of a background star close to a cluster member (Fig. 14; dotted lines). Both components of one pair are background objects. Of the 8 candidate cluster binaries 3 (093-04/093-05; 043-02/043-03; 024-05/024-06) were previously detected by Duchene et al. (1999) in their study of binarity among a sample of 67 IC348 objects. We also confirm their speculation that 083-03 and 023-03 are background objects with small projected separations to cluster members.
As shown in Fig. 14 (solid lines), several of the candidate binary pairs have spectral types and $`K_0`$ magnitudes consistent with a common age for the two components. For the candidate binaries E, F, C, D, and H, the lines connecting the two components have slopes consistent with the isochrones. The candidate binary B has a nearly vertical slope. However, given our estimate of the uncertainty in the spectral types of the binary components, the slope is also highly uncertain, and a common age for the binary components cannot be ruled out. While the component spectral types for the binary candidate G have similar uncertainties, the large separation in magnitude between the two components, if each are single stars, makes it unlikely that they share a common age. If, on the other hand, the brighter component is an approximate equal mass binary, the reduced brightness of each of the two stars is more consistent with the evolutionary models, and the triple system may be coeval. If more definitive studies reveal that the binary candidates B and G are not coeval, this may indicate that they are not physically related. Alternatively, a large age difference between the components may indicate that the binaries formed through capture.
If we define the binary fraction as the ratio of the number of companions detected to the number of targets observed (193 stars), the cluster binary fraction in the separation range $`0.8\mathrm{}8\mathrm{}`$ (240 $``$ 2400 AU) is 8%. This is comparable to the result of Duchene et al. (1999) who, based on a smaller sample of stars, found a 19% binary fraction for their entire sample; half of their binaries fall in the separation range of our study. However, there are several important differences between the two studies. We sample a lower range of primary masses ($`0.0150.8M_{}`$) than Duchene et al. (1999) ($`0.22M_{}`$). In addition, the mass ratios to which we are sensitive are set by the magnitude limit of the sample rather than by the magnitude difference between the binary components. In contrast to Duchene et al. (1999), who commented on the lack of substellar companions, we find candidate substellar companions (e.g., 022-05) and one candidate substellar binary (H).
### 9.2 Low-Mass Cluster Members
The very low-mass cluster population is highlighted in Figure 15. The 6 objects indicated have the largest water absorption strengths in the sample, corresponding to spectral types later than M9, and presumably the lowest masses. The errors on the derived properties for 3 of the objects (012-02, 102-01, 022-09) are modest, and imply masses $`0.025M_{}`$ in the context of both the B98 and DM98 models. The other 3 objects (024-02, 075-01, and 021-05) are in fact fainter than our effective limit for accurate spectral typing ($`K=16.5`$) and so have spectral type errors $`>2.5`$ subtypes (cf. section 7). Two of these objects, 024-02 ($`A_K=1.52`$) and 075-01 ($`A_K=2.3`$), are faint due to their large extinctions and are $`5`$ and $`7`$ times more extincted, respectively, than the cluster mean. Even with the larger errors for these objects, it appears very likely that all 3 are substellar cluster objects. However, because of its proximity to the main sequence, there is a small probability that 021-05 is a background M star.
### 9.3 Mass Function
To estimate a mass function for our sample, we used two combinations of evolutionary models and temperature scales: the B98 models in combination with the Luhman (1999) intermediate temperature scale and the DM98 models in combination with the dwarf temperature scale. The lower mass limit to which we are complete is determined by our spectral typing limit. As discussed in section 7, we have fairly accurate spectral types for all sources to $`K=16.5`$. For a mean cluster reddening of $`A_K0.3,`$ this corresponds to $`K_016.2`$ or $`M_K8.8`$ at the assumed distance of IC348. Thus, with the DM98 models, we are, for example, complete to $`0.017M_{}`$ at the mean extinction of the cluster and ages $`<3`$ Myr.
For the B98 models, some extrapolation was needed to both younger ages ($`<2`$ Myr), in order to account for the brighter cluster population, and to lower masses ($`<0.025M_{}`$) in order to estimate our mass completeness limit. In extrapolating below 2 Myr, we used the 1 Myr isochrone from the Baraffe et al. (1997) models as a guide. For the lower masses, we used the planetary/brown dwarf evolutionary theory of Burrows et al. (1997) to extrapolate the isochrone appropriate to the mean age of the subcluster (3 Myr). Several similarities between the Burrows et al. and B98 models suggest the utility of such an approach. Like B98, the Burrows et al. theory is non-gray, and the evolutionary tracks in the luminosity vs. $`T_{\mathrm{eff}}`$ plane at masses $`<0.04M_{}`$ are qualitatively similar. Two possible extrapolations are given to illustrate the uncertainty in the result.
In the B98 models, a $`0.025M_{}`$ object at 3 Myr has $`T_{\mathrm{eff}}=2628`$K, and $`M_K=7.56`$. In comparison, in the Burrows et al. theory, a $`0.025M_{}`$ object at 3 Myr is slightly hotter ($`T_{\mathrm{eff}}`$=2735K) but has a comparable absolute $`K`$ magnitude ($`M_K=7.6`$ assuming BC$`{}_{K}{}^{}=3.0`$); a 3 Myr old object that is 1.2 magnitudes fainter ($`M_K=8.8`$) has an effective temperature $`300`$K cooler and is $`0.011M_{}`$ lower in mass. Applying the same mass and temperature differentials to the 3 Myr old, $`0.025M_{}`$ object from B98 implies that a 3 Myr old, $`M_K=8.8`$ object in the B98 theory has $`T_{\mathrm{eff}}=2330`$K and a mass of $`0.014M_{}`$.
As an alternate estimate, we can extrapolate the 3 Myr isochrone based on a match in $`T_{\mathrm{eff}}`$ rather than mass. As described above, the effective temperature of a $`0.025M_{}`$, 3 Myr old object in B98 theory is 2628K. From the $`T_{\mathrm{eff}}=2628`$K point in the 3 Myr isochrone of the Burrows et al. models, $`\mathrm{\Delta }M_K=1.2`$ corresponds to a change in temperature and mass of $`\mathrm{\Delta }T_{\mathrm{eff}}=425`$K and $`\mathrm{\Delta }m=0.010M_{}`$. Applying these mass and temperature differentials to the 3 Myr old, $`0.025M_{}`$ object from B98 implies an effective temperature of $`2200`$K and mass $`0.015M_{}`$ for a 3 Myr old, $`M_K=8.8`$ object. Thus, with either estimate, our spectral typing limit of $`M_K=8.8`$ corresponds to a mass completeness limit of $`0.015M_{}`$ at the average age and reddening of the cluster members. The effective temperature appropriate to this mass limit is less certain.
Formally, the appropriate effective temperature affects our estimate of the lower limit to the final mass bin of our sample. Note, however, that our spectral typing limit of $`0.015M_{}`$ is close to the deuterium burning limit (Burrows et al. 1993, Saumon et al. 1996) and in the age range in which objects fade fairly rapidly with age. For example, in the Burrows et al. models, a 3 Myr old, $`0.010M_{}`$ object is half as luminous as a $`0.015M_{}`$ object at the same age. Given the rapid fading, it is unlikely that we have detected objects much less massive than $`0.015M_{}`$, which we adopt as the lower limit of the final mass bin of the sample. Note that our spectral typing limit of $`K=16.5`$ implies that we have somewhat underestimated the population of the final mass bin if that bin is characterized by the same spread in age and reddening that is measured at higher masses.
The mass functions for the age range $`010`$ Myr that result from the assumptions and extrapolations discussed above are shown in Fig. 16. The result for both the B98 models (solid symbols) and DM98 models (dotted symbols) are shown. Note that the objects indicated previously as potential background objects ($`K_0=15.516.5`$, M6$``$M8) are not included in the mass function for the B98 models, whereas some are included in the mass function for the DM models. Since these objects represent only a small fraction of the objects in each bin, whether or not these are included as members makes little difference to the slope of the mass function.
The DM98 models indicate a flattening at $`0.25M_{},`$ whereas the B98 models imply an approximately constant slope over the entire mass range $`0.70.015M_{}.`$ In either case, the mass function appears to decrease from $`0.25M_{}`$, through the hydrogen burning limit ($`0.08M_{}`$), down to the deuterium burning limit ($`0.015M_{}`$). The slope of the mass function in this range is consistent with $`dN/d\mathrm{log}MM^{0.5}`$ for B98 and $`dN/d\mathrm{log}MM^{0.6}`$ for DM98. The slow, approximately continuous decrease in the mass function in this interval differs from the result obtained by Hillenbrand (1997) for the Orion Nebula Cluster. The sharp fall off in the Orion Nebula Cluster mass function below $`0.2M_{}`$ ($`dN/d\mathrm{log}MM^{2.5}`$) is not reproduced here. Instead, we find that the slope of the IC348 mass function is more similar to that derived for the Pleiades in the mass range $`0.30.04M_{},`$ $`dN/d\mathrm{log}MM^{0.4}`$ (Bouvier et al. (1998)). The slope is similar to that inferred for the substellar population of the solar neighborhood from 2MASS and DENIS data. As determined by Reid et al. (1999), the observed properties of the local L dwarf population are consistent with a mass function $`dN/d\mathrm{log}MM^\alpha ,`$ with $`\alpha `$ -1 to 0, although a mass function similar to that for IC348 is not strongly precluded especially given the uncertainty in the age distribution of objects in the solar neighborhood.
Given the low masses to which we are sensitive, it is also interesting to compare our result to the mass function that is emerging for companions to nearby solar-type (G$``$KV) stars at separations $`<5`$ AU (e.g., Marcy et al. (2000)). While initial results indicated that the substellar companion mass function might be a smooth continuation of the stellar companion mass function (e.g., $`dN/d\mathrm{log}MM^{0.6}`$; Mayor et al. (1998)), proper motion data from Hipparcos have revealed that a significant fraction of companions in the $`0.0150.08M_{}`$ range are low inclination systems, and hence have larger (stellar or near-stellar) masses (Marcy et al. (2000); Halbwachs et al. 2000). When corrected for these low inclination systems, the companion mass function appears to be characterized by a marked deficit in the $`0.0150.08M_{}`$ mass range (the “brown dwarf desert”; Marcy et al. (2000); Halbwachs et al. 2000). In contrast, the mass function for IC348 appears to decrease continuously through the stellar/substellar boundary and the mass range $`0.080.015M_{}`$.
The low mass end of the IC348 sample extends into the mass range ($`1020M_J`$) in which objects transition from higher mass objects that burn deuterium early in their evolution to lower mass objects that are incapable of deuterium burning due to the onset of electron degeneracy pressure during the contraction phase (e.g., Grossman et al. (1974); Burrows et al. (1993)). According to the calculations of Saumon et al. (1996), $`15M_J`$ objects deplete their deuterium abundances by a factor of 2 after 30 Myr of evolution, while objects $`12M_J`$ retain all of their initial deuterium and derive no luminosity from thermonuclear fusion at any point in their evolution. They suggest the deuterium burning limit as a possible interpretive boundary between objects that are regarded as brown dwarfs and those regarded as planets.
If we assume that the hydrogen and deuterium burning mass limits delimit the brown dwarf population, with either the DM98 or B98 models, we have fully sampled the brown dwarf population, at ages up to the mean age of the subcluster and extinctions up to the cluster average. Thus, we can conclude with near certainty that the fraction of the subcluster mass contributed by brown dwarfs is low, only a few percent of the cluster mass. With the B98 tracks, we find a total of $`22`$ cluster substellar candidates which represents a significant fraction, $`20`$%, of all cluster M dwarfs by number, but only a small fraction, $`4`$%, by mass. For comparison, with the DM98 tracks, we find $`30`$ cluster substellar candidates which represents $`30`$% of all cluster M dwarfs by number and $`6`$% by mass.
These limits on the substellar contribution to the total cluster mass have interesting implications when compared with current limits placed by microlensing studies on the substellar content of the Galactic halo. Based on the search for microlensing toward the LMC the current EROS limits on the fraction of the halo mass that resides in brown dwarf mass objects is $``$ 10% (Lasserre et al. 2000). Scaling our results for IC348 by the stellar fraction of the halo mass ($`1`$ %), we find that if the halo has the same IMF as IC348, then substellar objects contribute negligibly to the halo mass ($`<0.1`$%). The several orders of magnitude difference between these limits leaves room for some interesting possibilities. If future microlensing results find confirmation for a halo mass fraction of even $`1`$ % in substellar objects, that would indicate that low mass star formation in the halo proceeded significantly differently from that currently occuring in Galactic clusters.
What do the IC348 results tell us about the star formation process? The absence of structure in the mass function at the hydrogen burning limit (e.g., a turnover) is perhaps expected. It is difficult to imagine how hydrogen burning, which demarcates the end of the pre-main sequence phase, could influence the determination of stellar masses, an outcome which is probably determined at much earlier times.
We also find no obvious feature in the IMF at the deuterium burning limit (e.g., a strong increase or decrease), a potentially more relevant mass scale for star formation since deuterium burning occurs at pre-main-sequence ages. This result may appear puzzling in the context of some current theories for the origin of stellar masses. For example, in a canonical theory of the formation of solar-type stars, it is the onset of deuterium burning that is believed to set in motion the sequence of events by which a star comes to have a role in determining its own mass. The onset of deuterium burning first induces a fully convective stellar interior. The convective interior, combined with the rapid stellar rotation that is likely to result from the accretion of angular momentum along with mass, is believed to generate a strong stellar magnetic field. The strong field is, in turn, believed to drive a magnetocentrifugal wind that ultimately sweeps away the cloud from which the star formed and possibly reverses the infall itself, thereby helping to limit the mass of the star. The self-deterministic aspect of such a mass-limiting wind is a critical element in explanations for the generic origin of stellar masses (e.g., Shu (1995)) and some theories of the IMF (e.g., Adams & Fatuzzo (1996)).
In this picture, as masses close to the deuterium burning limit are approached, one might expect that, the deuterium burning trigger being absent, low mass objects might not be able to reverse the infall and, consequently, it would be difficult to produce any objects of such low mass. This appears to be inaccurate both theoretically and observationally. Not only are young objects in this mass range fully convective without the aid of deuterium burning (Burrows, personal communication) and may thereby generate magnetic fields in advance of or in the absence of deuterium burning, but we also find no deficit of objects near the deuterium burning limit. This nevertheless raises the important question of what physical processes determine the masses of objects much below a solar mass.
Fragmentation is one possibly significant process at this mass scale. Coindentally, our survey mass limit is close to the characteristic mass for opacity-limited fragmentation under the low temperature, chemically enriched conditions current prevailing in molecular clouds ($`0.01M_{}`$; e.g., Silk 1977). In this picture, if cooling is efficient as collapse proceeds, the inverse dependence of the Jeans mass on density leads to fragmentation on increasingly small scales as collapse continues, halting only when objects become optically thick to their own radiation and the cooling efficiency is thereby impaired. If the characteristically low mass objects that form as the result of this process represent the “seeds” of star formation from which more massive objects must grow, we might expect to find a large number of objects with this mass. Perhaps significantly, we find no such large excess, but rather a smooth continuation from the stellar mass regime down to this mass scale. This implies that if fragmentation plays an important role in the formation of stars and brown dwarfs, that the subsequent events (e.g., merging, accretion) are efficient at erasing the characteristic mass scale for fragmentation. Future IMF studies that probe masses below the characteristic fragmentation mass can provide more stringent constraints on the role of fragmentation in the star formation process.
## 10 Summary and Conclusions
Using HST NICMOS narrow band imaging, we have measured the 1.9 $`\mu `$m water band strengths of low-mass objects in the IC348a subcluster. With the magnitudes and spectral types thereby obtained, we are able to separate cluster members from background objects. Comparisons with recent evolutionary tracks (B98, DM98) imply that our study probes a mass range extending from low-mass stars ($`0.7M_{}`$) down to the bottom of the deuterium burning main sequence ($`0.015M_{}`$). The mean age of the subcluster is 3 Myr with the B98 tracks and 1 Myr with the DM98 tracks. These results are subject to uncertainties in the evolutionary tracks and the appropriate conversions between theoretical ($`L_{}`$, $`T_{\mathrm{eff}}`$) and observed (e.g., spectral types, magnitudes) quantities which remain somewhat uncertain. We also confirm an age spread to the cluster, as found previously (Lada & Lada (1995); Herbig (1998); LRLL ), from $`<1`$ to $`1020`$ Myr.
Assuming that the hydrogen- and deuterium-burning mass limits delimit the brown dwarf population, we have fully sampled the brown dwarf population at ages up to 3 Myr and extinctions up to the cluster average ($`A_K=0.3`$). We find $`2030`$ cluster substellar candidates (depending on the choice of evolutionary tracks) which represents a significant fraction, $`25`$%, of all cluster M dwarfs by number, but only a small fraction, $`5`$%, by mass. The mass function derived for the subcluster, $`dN/d\mathrm{log}MM^{0.5},`$ is similar to that recently obtained for the Pleiades over a more limited mass range (Bouvier et al. (1998)), and apparently less abundant in low mass objects than the local field population (Reid et al. (1999)). In contrast, the derived mass function appears significantly more abundant in brown dwarfs than the mass function of companions to nearby solar-type stars at separations $`<5`$ AU (Marcy et al. (2000)).
The apparent difference may indicate that substellar objects form more readily in isolation than as companions. Alternatively, the difference may represent the result of evolutionary effects such as accretion (by the star) or dynamical ejection, which will tend to deplete the companion population and, in the latter case, contribute low mass objects to the local field population. Given the population statistics from precision radial velocity studies, if these evolutionary mechanisms are the underlying physical cause for the different IMFs, they must preferentially deplete the brown dwarf population compared to the lower mass planetary companion population, which appears to be present in significant numbers.
More generally, we find that the imaging photometric technique used in this study is a potentially powerful approach to the study of low mass populations in young clusters. As demonstrated here, it is possible to study a large range in mass ($`0.50.015M_{},`$ a factor of $`>30`$ in mass) with a single technique. To summarize, the utility of this approach derives from the multi-object approach inherent in a filter photometric method; the sensitivity of the index due to the rapid variation of the water band strength with late-M spectral type; the approximate orthogonality of the reddening vector to the variation with spectral type so that reddening errors do not introduce significant spectral type errors; and the long wavelength of the index which improves the sampling of embedded populations.
To stress this latter point, we can consider the depth to which one would have to carry out spectroscopy in the $`I`$-band to recover similar information for IC348. Our completeness limit for spectral typing is $`K16.5`$. With this level of completeness, we have sampled a significant fraction of the low-mass cluster population. For example, to $`A_K=0.5`$, the B98 and Burrows et al. (1997) tracks imply that at ages of 20, 3, and 1 Myr, we are complete to 35, 16, and 9 $`M_J`$. For a more extreme extinction of $`A_K=2`$, the B98 model suggests that at ages of 20, 3, and 1 Myr, we are complete to 100, 32, and 25 $`M_J`$. In contrast, for a spectroscopic study in the $`I`$-band, $`A_I/A_K5`$ and for the late spectral types probed in the present study, $`IK4.5`$. Consequently, for extinctions of $`A_K=`$0.5 and 2, the corresponding limiting magnitude is $`I=23`$ and 29. In contrast, optical spectral typing with existing 10-m telescopes is currently limited to sources brighter than $`I19.5.`$
Although the present study made use of narrow band filters and the ability to work above the Earth’s atmosphere with HST, the technique used here might find useful extrapolation to both broader filters and to ground-based observations. With broader filters, it would be possible to study objects at lower, planetary masses, as well as more distant, richer clusters where the spatial multiplexing advantage of a filter photometric technique could be used to better advantage. We will explore these possibilities in a future paper (Tiede et al. (2000)).
We are grateful to Nick Bernstein and Alex Storrs for their extensive help in getting our program scheduled and executed; to Matt Lallo and Russ Makidon for their explanation of HST pointing errors; to Marcia Rieke and Paul Martini for their help with the photometric calibration; to Wolfram Freudling for advice on NICMOSlook; to Tod Lauer who helped us investigate the impact of intra-pixel sensitivity on our photometry; to France Allard and Peter Hauschildt for sharing their atmosphere models; to Adam Burrows who shared his brown dwarf evolutionary models; to Martin Cohen who helped us estimate the background population; and to Kevin Luhman for useful advice and sharing his IC348 results with us in advance of publication. We are also grateful to Charles Lada and Tom Greene for useful discussions regarding this project, and to Arjun Dey and the anonymous referee whose comments significantly improved the manuscript. Support for this work was provided by NASA through grant number GO-07322.02-96A from the Space Telescope Science Institute, which is operated by AURA, Inc., under NASA contract NAS5-26555.
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# Particle versus Field Structure in Conformal Quantum Field Theories
## 1 A few introductory remarks
Ideas about the use of conformal quantum field theory entered particle physics for the first time at the height of the Kramers-Kronig dispersion relations . They were met with reactions ranging from doubts to outright rejection and the subject lay dormant for another 10 years when it reemerged on the statistical mechanics side in connection with second order phase transitions.
In the next section we will show that these early doubts of the old-time particle physicists were partially justified, because the particle structure in CQFT is indeed incompatible with interactions. However far from supplying a coffin nail for its utility in high energy physics, this no-go theorem also contains the message that one must use finer concepts in order preserve the usefulness of conformal quantum field theory as a theoretical laboratory for particle physics. There are massive particle-like objects (“infraparticles” ) which have a continuous mass distribution with an accumulation of spectral weight at $`p^2=m^2`$ whose generating local fields have an anomalous non-integer (non-semi-integer in the case of Fermion fields) contribution to their long distance behavior. In a CQFT long and short distance behavior coalesce and the accumulation of spectral weight at $`p^2=0`$ which becomes related to the anomalous dimension of operators is the vestige of the particle interaction in the massive parent theory from which the CQFT arose by taking the scale-invariant limit. This structure is the collective effect of a total collapse of all multiparticle thresholds on top of each other. The standard LSZ large time scattering limit does not commute with this scaling limit, in fact the LSZ limit of such fields vanishes. It is believed that in order to re-extract from such a situation anything which resembles particle physics one has to apply a more general form of scattering theory which is based on expectation values and probabilities for inclusive cross sections (where outcoming “stuff” below a prescribed energy-momentum resolution is not registered) rather than on amplitudes. But it is presently not clear how one can achieve this. In the case of infraparticles (the electron in QED which is inexorably linked to its photon-cloud) where one also meete a situation of coalescing thresholds, this generalized scattering theory is known to be very useful .
Recently there has been a quite different and conceptually<sup>1</sup><sup>1</sup>1The attribute “conceptually” here refers to the local quantum physical aspects and not to differential-geometric ones. less ambitious but formally quite attractive idea which promises to strengthen the utility of CQFT for particle physics and which is presented in the third section. It basically consists in finding a theory which radically reprocesses the spacetime interpretation and degrees of freedom of CQFT in such a way that now the “energy momentum vector” $`R_\mu `$ of the Dirac-Weyl compactified world $`\overline{M}`$ becomes the bona fide energy momentum instead of $`P_\mu `$ which in standard canonical or functional terminology means that $`R_\mu `$ is the one related to an action and not $`P_\mu `$. If one insists that this total reshuffling of physical interpretation should leave the basic mathematical building blocks (a certain generating set of algebras and the symmetry group structure) untouched, then there is only one answer: an associated anti De Sitter (AdS) theory . The nontrivial reprocessing leads to a mathematical isomorphism as described in i.e. it goes far beyond that picture about the AdS-CQFT correspondence which is limited to the (infinitely remote) boundary of AdS (see in particular the remarks at the end of ). The AdS appearance of the AdS structure as a kind of reprocessed CQFT is less surprizing if one recalls the 6-dimensional lightcone formalism which one uses in order to obtain an efficient description of the conformal compactification $`\overline{M}`$ of Minkowski space $`M`$ and the construction of its covering $`\stackrel{~}{M}`$ .
In this way one obtains a (perturbative) new constructive non-Lagrangian access to CQFT which opens a new window into the realm of CQFT beyond those few 4-dimensional Lagrangian candidates for which one had to use a combination of gauge theory with supersymmetry. This means that one has no guaranty that the conformal side at all permits a description in terms of an action.
## 2 Particle Structure and Triviality
We start with recalling an old theorem which clarifies the relation between the particle-versus-field content of conformal field theories. To be more precise the following statement is a result of the adaptation of a combination of several theorems
###### Theorem 1
The existence of one-particle states in conformally invariant theories forces the associated interpolating fields to be canonical free fields. The only particle-like structures consistent with interactions are hidden in the structure of those interpolating fields which have anomalous dimensions and whose mass spectrum is continuous with an accumulation of weight at $`p^2=0,p_0>0.`$
The easiest way to get a first glimpse at this situation is to look at conformal two-point functions
$$\psi (x)\psi ^{}(y)=\{\begin{array}{c}c\frac{1}{\left(xy\right)^2},dim\psi =1\\ c(\frac{1}{\left(xy\right)^2})^{d_\psi },dim\psi =d_\psi >1\end{array}$$
(1)
In the first case the application of the LSZ large time scattering limit yields
$$\psi (x)\psi ^{}(y)=\psi _{in}(x)\psi _{in}^{}(y)$$
(2)
which preempts the equality $`\psi =\psi ^{in}=\psi ^{out},`$ whereas in the anomalous case the large distance fall-off is too strong in order to be reconcilable with the mass shell structure of a zero mass particle which means
$$\psi (x)\stackrel{LSZ}{}0$$
(3)
It is worthwhile to reconsider the argument which leads to the absence of interaction in the space created by the interpolating field $`\psi .`$ The crucial observation is that the presence of a zero mass scalar particle state vector $`|p`$ with
$$p\left|\psi \right|00$$
(4)
forces $`\psi `$ to have a two-point function with a canonical scale dimension dim$`\psi =1.`$ The special feature of conformal invariance is that this implies that the two-point function is free i.e.
$$0\left|\psi ^{}(x)\psi (y)\right|=c\frac{1}{\left[(xyi\epsilon )\right]^2}$$
(5)
Such a conclusion relating canonical short distance dimension with absence of interactions cannot be drawn in the massive case. However the following theorem which was proven in the late 50<sup>ies</sup> by Jost and the present authors, and can be found in , holds for both cases:
###### Theorem 2
The freeness of the $`\psi `$ two-point function implies the field $`\psi `$ to be a free field in Fock space.
The guiding idea is to show that a localized operator or pointlike field which vanishes on the vacuum, vanishes automatically on all states i.e. is the zero operator. This is a consequence of the Reeh-Schlieder theorem which in conformal field theory is also known under the name state-field relation). It says that the operators from a region with a nontrivial causal complement (or fields smeared with test functions with support in such a region) act cyclically on the vacuum (and on any other finite energy state). If we denote by $`𝒜(𝒪)`$ either the polynomial -algebra of unbounded smeared fields with supports of testfunctions in $`𝒪`$ or the affiliated bounded operator algebra, this cyclicity property reads
$$\overline{𝒜(𝒪)\mathrm{\Omega }}=$$
(6)
where the bar denotes the closure and $`H`$ is the Hilbert space generated by all fields (bosonic and fermionic). Since (for fermionic $`\psi `$ there will be a change of sign)
$$\psi (x)𝒜(𝒪)\mathrm{\Omega }=𝒜(𝒪)\psi (x)\mathrm{\Omega }$$
(7)
if we choose $`O`$ spacelike with respect to $`x,`$ the vanishing of the “current” $`j(x)=(_\mu ^\mu +m^2)\psi (x)`$ on the vacuum implies the vanishing on the dense set $`𝒜(𝒪)\mathrm{\Omega }`$ and hence (operators in physics are closable) on all $`.`$The next step consists in proving that the commutator of two $`\psi s`$ on the vacuum is a c-number
$$\left([\psi (x),\psi (y)]i\mathrm{\Delta }(xy)\right)\mathrm{\Omega }=0$$
(8)
It then follows according to the previous argument that the bracket vanishes identically. We prove this last relation by using the frequency decomposition $`\psi =\psi ^{()}+\psi ^{(+)}`$ (which follows from $`j0)`$ in the commutator
$$[\psi (x),\psi (y)]\mathrm{\Omega }=([\psi ^{(+)}(x),\psi ^{(+)}(y)]+\psi ^{()}(x),\psi ^{(+)}(y)\psi ^{()}(y),\psi ^{(+)}(x))\mathrm{\Omega }$$
(9)
where we omitted all annihilation terms. The on-shell creation with subsequent on-shell annihilation as in the last two terms and the physical spectrum condition only admits the vacuum as its energy momentum content and therefore they yield a c-number which, by a finite renormalization of $`\psi `$ if necessary, yields
$$(\psi ^{()}(x),\psi ^{(+)}(y)\psi ^{()}(y),\psi ^{(+)}(x))\mathrm{\Omega }=i\mathrm{\Delta }(xy)\mathrm{𝟏}\mathrm{\Omega }$$
(10)
Since this and the full commutator is causal, the first term on the right hand side has to vanish all by itself. But on the other hand it is the separate Fouriertransform of momenta which lie on the forward mass shell and hence it is the boundary value of an analytic function in two complex 4-vectors of the form $`z=\xi i\eta ,\eta `$ from the forward light cone. However an analytic function which vanish on an open set on its boundary vanished identically (generalized Schwartz reflection principle). The resulting relation on the vacuum holds according to the previous arguments for the operators and therefore we obtained the characterizing relation for a free field. The generalization to any spin including half-integer values is now a routine matter. A closer look at the zero mass situation reveals that contrary to the massive case where the difference of two on-shell vectors is either spacelike or zero, the difference of two lightlike vectors may in addition be lightlike but this only happens for parallel vectors. Since this special configurations should not matter in the sense of L<sup>2</sup>-integrability of zero mass particle wave functions one again expects at least for $`d>1+1`$ the above result. However a mathematical proof of this result turned out to be quite nontrivial .
It is very helpful to place the above theorem into the setting of a more general theorem relating interactions and particle properties in general local quantum physics which states that operators localized in sub-wedge regions in interacting theories which possess nontrivial matrix elements between vacuum and one-particle states necessarily show the phenomenon of vacuum polarization i.e. operators which create polarization-free one-particle states exist only in interaction free field theories. Polarization-free-generators (PFG) which create pure one-particle states from the vacuum do however exist in any QFT if their localization region is a semi-infinite wedge region or larger . Since in conformal theories the wedge region is conformally equivalent to a compact double cone, a conformal one-particle structure according to this more general theorem is only possible in conformal free field theories.
The above argument is typical for a real-time structure which cannot be unraveled in the euclidean formulation.
## 3 Trying to make the best out of it
The negative result on the compatibility of zero mass particle structure with nontriviality of conformal theories should not be misread as an incompatibility with an intuitive idea about what constitutes particle-like excitations. The point here is that conformal theories in particle physics should be considered as the zero mass (scaling) limits of massive theories with mass gaps for which the LSZ scattering theory can be derived. In the scaling limit all the multiparticle thresholds in momentum space coalesce on top of each other and build up the possibly anomalous dimension. In this limit the Wigner particle theory (irreducible representation of the Poincaré group) and with it the prerequisite of the LSZ scattering theory gets lost in the presence of interactions, a fact which we have demonstrated above where it was shown that the field is either free or the LSZ limits are zero. So the right question would be: can one think of a more general scattering theory which may recuperate some of the lost structure in the aforementioned collapse of multiparticle cuts on top of each other? There is indeed another particle concept (“infraparticle”) which goes together with a generalized scattering theory build on inclusive scattering probabilities instead of amplitudes . This concept is expected to distinguish those anomalous dimensional fields which are of relevance in particle physics (which originate from the previous collapse in the scaling limit) from mere mathematical constructs as e.g. generalized free fields with anomalous dimensions. But we think that for the problem at hand, namely the formulation of a theory of anomalous dimension, we do not need to enter this deep and difficult issue of particle-like interpretation since here we restrict our interests in conformal theories as a simplified theoretical laboratory for field- and algebra- aspects and not for the study of particles and their scattering theory. We believe that the setting of local observable algebras which fulfill in addition to Einstein causality also Huygens principle for timelike distances contains all scale limits of theories which are of interest for particle physics and that interaction in this setting is characterized by the appearance of charge-carrying fields with anomalous dimensions. In view of the above No-Go theorem we will consider the noncanonical (anomalous dimension) nature of those fields as our pragmatic definition of interaction in this conformal setting. But we defer this analysis to a following longer paper which contains the relevant mathematical machinery .
As a consequence the observable algebra of an interacting conformal field theory (conserved currents etc.) should not have the structure of composites of free fields (e.g. free currents) since otherwise the fields carrying the superselected charges may not have anomalous dimensions. Apart from normalization constants the 2- and 3-point functions of conformal observable fields (currents) are indistinguishable from those formed with free composites with the same integer dimensions. If all correlations would be indistinguishable from those of free composites (total protection) then also the charge-carrying fields associated with such observables can be shown to be free.
A weak form of what in the case of conformal SYM theories has been called (partial) “protection” would be one where the relative normalization between 2-and 3-point functions is that of free composites (partial protection). Apparently perturbative supersymmetry causes partial protections . Although such models hardly represent realistic particle physics, they are the only Lagrangian candidates for d=1+3 nontrivial conformal field theories and may yet turn out to be the first 4-dimensional mathematically completely controllable models. The interest and fascination in conformal field theories originates to a large part from the well-founded belief that the simplest nontrivial 4-dimensional conformal field theories which will break the age old existence deadlock<sup>2</sup><sup>2</sup>2In any area of Theoretical Physics there always have been plenty of nontrivial mathematically controllable illustrations which demonstrate the nontrivial physical content of the conceptual basis of those areas, not so in 4-dim. QFT. This annoying totally singular situation has been sometimes overemphasized at the cost of practical calculations, but most of the time it went totally ignored. for nontrivial quantum field theories in physical spacetime. For this one wants to have as much protection as possible without ending with a free conformal theory.
Instead of entering an ambitious program in order to extract the particle physics “honey” from CQFT which requires a heavy conceptual investment in the area of a generalized scattering theory, there is another way which is more faithful to the formal aspects with which QFT is often identified (erroneously in my opinion, if one uses them for a definition of QFT) namely canonical formalism and/or functional integrals. It starts from the observation that in addition to the translation generator $`P_\mu `$ there is another translation-analogue described by a Lorentz-vector $`R_\mu .`$ It has a timelike purely discrete spectrum and the L-invariant “mass” $`m_c`$ with $`m_c^2=R_\mu R^\mu `$ plays a similar role as the rigid rotation operator $`L_0`$ in chiral theories. In fact it describes a generalized rotation around the Dirac-Weyl compactified Minkowski space $`\overline{M}S^3\times S^1.`$ Therefore it is not surprising that the bottom of the spectrum of $`m_c`$ is the anomalous part of the scaling dimension common to a whole equivalence class of fields which carry the same superselected charge. But despite all analogies to $`P_\mu `$ this operator is not related to an imagined functional integral action of CQFT. Nevertheless one can ask the question: is there a theory whose Lagrangian can be associated with a Hamiltonian interpretation of $`R_0\mathrm{?}`$ In order for this new theory to be useful for particle physics it should keep the same algebraic and group-theoretical building blocks as CQFT i.e. one seeks a mathematical isomorphism which goes hand in hand with that total physical reprocessing which is necessary to accomplish such an impossible looking task. The unique answer is the AdS-CQFT correspondence which was proven to be a such a “radical” isomorphism .
Although this step does not completely answer the question posed at the beginning of how to extract and analyze the particle content of CQFT, it goes a long way to open up conformal field theory as a genuine theoretical laboratory for particle physics. And last not least it facilitates the unsolved problem number one: find a nontrivial physically relevant (i.e. one which fits at least the conceptual framework of local quantum physics, even if it falls short in describing nature) and mathematically controllable model in 4-dimensional QFT.
The presented arguments suggest strongly that there exists a whole world of non-Lagrangian non-supersymmetric CQFT (in the sense that they cannot be accessed in the standard perturbative way) besides the Lagrangian SYM family. In fact the perturbative calculations in the literature already give some support in this direction. This is most visible in although these authors, evidently under the strong spell of the string-theoretic origin of the AdS-CQFT, do not interprete their calculations from this viewpoint.
The possible non-Lagrangian nature of most CQFT is in a certain way explained by Rehren’s deep observation that due to the isomorphic nature of the AdS-CQFT relation there must be degrees of freedom on the conformal side which cannot be described in terms of local fields namely those which originate from the AdS bulk (and not from the boundary) and which are necessary in order to return $`CQFTAdS`$. This leaves the interesting question of what should one make of the original observation by which the protagonists of the AdS-CQFT correspondence found this relation which is the relation between two Lagrangian field theories namely the conformal SYM model with some form of AdS supergravity . Since this is based on consistency checks within string theory which owes its widespread acceptance to perturbative mathematical consistency and a kind of globalized social contract but certainly not to its harmonious coexistence with the principles underlying particle physics, there is reason for some scepticism; in particular because such degrees of freedom would be easily overlooked in perturbative calculations on the CQFT side. It cannot be overstressed that this correspondence is very different and much more radical then those which arise from a different choice of “field coordinates”. It is impossible to understand its full content in terms of pointlike physical fields.
## 4 Some concluding remarks
If, as argued in this letter, the AdS theories are a useful new calculational tool which open up CQFT to particle physics studies within the standard Lagrangian quantization framework, than perhaps with an additional conceptual investment one could directly understand the structure underlying the anomalous dimension spectra within CQFT i.e. without the described reprocessing on the AdS side. This turns out to be true and will be the subject of a subsequent paper since the necessary conceptual investment does not fit the format of a letter like this.
Acknowledgements: I am indebted to Detlev Buchholz and Karl-Henning Rehren for a helpful exchange of emails. Furthermore I would like to thank Francesco Toppan for interesting questions which helped in shaping the presentation.
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# On the Entangling Power of Quantum Evolutions
## Abstract
We analyze the entangling capabilities of unitary transformations $`U`$ acting on a bipartite $`d_1\times d_2`$-dimensional quantum system. To this aim we introduce an entangling power measure $`e(U)`$ given by the mean linear entropy produced acting with $`U`$ on a given distribution of pure product states. This measure admits a natural interpretation in terms of quantum operations. For a uniform distribution explicit analytical results are obtained using group-theoretic arguments. The behavior of the features of $`e(U)`$ as the subsystem dimensions $`d_1`$ and $`d_2`$ are varied is studied both analytically and numerically. The two-qubit case $`d_1=d_2=2`$ is argued to be peculiar.
From the beginning it has been argued that entanglement is one of the crucial ingredients that allows Quantum Information processing to outperform, for certain tasks, any classically operating device. In this sense entanglement represents a uniquely quantum resource whose production is a sort of elementary prerequisite for any Quantum Computation (QC). Such a basic task is accomplished by unitary transformations $`U`$ i.e., quantum evolutions acting on the state-space of the multi-partite system that describe non-trivial interactions between the degrees of freedom of the different subsystems. Even though almost all the unitaries satisfy this latter requirement , it is quite natural to ask how different $`U`$’s are efficient, according to some criterion to be specified, as entanglers, and then by using such a criterion to analyze the full manifold of bi-partite quantum evolutions.
In this paper we address this issue by introducing over the space of bi-partite unitaries a measure for their entangling power. This is done by considering how much entanglement is produced by $`U`$ on the average acting on a given distribution of unentangled quantum states. The kind of situation we have in mind is a procedure for entanglement production in which one randomly generates product states (the ”cheap” resource ) according to some probability distribution $`p`$ and then applies the transformation $`U.`$ The average entanglement obtained with the above scheme will be our measure $`e_p(U)`$ of the quantum evolution $`U.`$
It is important to stress that these $`U`$’s can represent different objects, both from the logical and physical point of view. Some prototypical instances are given by: a) A quantum computation using a pair of quantum registers. Here the entangling power measure will quantify how the computation $`U`$ is efficient in making the first (say memory) and the second (say computational) registers entangled. This kind of entanglement, that represents mutual information between the two registers, has been recently proved to play a role in QC viewed as a communication process . b) The global evolution of a system plus its environment. In this case $`e_p(U)`$ measures the decohering power of the system-environment coupled evolution $`U.`$ Engineering weak decoherence then amounts to design an optimal $`U`$ with respect to the criterion of minimal entangling power. c) A single two-subsystem e.g., two-qubits, gate in a quantum-network. Now the entangling $`U`$ are the two-qubit gates needed to get universal QC .
To formalize our setting let us consider a bipartite quantum system with state space $`=_1_2`$ where dim$`_i=d_i(i=1,2)`$ and $`U𝒰()U(d_1d_2).`$ If $`E`$ is an entanglement measure over $``$ we define the entangling power of $`U`$ (with respect to $`E`$) as
$$e_p(U):=\overline{E(U|\psi _1|\psi _2)}^{\psi _1,\psi _2}$$
(1)
where the bar denotes the average over to all the product states $`|\psi _1|\psi _2.`$ distributed according some probability density $`p(\psi _1,\psi _2)`$ over the manifold of product states.
We shall use as entanglement measure of $`|\mathrm{\Psi }`$ the linear entropy
$$E(|\mathrm{\Psi }):=1\text{tr}_1\rho ^2,\rho :=\text{tr}_2|\mathrm{\Psi }\mathrm{\Psi }|.$$
(2)
This quantity measures the purity of the reduced density matrix $`\rho ,`$ it can be regarded as a kind of ”linearized” version of the von Neumann entropy $`S(\rho )=\text{tr}\rho \mathrm{ln}\rho ,`$ which is known to provide the essentially unique measure of entanglement for bi-partite pure quantum states. One has that $`0E(|\mathrm{\Psi })11/d`$ where the lower (upper) bound is reached iff $`|\psi `$ is a product state (maximally entangled). The measure (2) has, with respect to $`S(\rho ),`$ the definite advantage of being a polynomial in $`|\psi .`$
Now we introduce some notations. We shall denote by $`T_{ij},(i,j=1,\mathrm{},4)`$ the transposition between the $`i`$-the and the $`j`$-th factor of $`^{\mathrm{\hspace{0.17em}2}}:=_2(\text{ }\mathrm{C}^{d_1}\text{ }\mathrm{C}^{d_2})(\text{ }\mathrm{C}^{d_1}\text{ }\mathrm{C}^{d_2}).`$ Notice that $`T_{12}`$ and $`T_{34}`$ are well defined elements of $`𝒰(^{\mathrm{\hspace{0.17em}2}})`$ only when $`d_1=d_2,`$ in this latter case such operators will be referred to as swaps. Moreover – when $`_i_j`$ – one defines the projectors $`P_{ij}^\pm :=2^1(𝟙\pm 𝕋_{𝕚𝕛})`$ over the totally symmetric (antisymmetric) subspaces of $`_i_j,`$ the latter being thought of as embedded in $`^{\mathrm{\hspace{0.17em}2}}.`$ The space End $`(^{\mathrm{\hspace{0.17em}2}})`$ is endowed with the Hilbert-Schmidt scalar product $`<A,B>:=\text{tr}(A^{}B).`$ Finally with $`𝒮()`$ we shall denote the space of density matrices over $`.`$
Proposition 0 The entangling power (1) is given by
$`e_p(U)=2\text{tr}[U^{\mathrm{\hspace{0.17em}2}}\mathrm{\Omega }_pU^{\mathrm{\hspace{0.17em}2}}P_{13}^{}],`$ (3)
where $`\mathrm{\Omega }_p:=𝑑\mu (\psi _1,\psi _2)(|\psi _1\psi _1||\psi _2\psi _2|)^{\mathrm{\hspace{0.17em}2}}𝒮(^{\mathrm{\hspace{0.17em}2}})`$ and $`d\mu `$ denotes the measure over the product state manifold induced by the probability distribution $`p(\psi _1,\psi _2).`$
Proof. Let us observe that Eq. (2) can be written in a linear form using the identity tr$`[(AB)T]=\text{tr}(AB)`$ where $`T`$ is the swap. Then $`E(|\mathrm{\Psi })=1\text{tr}(|\mathrm{\Psi }\mathrm{\Psi }|^{\mathrm{\hspace{0.17em}2}}T_{13}).`$ Form this remark and the definition (1) it follows immediately (3)$`\mathrm{}`$
The result above express $`e_p(U)`$ as the expectation value over $`\mathrm{\Omega }_p`$ of the positive operator $`2U^{\mathrm{\hspace{0.17em}2}}P_{13}^{}U^{\mathrm{\hspace{0.17em}2}}.`$ This latter operator can be viewed as the effect associated to the completely positive (CP) map $`\mathrm{\Phi }_U`$ on $`𝒮(^{\mathrm{\hspace{0.17em}2}})`$ given by $`\mathrm{\Phi }_U:\mathrm{\Omega }2P_{13}^{}U^{\mathrm{\hspace{0.17em}2}}\mathrm{\Omega }U^{\mathrm{\hspace{0.17em}2}}P_{13}^{}.`$ This remark allows us to interpret the entangling power (5) as probability of success of a two-party (A and B) quantum protocol (see Fig. 1). Suppose $`A`$ ($`B`$) owns spaces $`_1`$ and $`_3`$ ($`_2`$ and $`_4`$)
a) A and B generate pairs of states $`|\psi _1|\psi _2`$ according the distribution probability $`p(\psi _1,\psi _2)`$ \[$`\mathrm{\Omega }_p`$ is prepared\] b) Apply to each member of the pair the joint transformation $`U`$ \[action of $`U^{\mathrm{\hspace{0.17em}2}}`$\] c) Perform a projective measurement of $`\sqrt{2}P_{13}^{}.`$
Eq. (3) nicely displays several properties required for any entangling measure for bi-partite unitary evolutions. i) $`e_p(U_1U_2U)=e_p(U)(U_iU(d_i))`$. Indeed from the $`U(d_1)`$-invariance of $`P_{13}^{}`$ one finds $`e_p(U_1U_2U)=2\text{tr}[U^{\mathrm{\hspace{0.17em}2}}\mathrm{\Omega }_pU^{\mathrm{\hspace{0.17em}2}}(U_1U_2)^{\mathrm{\hspace{0.17em}2}}P_{13}^{}(U_1U_2)^{\mathrm{\hspace{0.17em}2}}]=2\text{tr}[U^{\mathrm{\hspace{0.17em}2}}\mathrm{\Omega }_pU^{\mathrm{\hspace{0.17em}2}}(U_1)_{13}^{}P_{13}^{}(U_1)_{13}]=e_p(U).`$ Where $`(U_1)_{13}:=U_1𝟙𝕌_\mathrm{𝟙}𝟙.`$ ii) When $`d_1=d_2,`$ by denoting with $`T`$ the transposition between the two factors of $`,`$ one has $`e_p(TU)=e_p(U).`$ This stems from $`T^{\mathrm{\hspace{0.17em}2}}P_{13}^{}T^{\mathrm{\hspace{0.17em}2}}=P_{24}^{}.`$ This leaves $`e_p(U)`$ unchanged, indeed this label change amounts simply to the replacement tr$`{}_{1}{}^{}`$ tr<sub>2</sub> in Eq. (2). Since, for pure states, the two reduced density matrices are isospectral the linear entropy is unchanged. Moreover if $`\mathrm{\Omega }_p`$ is swap invariant i.e., $`p(\psi _1,\psi _2)=p(\psi _2,\psi _1)`$ one also has $`e_p(UT)=e_p(U).`$ iii) One has $`e_p(𝟙)=\mathrm{𝟘}.`$ This simply because $`\mathrm{\Omega }_pP_{13}^{}=0.`$ Indeed, form the definition (3) one has $`\mathrm{\Omega }_pT_{13}=\mathrm{\Omega }_p.`$ iv) From the previous remarks it follows that the entangling power is constant along the orbits in $`𝒰()`$ of the left action of the subgroup of the bi-local operations $`U_1U_2.`$ In particular $`e`$ vanishes on all the elements of such a group. In the symmetric case $`d_1=d_2`$ the group is extended by the swap $`T.`$
Different distributions $`p(\psi _1,\psi _2)`$ would result in very different $`e(U).`$ An extreme example of this obvious remark is provided by transformations $`U`$ that simply permutes elements of a given basis $`|i|j`$ of $``$. If $`p`$ is supported just on this basis the associated $`e_p(U)`$ vanishes identically, while we shall show later for a different probability distribution that $`U`$’s can even be maximally entangling. Another example is given in the context of the case b) mentioned in the introduction. Suppose $`U`$ admits a decoherence-free subspace $`𝒞_1`$ the if $`p`$ is, for any $`|\psi _2,`$ supported in $`𝒞`$ then again $`e_p(U)=0.`$
From now on we focus on the case in which $`p`$ is the uniform distribution $`p_0`$. With this term we refer to the unique $`U(d_1)\times U(d_2)`$-invariant probability distribution i.e., $`p(\psi _1,\psi _2)=p(U_1\psi _1,U_2\psi _2).`$ When all the product state are considered to be equally easy to be prepared, this latter assumption on $`p`$ is quite natural from the physical point of view . Moreover, in view of its symmetry, the uniform $`p`$ will result in a great computational simplification that will allow for an explicit analytical evaluation of the average over the product-state manifold that appears in Eq. (1).
Let us begin by proving an easy group-theoretic Lemma that will play an essential technical role in the following.
Lemma $`\mathrm{\Omega }_{p_0}=4C_{d_1}C_{d_2}P_{13}^+P_{24}^+,C_d^1:=d(d+1).`$
Proof. Since the uniform distribution factorizes we can consider separately the average $`\omega _{13}`$ with respect to $`|\psi _1`$ ( on the first and the third factor of $`^{\mathrm{\hspace{0.17em}2}}`$) and the one $`\omega _{24}`$ with respect to $`|\psi _2`$ ( on the second and the fourth factor of $`^{\mathrm{\hspace{0.17em}2}}`$) then $`\mathrm{\Omega }_{p_0}=\omega _{13}\omega _{24}.`$ Let first observe that in view of definition (3) one has that $`\mathrm{\Omega }_{p_0}`$ is supported in $`P_{13}^+P_{24}^+^{\mathrm{\hspace{0.17em}2}}`$ i.e., $`\mathrm{\Omega }_p`$ is symmetric under the exchange of the first (second) and the third (fourth) factor. Moreover since the uniform distribution is $`U(d_1)\times U(d_2)`$ invariant one has $`[U_1^{\mathrm{\hspace{0.17em}2}},\omega _{13}]=0,U_1U(d_1),`$ and analogously for $`\omega _{24}.`$ Since the $`U_1^{\mathrm{\hspace{0.17em}2}}`$’s act on the totally symmetric subspace irreducibly, it follows from the above commutation relation and the Schur Lemma that $`\omega _{13}=2CP_{13}^+.`$ The normalization constant is found by the condition tr$`\omega _{13}=1.`$ Reasoning in the same way for $`\omega _{24}`$ one gets the desired result. $`\mathrm{}`$
Proposition 1 The entangling power (1), with respect the uniform distribution, is given by
$`e_{p_0}(U)`$ $`=`$ $`1C_{d_1}C_{d_2}{\displaystyle \underset{\alpha =0,1}{}}I_\alpha (U)`$ (4)
$`I_\alpha (U)`$ $`=`$ $`t(\alpha )+<U^{\mathrm{\hspace{0.17em}2}}(T_{1+\alpha ,3+\alpha })U^{\mathrm{\hspace{0.17em}2}},T_{13}>,`$ (5)
where $`t(\alpha ):=\text{tr}T_{1+\alpha ,3+\alpha }.`$
Proof. It is just a calculation. Insert $`\mathrm{\Omega }_{p_0}=C_{d_1}C_{d_2}(𝟙+𝕋_{\mathrm{𝟙𝟛}})(𝟙+𝕋_{\mathrm{𝟚𝟜}})`$ in Eq. (3). Notice that one has tr $`T_{13}=d_1d_2^2,\text{tr}T_{24}=d_1^2d_2.`$ $`\mathrm{}`$
From the relations $`[T_{1+\alpha ,3+\alpha },(U_1U_2)^{\mathrm{\hspace{0.17em}2}}]=0,(\alpha =0,1)`$ it follows that both the functions $`I_0`$ and $`I_1`$ are invariant under the two-sided i.e.,left and right, action of bi-local unitaries e.g., $`I_1(U)=I_1(UU_1U_2).`$ Moreover in the symmetric case $`d_1=d_2`$ it follows from the above that the entangling power (5) can be written in a manifestly swap invariant form $`e_{p_0}(U)=1C_d^2_{i=0}^1I_d(T^iU),`$ where $`I_d(U):=d^3+<U^{\mathrm{\hspace{0.17em}2}},T_{13}U^{\mathrm{\hspace{0.17em}2}}T_{13}>.`$
The entangling power $`e`$ defines a random variable over $`𝒰()`$ if the latter, endowed with the Haar measure, is considered as a probability space. Therefore it makes sense to consider the associated density of probability distribution $`q(e).`$ Moreover, since the the manifold of unitary transformations over $``$ is compact, the obviously continuous mapping $`Ue(U)`$ must achieves extrema, in particular $`\overline{U}𝒰():e(\overline{U})=\mathrm{max}_Ue(U).`$ Such maximally entangling $`\overline{U}`$’s will be referred to as optimal. In Fig. 2 are reported the $`q(e)`$’s obtained numerically for the cases $`d\times d`$ with $`d=2,3,4.`$ While in the cases $`d2`$ the function $`q(e)`$ vanishes on both the lower and the upper sides of the allowed range of $`e,`$ it it remarkable that the two-qubit case $`d=2`$ shows a peculiar feature: $`q(e)`$ is a monotonic function of $`e.`$ This implies that most of the two-qubit gates $`U`$ correspond to nearly optimal ones. Moreover, as will be discussed later in details, the entangling power of optimal $`U`$’s does not correspond to an upper bound that is instead reached by all the other cases for $`d.`$ In this sense the prototypical quantum information case of two qubits is quite singular. A first very natural question is how on the average an operator is entangling i.e., the mean of the $`q(e).`$
Proposition 2 The average of the entangling power $`e_{p_0}(U)`$ over $`U(d_1d_2)`$ is given by
$$\overline{e_{p_0}(U)}^U=\frac{(d_11)(d_21)}{d_1d_2+1}$$
(6)
Proof. To prove Eq. (6) we first notice that, in view of definition (5), to compute the mean of the entangling power amounts to compute the average of the entanglement measure over all the states $`|\mathrm{\Psi }`$ ( not just over the product states). The trace of the square of the reduced density matrix of $`|\mathrm{\Psi }`$ is given by $`\text{tr}(|\mathrm{\Psi }\mathrm{\Psi }|^{\mathrm{\hspace{0.17em}2}}T_{13}).`$ Now we take the average with respect $`|\mathrm{\Psi }`$ using again Lemma $`[d_1d_2(d_1d_2+1)]^1\text{tr}[(𝟙+𝕋_{\mathrm{𝟙𝟛}}𝕋_{\mathrm{𝟚𝟜}})𝕋_{\mathrm{𝟙𝟛}}]=[𝕕_\mathrm{𝟙}𝕕_\mathrm{𝟚}^\mathrm{𝟚}+𝕕_\mathrm{𝟙}^\mathrm{𝟚}][𝕕_\mathrm{𝟙}𝕕_\mathrm{𝟚}(𝕕_\mathrm{𝟙}𝕕_\mathrm{𝟚}+\mathrm{𝟙})]^\mathrm{𝟙}.`$ This expression inserted in the definition of the entanglement measure proves Eq. (6). $`\mathrm{}`$
For proving bounds on the entangling power (1) it is useful to consider one of the states, say $`|\psi _2,`$ of the input product as fixed. In this case a pair of CP-maps associated with $`U`$ are naturally defined. Indeed one has (explicit dependence on $`U`$ and $`|\psi _2`$ is omitted) $`\mathrm{\Phi }:𝒮(_1)𝒮(_1):\rho _{i=1}^{d_2}A_i\rho A_i^{}`$ and $`\stackrel{~}{\mathrm{\Phi }}:𝒮(_1)𝒮(_2):\rho _{i=1}^{d_1}\stackrel{~}{A}_i\rho \stackrel{~}{A}_i^{}`$ where $`A_i:_1_1`$ and $`\stackrel{~}{A}_j:_1_2`$ are given by $`A_i:=j|U|\psi _2,(j=1,\mathrm{},d_2)\stackrel{~}{A}_i:=_{j=1}^{d_2}|ji|A_j,(i=1,\mathrm{},d_1).`$ Therefore one can also define, for fixed $`|\psi _2_2`$ the (partial) entangling power of $`U`$ as $`\stackrel{~}{e}_p(\mathrm{\Phi }):=\overline{E(\mathrm{\Phi }|\psi \psi |)}^\psi .`$ Notice that the equation above can also be written is the form (3) with a special choice for $`p(\psi _1,\psi _2)`$ i.e., with $`\mathrm{\Omega }_p=\mu (\psi _1)(|\psi _1\psi _1||\psi _2\psi _2|)^{\mathrm{\hspace{0.17em}2}}.`$ The definition of $`\stackrel{~}{e}_p(\mathrm{\Phi })`$ of course makes sense for general CP-maps, in this case the expression for $`\stackrel{~}{e}_p(\mathrm{\Phi })`$ analogous to Eq. (3) is given by $`\stackrel{~}{e}_p(\mathrm{\Phi })=2\text{tr}[\mathrm{\Phi }^{\mathrm{\hspace{0.17em}2}}(\omega _p)P_{13}^{}]`$ where $`\omega _p:=𝑑\stackrel{~}{\mu }(\psi )|\psi \psi |^{\mathrm{\hspace{0.17em}2}}𝒮(_1^{\mathrm{\hspace{0.17em}2}}).`$
Proposition 3 The entangling power of the CP-map $`\mathrm{\Phi }:𝒮(_1)𝒮(_1)`$ with respect to the uniform distribution is given by
$$\stackrel{~}{e}_{p_0}(\mathrm{\Phi })=1C_{d_1}(\text{tr}_2\stackrel{~}{X}^2+\text{tr}_1X^2)$$
(7)
where $`X:=_{j=1}^{d_2}A_iA_i^{}`$ and $`\stackrel{~}{X}:=_{j=1}^{d_1}\stackrel{~}{A}_i\stackrel{~}{A}_i^{}.`$
Proof. One has that $`\text{tr}_1\rho ^2`$ is given by $`_{i,j=1}^{d_2}\text{tr}_1(A_i|\psi _1\psi _1|A_i^{}A_j|\psi _1\psi _1|A_j^{}),`$ this last expression can be rewritten as $`_{i,j=1}^{d_2}\text{tr}_1[(A_j^{}A_i)(A_i^{}A_j)|\psi _1\psi _1|^{\mathrm{\hspace{0.17em}2}}].`$ Now we perform the average with respect $`|\psi _1;`$ using Lemma $`\omega _{p_0}=C_{d_1}(𝟙+𝕋).`$ Using again the identity tr$`[(AB)T]=\text{tr}(AB)`$ one gets
$`C_{d_1}^1\overline{\text{tr}_1\rho ^2}^{\psi _1}={\displaystyle \underset{i,j=1}{\overset{d_2}{}}}|\text{tr}_1(A_j^{}A_i)|^2+\text{tr}_1({\displaystyle \underset{1=}{\overset{d_2}{}}}A_iA_i^{})^2`$ (8)
\[Notice that the two terms in the equation above corresponds to the $`I_\alpha `$’s in Eq. (5).\] It is then straightforward algebra to check that the first term in the equation above can be written as $`\text{tr}_2\stackrel{~}{X}^2`$ $`\mathrm{}`$
We now provide bounds on the entangling power. We assume that $`d_1d_2.`$
Proposition 4 For any $`U𝒰()`$ one has
$$0e_{p_0}(U)\frac{d_2d_2/d_1}{d_2+1}.$$
(9)
Proof. The lower bound is obvious in view of the definition (5), it is achieved by all the unitaries obtained composing bi-local transformations of $`U(d)\times U(d)`$ with the swap. Let us consider first the operator $`X`$ in Eq. (7). By denoting with $`\mathrm{\Phi }`$ the CP-map associated with the $`A_j`$’s one has $`\mathrm{\Phi }(𝟙/𝕕_\mathrm{𝟙})=𝕏/𝕕_\mathrm{𝟙}`$ and then $`1\text{tr}_1\mathrm{\Phi }(𝟙/𝕕_\mathrm{𝟙})^\mathrm{𝟚}\mathrm{𝟙}\mathrm{𝟙}/𝕕_\mathrm{𝟙}`$ (general bound on linear entropy) it follows that $`d_1\text{tr}_1X^2.`$ This latter inequality provides a bound on the second term of Eq. (7). Reasoning in the same way with the operator $`\stackrel{~}{X}`$ and the associated CP-map $`\stackrel{~}{\mathrm{\Phi }}`$ one finds for the first term of (7) the lower bound $`d_1^2/d_2.`$ Putting these two results together, and in view of the assumption $`d_1d_2,`$ inverting $`d_1`$ with $`d_2`$ one gets the desired result (9) $`\mathrm{}`$
Another issue is to understand whether the upper bound in Eq. (9) is achieved by an optimal unitary transformation $`U.`$ As we shall show in the following the answer seems to be affirmative for $`d>2,`$ Let us stress that this is not obvious at all in that the upper bound (9) has been obtained by providing separate bounds on the two terms appearing in $`e_{p_0}(U).`$ Notice that the proof of Proposition 3 allows us to state the condition on $`U`$ in order to saturate the bound (9) as : for any initial state $`|\psi _1`$ the associated CP-maps $`\mathrm{\Phi }`$ and $`\stackrel{~}{\mathrm{\Phi }}`$ (depending both on $`|\psi `$ and $`U`$) must be unital i.e., they map totally mixed states onto totally mixed states. It might well be that no $`U`$’s yields unitality for both CP-maps at once. This in fact turns out numerically to be the case for $`d=2,`$ in which one has that the optimal $`U`$’s are such that \[see Fig. 2\] $`e_{p_0}(U)=2/9<1/3.`$ An optimal operators for qubit is given (not surprisingly) by the controlled-not $`U:=|00|𝟙+|\mathrm{𝟙}\mathrm{𝟙}|𝕏,`$ where $`X:=|01|+|10|.`$ More interestingly the operators providing a natural $`d`$-dimensional generalization of the controlled-not are in general not optimal. This is shown by the following calculation.
Let us consider, for $`d_1=d_2,`$ $`U=_{\alpha =1}^d|\alpha \alpha |U_\alpha ,`$ where the $`|\alpha `$’s are a $`d`$-dimensional orthonormal basis and the $`U_\alpha `$’s are unitaries which, without any loss of generality can be taken to be orthogonal with respect to the Hilbert-Schmidt scalar product i.e., $`<U_\alpha ,U_\beta >=d\delta _{\alpha ,\beta }.`$ By using Eq. (5) is easy to prove that for these unitaries one has $`e_{p_0}(U)=d(d1)/(d+1)^2,`$ that is $`d/(d+1)`$ times smaller than the bound $`(\text{9})`$ and, for $`d>2`$ even smaller.
We also performed numerical maximization of $`e_{p_0}(U)`$ trying to find exact expressions for the optimal unitary transformations. For $`d_1=d_2=d=odd`$ we have found that the following “classical” unitary transformation (which only permutes the $`d^2`$ bases states) is optimal and reaches the bound: $`U|i|j=|i+j|ij`$ where $`i,j=0\mathrm{}d1`$ and the sums are mod$`d.`$ \[Notice that the above expression for even $`d`$ does not define a permutation of the basis of $``$.\] Similarly a more complicated construction gives an optimal permutation that achieves the bound for the case $`d_1=d_2=d=4n`$. Thus for equal dimensions the only case that remains to be solved is $`d_1=d_2=d=4n+2`$, e.g. $`6\times 6.`$ We have also found unitary transformations satisfying the bound for a very asymmetric case, namely $`d_2=nm`$ with $`n,md_1`$. This last example is of the type of controlled unitary operation from the larger to the smaller system. The previous constructions for equal dimensions can be viewed as a concatenation of two such controlled operations, in which the control role is played alternatively by one of the subsystems. It may also be interesting that for the cases 2 $`\times `$ odd the bound can be shown not to be reached by permutations.
Let us finally discuss briefly the numerical evidences. First, even for dimensions $`2\times 3`$ the bound appears not to be reached, rather for optimal $`U`$’s we get the value $`1/3`$ (instead of $`3/8`$). For all other cases that we have checked the bound seems to be reached, namely for $`2\times 4`$ up to $`2\times 7`$, and $`3\times 4`$ up to $`3\times 6`$. In conclusion one might conjecture that the only cases where the optimal transformations do not reach the bound (9) are $`2\times 2`$ and $`2\times 3`$ .
Conclusions. In this paper we introduced a measure for the entangling power $`e(U)`$ of unitary transformations $`U`$ acting on the state-space $``$ of a bi-partite $`d_1\times d_2`$ quantum system. In terms of this measure we moved a first step towards the analysis of the manifold of bi-partite unitary transformations. We analyzed the induced probability distribution $`q(e)`$ over $`𝒰()`$ as $`d_1`$ and $`d_2`$ varies, and we found an analytical form of optimal transformations for some cases. Although we believe that both the questions addressed and the approach we adopted are quite natural and physically motivated the role, if any, that the entangling power will play in Quantum Information theory is still an issue for future work.
The authors thank M. Rasetti and J. Pachos for useful discussions. Ch. Z. is supported by the EU project IST-Q-ACTA.
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# CARBON AND OXYGEN GALACTIC ABUNDANCE GRADIENTS: A Comparative Study of Stellar Yields
## 0.1 INTRODUCTION
In a chemical evolution model that assumes an infall of H and He gas without mass loss by galactic winds or radial flows, the predicted C/O abundance ratio depends mainly on the stellar yields and on the initial mass function. By fixing the initial mass function, this work explores the use of C/O to test stellar yields.
The stellar yield of an element $`X_i`$, being the mass fraction of a star of initial mass $`m`$ converted to $`X_i`$ and ejected to the interstellar medium (ISM), strongly depends on the assumptions of the stellar evolution modeling. By the 80’, no stellar yield set for massive stars ($`m>8`$ M) of different metallicities had been computed. In this decade, Maeder (1992, M92) published He, C, O, and $`Z`$ stellar yields for two initial metallicities ($`Z=0.001`$ and $`0.02`$) for objects with initial mass in the 9-120 M range. Later on, Woosley & Weaver (1995, WW) calculated stellar yields for 80 elements and isotopes for stars with initial masses between 11 and 40 M and with five initial metallicities ($`Z/Z_{}=`$ 0, $`10^4`$, 0.01, 0.1, and 1), without initial abundance ratios scaled to solar. More recently, Portinari, Chiosi, & Bressan (1998, PCB) published stellar yields of 17 elemental species for a wide range of metallicities ($`Z=0.0004`$, 0.004, 0.008, 0.02, and 0.05) and masses (6-120 M).
For low and intermediate mass stars (LIMS, $`0.8<m/\mathrm{M}_{}<8`$), stellar yields have been calculated by Renzini & Voli (1981, RV), van den Hoek & Groenewegen (1997, HG), and Marigo, Bressan, & Chiosi (1996, 1998, MBC). While HG produce a grid of stellar yields for five metallicities ($`Z=`$0.001, 0.004, 0.008, 0.02, and 0.04), RV and MBC sampled yields for only a pair of metallicities ($`Z=`$0.004, 0.02 and $`Z=`$0.008, 0.02 respectively). A combination of PCB and MBC provides the only complete grid of stellar yields computed from the same evolutionary tracks (Padova group).
Carigi (1994) has shown that the evolution of \[C/O\] with \[O/H\] in the solar vicinity can be explained by M92 yields due to the dependence of the C and O yields on initial $`Z`$. Prantzos, Vangioni-Flam, & Chauveau (1994) concluded that the metallicity dependent M92 yields are also able to reproduce the growth of \[C/O\] in the solar neighborhood. Extending the modeling beyond the solar vicinity, this work follows the \[C/O\] abundance and evolution with chemical evolution models for the galactic disk, with the aim to study the behavior of differences of sets of metal-dependent yields from both, massive stars and LIMS.
In general, chemical evolution models of the solar vicinity and the Galactic disk assume a given yield set, either M92 yields (e.g. Giovagnoli & Tosi 1995, Prantzos & Aubert 1995, Carigi 1996), WW yields (e.g. Timmes, Woosley, & Weaver 1995, Chiappini, Matteucci & Gratton 1997, Allen, Carigi, & Peimbert 1998, Prantzos & Silk 1998), or Padova yields (e.g. Portinari et al. 1998, Tantalo et al. 1998). Since these yield sets have substantial differences, it is not straightforward to separate the effects of yields from other model assumptions when intercomparing the different studies. In this work, consistent chemical evolution models are build to find out which yield sets reproduce better the ISM abundances. The predicted abundances from the different sets of stellar yields are compared for the four elements in common among all sets (H, He, C, and O).
In this paper, all chemical abundances are presented by number, with the exception of $`\mathrm{\Delta }Y/\mathrm{\Delta }O`$, and $`\mathrm{\Delta }Y/\mathrm{\Delta }Z`$, that are given by mass.
The paper has been organized as follows: The observational constraints on the models are presented in §2. §3 describes the assumptions adopted in the chemical evolution models. The model predictions are then shown and briefly discussed in in §4. The global discussion and conclusions are presented in §5 and §6, respectively.
## 0.2 OBSERVATIONAL CONSTRAINTS
In this work, the compilation by Peimbert (1999) and the recent data of Esteban et al. (1998, 1999a, 1999b) are used as abundance constraints for the models. He/H, C/H, and O/H values, for the Galactic H ii regions M17, M8, and Orion are taken from Peimbert’s Table 1, but corrected for dust depletion, adding 0.08 dex and 0.10 dex to the O and C abundances respectively, as suggested by Esteban et al. (1998). In the modeling, I attempt to fit the abundance gradients to be within the average values by Peimbert (1999) and those obtained from the observations by Esteban et al. (1998, 1999a, 1999b), considering temperature fluctuations ($`t^2>0.00`$). Average O/H gradients computed by Peimbert (1999) are based on the gradients from Esteban et al. (1998, 1999a, 1999b), Shaver et al. (1983), and Deharveng et al. (1999). Since in the literature there are C/O values based on recombination lines only for M17, M8, and Orion, the average C/H gradient is obtained from Esteban et al. (1998, 1999a, 1999b) and Peimbert et al. (1992). The relative enrichment $`\mathrm{\Delta }Y/\mathrm{\Delta }O`$ and $`\mathrm{\Delta }Y/\mathrm{\Delta }Z`$ values for M17, M8, and Orion are taken from Peimbert (1999, $`t^2>0.00`$), which have already been corrected for dust (also as suggested by Esteban et al. 1998).
O/H abundances and gradient from B-stars are in agreement with those from H ii regions (Gummersbach et al. 1998, Smartt & Rolleston 1997), but the C/H values are lower by at least 0.3 dex, and the C/H gradient is twice as flat (Gummersbach et al. 1998). The gradients from B-stars shown in Table 3 were computed using the data of Gummersbach et al. (1998) and Smartt & Rolleston (1997) for the galactocentric range considered here ($`4<r<10`$ kpc). The innermost B-star of both data sets was eliminated from the fit, given their uncertain abundances as discussed by the respective authors. The He/H data by Gummersbach et al. (1998) are not considered in this work given that the large dispersion in abundances ($`1.1<\mathrm{log}(\mathrm{He}/\mathrm{H})<0.5`$, corresponding to $`14<`$ $`\mathrm{\Delta }Y/\mathrm{\Delta }O`$$`<150`$) does not provide a tight enough constraint for the chemical evolution modeling. The large differences in He abundance observed among B-stars may come from a contamination in the outer stellar layers by helium produced by the stars, combined with errors and other uncertainties in the abundance determinations. Although B-stars and H ii regions show a different slope of the C/O gradient, both of them are negative. In this study, the presence (not the slope magnitude) of a negative C/O gradient in the galactic disk is used as another observational constraint.
The model C/O history for the solar vicinity is constrained with C/O abundances from dwarf stars located 1 kpc around the Sun (Gustafsson et al. 1999) and the solar value from Grevesse & Sauval (1998). Since C/H values from B-stars are lower than those from H ii regions, two different present-day C/O values are presented. The increase of C/O with metallicity in the solar vicinity provides yet another observational constraint.
Planetary nebulae abundances are not be used here as constraints because C is produced by their progenitors. The models by Renzini & Voli (1981), van den Hoek & Groenewegen (1997), and Marigo et al. (1996, 1998) predict C enrichment in the envelope during the evolution of stars with masses lower than $``$ 8 M. Moreover, Peimbert, Torres-Peimbert, & Luridiana (1995) and Peimbert, Luridiana, & Torres-Peimbert (1995) found C/H values higher than the solar value, by at least 0.1 dex, for the vast majority of planetary nebulae in their sample.
Another important observational constraint is gas consumption, measured as the ratio of gas to total surface mass densities, $`\sigma _{gas}/\sigma _{tot}`$. In this work the observed $`\sigma _{gas}`$ distribution, as compiled by Matteucci & Chiappini (1999), is used together with an exponential $`\sigma _{tot}`$ distribution with a 3 kpc scale-length and a $`\sigma _{tot}(r_{})=45M_{}pc^2`$ amplitude (Kuijken & Gilmore 1991). A galactocentric distant for the Sun, $`r_{}`$, of 8 kpc is adopted here.
Summarizing, the just-described observational constraints are used in this study as follows: a) All models (for each and every yield set) are build to exactly reproduce: i) the observed gas fraction distribution of the galaxy, $`\sigma _{gas}/\sigma _{tot}`$; and ii) the observed O/H galactic gradient. b) In order to study the differences among available yields sets, the model predictions are then compared to: i) the observed rise of C/O with metallicity (or equivalently with time) in the solar neighborhood; and ii) the observed decrease of the C/O abundance with galactocentric distance.
## 0.3 CHEMICAL EVOLUTION MODELING
The scope of this work is to explore models assuming the different sets of metal-dependent stellar yields available in the literature.
All models are built to reproduce two observational constraints: the O/H abundance gradient from H ii regions and B-stars, and the $`\sigma _{gas}`$ distribution from 4 to 10 kpc. Models are very similar to the infall model of Allen, Carigi, & Peimbert (1998), but adopting somewhat different assumptions:
a) The star formation rate is set proportional to a power of $`\sigma _{gas}`$ and $`\sigma _{tot}`$: $`SFR(r,t)=\nu \sigma _{gas}^x(r,t)\sigma _{tot}^{x1}(r,t)`$, where $`\nu `$ and $`x`$ are constant in time and space, and fixed in such away that the observational constraints are reached after 13 Gyr, the age of the model.
b) Three sets of metal-dependent stellar yields from massive stars are used: i) Geneva yields: M92 for $`9m/\text{M}\text{}85`$ (high mass-loss rate); ii) Santa Cruz yields: WW for $`11m/\text{M}\text{}40`$ (models “B” for 30, 35 and 40 M), and WLW for $`m=60\text{M}\text{}`$ and $`m=85\text{M}\text{}`$; Since $`Z_{max}=Z_{}`$ for the Geneva and Santa Cruz groups, then yields ($`Z>Z_{}`$) = yields($`Z_{}`$). iii) PCB yields: Portinari et al. (1998) for $`9m/\text{M}\text{}85`$. In all cases, effects due to black hole formation are not considered.
c) Three sets of metal-dependent stellar yields for LIMS are used: i) RV yields: Renzini & Voli (1981) for $`1m/\text{M}\text{}3`$ ($`\alpha =0.0`$,$`\eta =0.3`$, case A), for $`3m/\text{M}\text{}8`$ ($`\alpha =1.5`$, $`\eta =0.3`$, case A). ii) HG yields: van den Hoek & Groenewegen (1997) for $`0.8m/\text{M}\text{}8`$ ($`\eta _{AGB}=0.4`$,$`M_{HBB}=0.8`$). iii) MBCP yields: Marigo et al. (1996) for $`0.8m/\text{M}\text{}3`$; Marigo et al. (1998) for 4, 4.5, and 5 M; Portinari et al. (1998) for 6 and 7 M.
d) The fraction of binary systems predecessors of SNIa is taken as $`A=0.07`$. Contrary to Carigi (1994), SNIb produced by binary systems are now not considered.
Carigi (1996) and Allen, Carigi, & Peimbert (1998) successfully modeled galactic abundance gradients considering the IMF by Kroupa, Tout, & Gilmore (1993). Therefore, in this work the same IMF is assumed for a 0.01 to 85 Mmass range. Moreover, this IMF is maintained constant in time and space, based on the conclusion of Chiappini, Matteucci, & Padoan (2000) that a constant IMF is still the best assumption to explain the observational constraints in the Milky Way.
In the present exercise, to avoid spurious mass extrapolation effects for a fair comparison of the predictions from the different yield sets, the stellar mass range was limited so as to be lower than 85 M(the maximum mass common to the three sets from massive stars).
Since the code does not use the instantaneous recycling approximation, it is necessary to estimate, at each time step, the stellar yields of all dying stars at the metallicities at their birth times. Therefore, the stellar yields from the different groups have been interpolated in both mass and metallicity, taking in consideration their different samplings. In particular, given the number of initial metallicities considered by stellar evolution models, and the yields $`Z`$-dependencies, I have assumed: a) a linear interpolation with mass for yields of massive stars and LIMS; b) a linear interpolation with metallicity for the Santa Cruz, PCB, RV, HG and MBCP yields; and c) three kinds of $`Z`$ interpolations (linear, as $`Z^{0.5}`$, and $`Z^2`$) for the Geneva yields. It is important to remark that the HG yields vary almost linearly with $`Z`$, but the $`Z`$-behavior of yields of massive stars is not well determined: the PCB yields dependency on $`Z`$ is not simple and the Santa Cruz yields are almost independent of $`Z`$.
Input parameters for the set of models are summarized in Table 1. Column (2) lists $`x`$, the SFR power on $`\sigma _{gas}`$, and column (3) the efficiency factor, $`\nu `$, of each model. The slope of the final abundance gradients is basically governed by $`x`$; while $`\nu `$ and $`x`$ together determine the O abundance. Models are divided in three groups (column 4), assuming M92, WW&WLW, or PCB yields. M92 and WW&WLW yields have been complemented with RV, HG, and MBCP yields. Three subgroups of models have been computed according to the type of interpolation used ($`Z^n`$, $`n=`$0.5, 1, or 2) between the two metallicities considered by Maeder (1992). The table lists just one model for the PCB yields (C-model), since the Padova yields (PCB+MBCP) form a complete and consistent set.
Since LIMS do not produce O, each set of yields from massive stars requires different values of $`x`$ and $`\nu `$ to match the same observational constraints. Since the M92-O yields for $`m>25`$ M decrease with metallicity, O gradients saturate quicker, calling for higher $`x`$ values. On the other hand, WW-O yields are almost independent of metallicity, and WLW-O yields are lower than M92-O yields at low metallicity, forcing higher $`\nu `$ values for the WW&WLW modeling to reach the observed O/H of H ii regions. Since PCB-O yields are lower than M92-yields at low $`Z`$, it is required that $`\nu _{\mathrm{C}\mathrm{model}}\nu _{\mathrm{M}\mathrm{models}}`$ when the interpolation is weighted to high metallicities.
Since the yields for supermetallic stars are similar to those by WW, then $`x_{\mathrm{C}\mathrm{model}}\nu _{\mathrm{W}\mathrm{models}}`$.
Since infall models of the present kind produce abundance gradients that flatten out with time (Carigi 1996), the M-models require higher $`x`$ exponents to reproduce the O/H gradient at a fixed age (13 Gyr).
## 0.4 MODEL RESULTS
All models match well several observational constraints at the solar vicinity, e.g. age-metallicity relation, G-dwarf distribution, \[C/Fe\] and \[O/Fe\] vs \[Fe/H\] evolution. For a SFR$`=\nu \sigma _{gas}^x\sigma _{tot}^{x1}`$ law, the present models adopt $`x`$ in the range 1.13 and 1.22. Predictions for the solar vicinity do not change significantly when the $`x`$ exponent varies from 1 to 2 (Carigi 1996). This work does not then discuss how well the observed relations are matched, since the predicted relations are very similar to those of the best model of Carigi (1996), where her Fig 4 already summarizes the comparison with the observational data.
### 0.4.1 C/O evolution of the solar vicinity assuming different yields
Figure 1 shows the C/O evolution in the solar vicinity as predicted by models that consider yields by Maeder (1992). Since the dependence of these yields for massive stars with $`Z`$ is not well known, the influence of the assumed type of interpolation on on the C/O ratio is explored in Figure 1a. The figure shows that, between 1.5 and 9 Gyr, the C/O rise in time slows down as the interpolation exponent $`n`$ increases ($`Z^n`$); out of this lapse, the increase is basically the same for all reasonable $`n`$s. It is then important to compute yields for intermediate metallicities to reduce these interpolation ambiguities. Besides interpolations effects, Figure 1b also explores the influence of LIMS on the C/O evolution, since these stars are important carbon producers. After 1 Gyr or so, RV predict less carbon production than MBC but more than HG.
The best models from the M92 yields among the different combinations of LIMS yields and interpolation types assumed (M-models = MRV1, MHG1, MMBC2) are shown in Figure 2. Results from models assuming the Santa Cruz and Padova yield sets are also presented in similar diagrams (Figure 3 and Figure 4 respectively). These results, from the modeling assuming the different yield sets, are discussed in the next two subsections, looking first at the C/O evolution in the solar vicinity, and later on at the local chemical distributions along the galactic disk.
In particular, the (h) panel of each figure presents the predicted C/O chemical enrichment history of the solar neighborhood, as compared to the following observational data: i) the observed abundance in the Sun and in the neighboring disk dwarf stars, ii) a C/O value at $`r_{}`$ derived from the observed gradient for H ii regions, and iii) a C/O average value for the nearest B-stars from Gummersbach et al. (1998). From these panels, one can see that the C/O abundance always increases with time in the M-models and the C-model, in agreement with most of the observations (dwarf stars, Sun, and H ii regions), but the results of the W-models fail to predict the observed recent rise in C/O.
The discrepancy among the C/O predictions with different yields can be understood. The M92 and PCB-O yields increase with mass and decrease with metallicity, while the C yields increase with both mass and metallicity for $`ZZ_{}`$. At the onset of the evolution, the WLW yields from very massive stars ($`m>40`$ M) enrich the ISM with ever more oxygen through O-yields that increase with mass. Once these stars die after a few million years, the C/O of the W-models behaves like in the M-models because, for $`m<40`$ M, both sets of yields are similar at low metallicities. Since the WW yields do not depend on metallicity, the W-models reach (at $``$ 3.5 Gyr) a plateau in C/O that decreases slowly with time.
The C-model predicts two plateaus in the C/O evolution of the solar neighborhood: a first one at $``$ 5 Gyr and another after $``$ 11 Gyr. The later one is due to the contribution of supermetallic stars. These stars have extreme winds that occur before C and O are synthesized, so their winds are He rich, while C and O are produced during the SN explosion. These characteristics of the yields from supermetallic stars mimics the behavior predicted by WW.
Table 2 summarizes the present-day predictions for the solar neighborhood from all models. M-models reproduce the observed C/H and C/O values from H ii regions. WRV and WHG models predict C/H and C/O values in agreement with B-stars, while the WMBC model results agree with the C/H and C/O from H ii regions.
Inspecting Figure 1 and Table 2 it can be noticed that the average C/O for B-stars is lower than for H ii regions and for dwarf stars, mainly due to the low C/H values observed in B-stars.
### 0.4.2 Chemical radial distributions in the local disk assuming different yield sets
Predictions for the Galactic disk are also summarized in the (a) to (g) panels of Figures 2-4, together with Table 3. From these figures and table it can be seen that:
a) The C/H gradients predicted by the best models with M92 yields are similar to the observed ones from H ii regions. The gradients from the W-models (considering WW&WLW yields) or the C-model are flatter by about 0.04 dex kpc<sup>-1</sup> and are close to the gradient traced by B-stars;
b) C/H abundance ratios from the M, WMBC, and C models agree with the H ii regions observations, but those from WRV and WHG models are lower (by at least 0.1 dex) in closer agreement with the C/H of B-stars;
c) M-models are the only ones that reach the high C/H and C/O values of M17;
d) The negative C/O gradients predicted by the M-models are in very good agreement with the observations. The gradients predicted by W-models or C-model, being positive, are a poor match to the observed gradient;
e) The C/O values from W-models agree with those determined from B-stars, while the C/O values from M and C models agree with the C/O values determined from H ii regions;
f) M-models and the WMBC-model reproduce the observed He/H abundances in H ii regions, with the former yielding the best agreement;
g) RV yields predict lower C and He than MBC but higher than HG;
h) M-models predict strong negative $`\mathrm{\Delta }Y/\mathrm{\Delta }O`$ gradients while W-models and C-model yield almost no gradient;
i) Predicted $`\mathrm{\Delta }Y/\mathrm{\Delta }Z`$ gradients are almost flat for all models;
j) With the exception of the WMBC model, all models predict helium to heavy-elements abundance ratios in the $`1.6<`$ $`\mathrm{\Delta }Y/\mathrm{\Delta }Z`$$`<1.0`$ range, values that are lower than those observed in M8 and Orion.
The C/O gradients from the W-models and the C-model are almost flat because WW yields are almost independent of $`Z`$ and the very metal-rich stars eject C and O in SN explosions and not in winds, both facts imply a saturation in C/O.
The He/H gradients predicted by the M-models are steeper than by the C-model because the very-massive and super-metallic stars eject a smaller amount of He than solar metallicity stars (PCB)
### 0.4.3 C/O vs O/H
In Figure 5 the predicted evolutions of C/O versus O/H for the solar neighborhood are presented. In addition to Galactic data, C/O and O/H values for extragalactic H ii regions are shown. Objects in the halo and in the disk of the solar neighborhood (cylinder of radio 1 kpc centered in the Sun), and in external galaxies call for an increase of C/O with O/H enrichment. Only models that reproduce this increase are shown in the figure, so the W-models are excluded. The agreement for log(O/H) $`<`$ -3.8 is not so good, with the exception of MHG1. Again, here one can notice that C yields for intermediate mass stars given by HG are lower than those obtained by RV or MBCP.
The log(C/O)-log (O/H) relation predicted by the MHG1-model of this work (Fig 5) is not surprisingly very similar to the one obtained by Henry, Edmunds, & Köppen (2000), since their best models, either analytical or numerical, are constructed also adopting the yields of Maeder (1992) and van den Hoek & Groenewegen (1997).
A powerful tool for comparing chemical evolution models with interstellar medium abundances is given by the relation
$$\mathrm{log}(\mathrm{C}/\mathrm{O})=a\mathrm{log}(\mathrm{O}/\mathrm{H})+b.$$
(1)
Table 4 presents the $`a`$ values predicted by the models, considering different concentric galactic rings: i) a very local disk, between $`r=6`$ and $`r=8`$ kpc, and ii) the complete area studied in this paper. Moreover, this table shows the $`a`$ values derived from four sets of observations, two galactic and two extragalactic ones; for B-stars it was assumed that the C and O abundances correspond to those of the interstellar medium. The observed $`a`$ values indicate that the negative C/O local gradient may be extended to the whole Galactic disk. As can be seen from the table, the observed $`a`$ values are in good agreement with the M-models and in disagreement with the W-models and the C-model, reiterating the results presented in Figure 5. According to the discussion of Figure 5, the best models are: MRV1, MHG1, and MMBC2.
## 0.5 DISCUSSION
The determinations of C/O from H ii regions in spiral and irregular galaxies (Garnett et al. 1999), H ii regions in our Galaxy (Peimbert, Torres-Peimbert & Ruiz 1992, Esteban et al. 1998, 1999a, 1999b), and dwarf stars in the solar vicinity (Gustafsson et al. 1999); all converge to the conclusion that C/O must increase with metallicity. M92 yields can explain this behavior, Padova yields only partially, while WW&WLW yields not at all. Other authors (Ferrini, Matteucci, Prantzos, Timmes, and Tosi, see Tosi 1996 for references) using different chemical evolution models, have predicted flat C/O gradients (Tosi 1996) because all of them have adopted either yields at a fix $`Z`$ or the WW yields, which are almost independent of $`Z`$.
Carigi (1994) and Prantzos et al. (1994) had already concluded that the evolution of \[C/O\] with \[O/H\] or \[Fe/H\], respectively, in the solar vicinity can be explained by the metal dependent yields of Maeder (1992). The present work have extended these previous ones to the Galactic disk using the recent C determinations in H ii regions and B-stars as observational constraints. Moreover, seven sets of yields have been considered here, composing the complete sample of stellar yields dependent on the initial stellar metallicity available in the literature.
A main difference among these sets of yields resides on the stellar-wind assumptions. WW do not consider stellar winds at any stage, and M92 and PCB assume a mass-loss rate proportional to $`Z^{0.5}`$ and $`m^{2.5}`$ during the post-main-sequence phases.
C and O yields from M92, PCB, and WW are similar for $`m<25`$ M, where stellar winds are not so relevant. In M92 and PCB, most of the net C is produced before the SN stage, and since the wind strength increases with metallicity, the C yield increases with $`Z`$. On the other hand, O is mainly ejected during the SN stage and, since the amount of O expelled by the SN depends on the stellar mass just before the SN explosion, the M92 and PCB-O yield decreases with $`Z`$. This fact causes the M-models and the C-model to reproduce the rise of the C/O in the solar vicinity.
The increase of C yields and the decrease of O yields with $`Z`$ can not be extrapolated to $`Z>Z_{}`$. According to PCB, the C to O yields ratio decreases from $`Z=Z_{}`$ to $`Z=2.5Z_{}`$. Supermetallic and massive stars have early and intense winds, which occur before these stars synthesize C and O. Therefore, their C and O yields are similar to those of metal-poor stars that mainly eject C and O in the SN explosion. This behavior of the C and O yields produce a saturation of the C/O abundance ratio at $``$ -0.1 dex, which produces flatter gradients for the C-model. Since M-models assume that the supermetallic stars behave like metal-solar stars, the C-model would predict a negative C/O gradient if the mass-loss rate for $`Z=2.5Z_{}`$ is similar to that for $`Z=Z_{}`$
It is important to note that in the W-models the wind contribution to the C and O yields was not included, since WLW present information to calculate only the He expelled by winds. If winds were considered in these models, the C yield would certainly be higher but, given the low weight at the high-mass end of the initial mass function, hardly enough to reproduce the present-day C/O abundances. When the WLW yields are not taken into account (by reducing the upper mass-limit to 40 M) the predicted C/O decreases at times earlier than 0.4 Gyr, but still remains low and basically constant for the rest of the evolution (Carigi & Peimbert 2000).
The M-models and C-model agree better with the younger objects in the solar vicinity than with the older ones. The agreement would improve if metal-poor stars eject more oxygen and less carbon than predicted by a $`Z^{0.5}`$ law. The M-models are in agreement with data for spiral and irregular galaxies while the W-models are not, as discussed by Garnett et al. (1999). Then again, the general agreement would be even better if the C and O yields for metal-poor stars were more dependent on metallicity than assumed by M92 and PCB. Furthermore, the C/O value of I Zw 18 determined by Garnett et al. (1995, C/O$`=0.60\pm 0.10`$ dex) would be hard to explain (not so the lower value of Izotov & Thuan 1999, $`<\mathrm{log}(\mathrm{C}/\mathrm{O})>=0.78\pm 0.03>`$). It should be noted that the models presented in Figure 5 were made to reproduce the Galactic disk and therefore, the comparison of the extragalactic H ii regions and halo objects is only indicative. Specific models for each galaxy or for the Galactic halo should be carried out, see for example the models by Carigi, Colín, & Peimbert (1999) and Carigi & Peimbert (2000).
Moreover, the M-models and the C-model agree well with the \[C/H\], \[C/O\], and \[C/Fe\] vs \[Fe/H\] relations observed in dwarf stars of the solar vicinity by Gustafsson et al. (1999), after correcting for the adopted solar abundances (Grevesse & Sauval 1998 as opposed to Anders & Grevesse 1989). The predicted relations are matched at high \[Fe/H\] and become slightly lower than those observed at lower metallicities; again, a behavior that calls for a stronger metal dependency of the C and O yields for metal-poor stars than a $`Z^{0.5}`$ law. W-models basically match the relation slopes, but again predict a lower C abundance than observed (by $``$ 0.3 dex in C/H and $``$ 0.2 dex in C/Fe).
Maeder (1993), Woosley & Weaver (1995) and Portinari et al. (1998) have assumed different values for the C<sup>12</sup>($`\alpha `$,$`\gamma `$)O<sup>16</sup> reaction rate. It is known that this rate is uncertain by about a factor of 2. This uncertainty causes minor differences in the behavior of C/O as a function of metallicity predicted by the models, but more severely affects the absolute C/O value. This change in the absolute C/O value does not alter the conclusions of this paper. In other words, to a very good approximation, a change in the C<sup>12</sup>($`\alpha `$,$`\gamma `$)O<sup>16</sup> rate modifies $`b`$ in equation (1), but not $`a`$.
The C/H values derived from Galactic B-stars are about 0.3 dex lower than those from Galactic H ii regions and from disk dwarf stars of the solar vicinity. Carbon abundances of B-stars are not easy to understand and may not be representative of the present-day C in the ISM, because: i) being so young ($`10^3`$ \- $`10^2`$ Gyr), B-stars would be expected to be richer in C than most dwarf stars, but the observed C/H in B-stars in the solar vicinity is very similar to the C of the poorest dwarfs; ii) the O and N abundances and gradients from B-stars and from H ii regions are very similar (surprisingly so for N), so it is difficult to understand why the C/O gradient from B-stars is almost flat, while the one from H ii regions is not; iii) a real C/O gradient is supported by the fact that other two nearby spiral galaxies, M101 and NGC 2403, also show significant C/O gradients, -0.04 and -0.05 dex/kpc respectively (Garnett et al. 1999); and iv) C/O from H ii regions, the Sun, and dwarf stars are consistent with each other, but not with the C/O from B-stars in the solar vicinity.
It is known that intermediate mass stars ($`m<8`$ M) produce carbon and not oxygen. Carigi (1994) found that yields by Renzini & Voli (1981) can not explain the C/O increase with the metallicity in the solar neighborhood. Prantzos et al. (1994) suggested that C-yields higher that those predicted by Renzini & Voli could reproduce the observed rise in C/O. The other two sets of yields for LIMS dependent on $`Z`$ that exist in the literature (van den Hoek & Groenewegen 1997, and Marigo et al. 1996, 1998) can not explain the recent rise of C/O in the solar vicinity. However, LIMS have an important role in the early evolution of C/O.
## 0.6 CONCLUSIONS
Based on basic chemical evolutionary models for the Galactic disk, assuming different sets of stellar yields dependent on metallicity, the following conclusions are reached:
a) Different sets of stellar yields predict different C/O values. The recent rise of C/O with O/H is due mainly to massive stars. The early rise of C/O with O/H is due to both massive stars and LIMS.
b) Models using stellar yields with stellar winds dependent on $`Z`$ (Maeder 1992, Portinari et al. 1998) reproduce the rise of C/O with time shown by H ii regions, the Sun, and dwarf stars within the solar vicinity.
c) Models considering Maeder’s yields are also successful in reproducing the C/O Galactic abundances and gradient determined from H ii regions.
d) The model assuming Padova yields reproduces C/O abundances but not C/O Galactic gradients from H ii regions, because stellar winds $`Z^{0.5}`$ become too strong for supermetallic stars.
e) Models assuming yields from WW&WLW reproduce only the C/O abundances from B-stars, but fail to reproduce the other observational constraints.
f) Since the C/H abundances in B-stars are lower than in disk dwarf stars of different ages in the solar neighborhood, while the O/H value is very similar, the C/O value from these B-stars may not be a good observational constraint.
g) The main difference between the sets of yields of massive stars arises because Geneva and Padova groups consider stellar winds dependent on metallicity while WW do not, and WLW only partially. One can then conclude that $`Z`$-dependent stellar winds must play an important role in the chemical enrichment history of the Galaxy.
h) Observations within the solar neighborhood, the Galactic disk, as well as in spiral and irregular galaxies imply that C/O must increase with metallicity. Modeling including not only winds, but actually metal-dependent winds, is necessary to properly follow the chemical evolution of galaxies.
i) To improve the agreement with the C/O Galactic abundances and the C/O evolution with metallicity, the present models call for a more complicated mass-loss rate law than $`Z^{0.5}`$, assumed by Maeder (1992) and by Portinari et al. (1998): $`\dot{M}_{wind}Z^n`$, such that $`n>0.5`$ if $`ZZ_{}`$, $`n0.5`$ when $`ZZ_{}`$ and $`0.5<n<0.7`$ for $`Z>Z_{}`$.
j) Models based on yields by Renzini & Voli (1981) predict less C and He than those by Marigo et al. (1996, 1998) and more than those by van den Hoek & Groenewegen (1997)
k) The $`\mathrm{\Delta }Y/\mathrm{\Delta }Z`$ value decreases along the sequence of models MBCP - RV - HG. At the same time, it increases along the PCB - M92 - WW&WLW model sequence.
l) The number of H ii regions with known C abundance is quite small, and the C/O gradient from H ii regions might change with future C determinations in more H ii regions. It is important to obtain C abundances from H ii regions located at other galactocentric distances to determine the Galactic C/O gradient with higher accuracy.
###### Acknowledgements.
I would like to thank Manuel Peimbert for valuable comments and useful suggestions. I also wish to thank Laura Portinari for providing me the massive stars yields of the Padova group. I am grateful to Don Garnett for supplying me the data for irregular and spiral galaxies. I also acknowledge the several excellent suggestions by the referee which have improved this final version. I thank Jesús González for a thorough reading of the manuscript.
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# 1 Abstract
## 1 Abstract
Derndinger \[Der80\] and Krupa \[Kru90\] defined the F–product of a (strongly continuous one–parameter) semigroup (of linear operators) and presented some applications (e.g. to spectral theory of positive operators, cf. \[EN00\]). Wolff (in \[Wol84\] and \[Wol00\]) investigated some kind of nonstandard analogon and applied it to spectral theory of group representations. The question arises in which way these constructions are related.
In this paper we show that the classical and the nonstandard F–product are isomorphic (theorem 4.6). We also prove a little “classical” corollary (4.7).
## 2 Basic Notation
### 2.1 Semigroups:
Let $`E`$ be a (real or complex) Banachspace. We denote the norm of an element $`f`$ of $`E`$ by $`|f|`$. $`(E)`$ is the set of bounded linear functions from $`E`$ to $`E`$. The elements of $`(E)`$ are called bounded linear operators, and the norm of an operator $`A`$ is is denoted by $`A`$. A one–parameter semigroup of bounded linear operators (or semigroup, for short) is a function $`T:_0^+(E)`$ such that $`T(0)=\text{Id}_E`$ and $`T(t_1+t_2)=T(t_1)T(t_2)`$. A semigroup $`T`$ is called strongly continuous (or continuous, for short), if for all $`fE`$, $`lim_{t0}|T(t)ff|=0`$ (i.e. $`lim_{t0}T(t)=\text{Id}_E`$ in the strong topology). $`T`$ is called uniformly continuous, if $`lim_{t0}T(t)\text{Id}_E=0`$ (i.e. $`lim_{t0}T(t)=\text{Id}_E`$ in the topology induced by the operator–norm).
For a continuous semigroup $`T`$ and $`fE`$ we define $`A(f):=lim_{h0}\frac{1}{h}(T(h)ff)`$, if this limit exists (in $`E`$). $`D(A)`$ is the set of $`fE`$ such that $`A(f)`$ exists. $`A`$ is called the Generator of $`T`$.
More about the (classical) theory of continuous semigroups can be found in the textbook \[EN00\]. We will only need the following results:
###### Lemma 2.1.
1. For every continuous semigroup $`T`$ there exist constants $`M1`$, $`\omega `$ such that for all $`t_0^+`$: $`T(t)Me^{\omega t}`$.
2. $`D(A)`$ is a dense linear subspace of $`E`$, and $`A`$ is a closed linear operator (i.e. the graph of $`A`$ is a closed subset of $`E\times E`$).
3. $`D(A)=E`$ if and only if $`T(t)`$ is uniformly continuous.
4. If $`fD(A)`$ and $`h>0`$, then $`|T(h)ff|h|A(f)|sup_{sh}T(s)`$
5. If $`fD(A)`$ and $`h>0`$, then $`|\frac{T(h)ff}{h}A(f)|sup_{sh}(|T(s)A(f)A(f)|)`$
6. For any $`\overline{f}`$ and $`h>0`$ there is a $`fD(A)`$ such that $`A(f)=\frac{T(h)\overline{f}\overline{f}}{h}`$ and $`|\overline{f}f|sup_{sh}(|T(s)\overline{f}\overline{f}|)`$
Remark: The last three items of the lemma follow from the following two facts:
If $`fD(A)`$ and $`t>0`$, then $`T(t)ff=_0^tT(s)Af\text{d}s`$.
For all $`fE`$ and $`t>0`$, $`_0^tT(s)f\text{d}s`$ is in $`D(A)`$, and $`A(_0^tT(s)f\text{d}s)=T(t)ff`$.
(The integral is defined as the limit of the Riemann–sums.)
Using $`|f(s)\text{d}s||f(s)|\text{d}s`$ the items 4 and 5 follow directly from $`T(t)ff=_0^tT(s)Af\text{d}s`$. We get item 6 by defining $`f:=\frac{1}{h}_0^hT(s)\overline{f}\text{d}s`$.
The famous theorem of Hille–Yosida states that uniformly continuous semigroups are exactly the semigroups of the form $`T(t)=e^{tA}`$, where $`A`$ is in $`(E)`$. If $`T(t)=e^{tA}`$, then $`A`$ is the generator of $`T`$, and $`D(A)=E`$. In this case, the constructions in this paper do not result in anything new. So we will mainly be interested in semigroups that are not uniformly continuous.
Our basic example for such a semigroup is the following:
###### Example 2.2.
$`C^\text{b}`$ is the (real) Banachspace of uniformly continuous, bounded functions from $``$ to $``$ (with the sup–norm, denoted by $`_{\mathrm{}}`$). $`T`$ is the translation semigroup, defined by $`T(t)f(x):=f(x+t)`$.
Then $`T`$ is continuous (since each $`fC^\text{b}`$ is uniformly continuous), but $`T`$ is not uniformly continuous: for example, let $`f_k(x)=\mathrm{sin}(kx)`$. Then $`f_kC^\text{b}`$ and $`|f_k|=1`$. For all $`t>0`$ there is a $`k`$ such that $`|T(t)f_kf_k|>1`$, i.e. $`T(t)\text{Id}_E>1`$. $`f`$ is in $`D(A)`$ if and only if $`f^{}`$ exists and is element of $`E`$. In this case, $`A(f)=f^{}`$.
### 2.2 Robinsonian (nonstandard) analysis:
We present the concept of nonstandard extensions following \[CK90\].
Let $`X`$ be a set. $`P(X)`$ denotes the set of all subsets of $`X`$. $`V_0(X)=X`$, $`V_{n+1}(X)=V_n(X)P(V_n(X))`$, $`V(X)=_{n\omega }V_n(X)`$. We assume that (a suitable copy of) $``$ exists in $`X`$, and that all vector spaces we are interested in are subsets of $`X`$. So all sets of numbers, subspaces of $`E`$ etc are in $`V(X)`$. For technical reasons, we want to treat the elements of $`X`$ as urelements, i.e. we want to avoid that for some $`yV(X)`$ and $`xX`$, $`yx`$. We can do that by assuming (without loss of generality) that $`X`$ is a base–set, i.e. $`X`$ does not contain the empty set and $`xV(X)`$ is empty for all $`xX`$.
A nonstandard universe is a triple $`V(X),V(Y),^{}:V(X)V(Y)`$ such that:
* $`X`$ and $`Y`$ are infinite base–sets, $`XY`$, $`{}_{}{}^{}X=Y`$, for all $`xX`$, $`{}_{}{}^{}x=x`$
* (non–triviality) for every infinite subset $`A`$ of $`X`$: $`A{}_{}{}^{}A`$
* (transfer principle) if $`\phi (x_1,\mathrm{},x_n)`$ is a $`\mathrm{\Sigma }_0`$ formula and $`a_1,\mathrm{},a_nV(X)`$, then $`V(X)\phi (a_1,\mathrm{},a_n)`$ iff $`V(Y)\phi (^{}a_1,\mathrm{},^{}a_n)`$
$`\mathrm{\Sigma }_0`$ formulas are first order formulas $`\phi `$ (in the Language $`\{\}`$) such that only quantifiers of the form $`xy`$ and $`xy`$ occur in $`\phi `$.
Note that the transfer principle is similar to the well–known Łoś’ theorem for ultraproduct–constructions. However, $`V(Y)`$ cannot be an ultrapower of $`V(X)`$, since otherwise Łoś’ theorem would apply to all first order formulas, including the formula “for all natural numbers $`n`$, there is a decreasing $``$–chain of length $`n`$”. This sentence cannot be true in $`(V(Y),)`$, since $``$ is well–founded and $`V(Y)`$ has nonstandard natural numbers. However, $`Y`$ can be an ultrapower of $`X`$, and this specific kind of nonstandard extension is the most important one in our context:
Let $`𝒰`$ be a countably incomplete ultrafilter over a set $`I`$ (i.e. there is a countable family $`A_n`$ of elements of $`𝒰`$ such that $`A_n𝒰`$). Let $`Y`$ be the ultrapower of $`X`$ with respect to $`𝒰`$. (Without loss of generality we can assume that $`X`$ and $`Y`$ are base–sets). Then $`V(Y)`$ can be made a nonstandard extension of $`V(X)`$ in a way that the map restricted to $`X`$ is the usual ultrapower injection, i.e. $`{}_{}{}^{}x`$ is the equivalence class of the constant function $`c_x`$. Non–triviality follows from the fact that $`𝒰`$ is countably incomplete, and the transfer principle is proved similar to Łoś’ theorem.
Nonstandard extension obtained in this way are called bounded ultrapowers.
In any nonstandard extension there are infinite numbers: The sets $`{}_{}{}^{}`$ and $`{}_{}{}^{}`$ are called the nonstandard natural and real numbers. Sometimes we will call $``$ and $``$ the standard natural (or real) numbers to emphasis the difference.
A nonstandard number $`r`$ is called finite if there is a standard natural number $`n`$ such that $`|r|<n`$ (in $`V(^{}X)`$). So a natural number is finite iff it is standard. A nonstandard number that is not finite is called infinite. A nonstandard real $`r`$ is called infinitesimal if $`r=0`$ or $`r0`$ and $`\frac{1}{r}`$ is infinite.
Let $`r,s{}_{}{}^{}`$. $`rs`$ stands for “$`rs`$ is infinitesimal”. For each finite $`r{}_{}{}^{}`$ there is a unique real number $`r^{}`$ such that $`r^{}r`$. This $`r^{}`$ is denoted by $`\mathrm{std}(r)`$
Similar notation will be introduced for elements of nonstandard normed vectorspaces (on page 3.1).
The following notions are central in nonstandard analysis:
$`yV(Y)`$ is called standard if there is a $`xV(X)`$ such that $`y={}_{}{}^{}x`$.
$`yV(Y)`$ is called internal if there is a $`xV(X)`$ such that $`y{}_{}{}^{}x`$.
$`yV(Y)`$ is called external if $`y`$ is not internal.
Note that for $`yY`$, $`y`$ is standard iff $`yX`$, so the notation is compatible with our use of standard natural (or real) numbers.
An example of an internal set that is not standard is the set of all nonstandard natural numbers less than some $`m`$, where $`m`$ is infinite.
If $`A,B`$ are internal (or standard), then so are $`AB`$, $`AB`$, etc.
The set of all infinite natural numbers in not internal. This is a special case of the so–called spillover principle:
###### Theorem 2.3.
(spillover principle) Let $`A`$ be internal.
* $`A`$ contains arbitrary large finite (i.e. standard) natural numbers iff $`A`$ contains arbitrary small infinite natural numbers.
* $`A`$ contains arbitrary small positive standard reals iff $`A`$ contains arbitrary large positive infinitesimals.
It is important to note that $`A`$ just has to be internal, it is not necessary that $`A`$ is internal.
If $`A`$ is a set in $`V(X)`$, then $`P^{\text{fin}}(A)`$, the set of all finite subsets of $`A`$, is in $`V(X)`$ as well, and is mapped to $`{}_{}{}^{}P_{}^{\text{fin}}(A)V(Y)`$. If $`A`$ was infinite, this set will contain more elements than just the sets $`{}_{}{}^{}B`$, where $`BA`$ finite. These new elements are called hyperfinite (note that they do not have to be finite). So a set $`B`$ in $`V(Y)`$ is called hyperfinite, if there is some $`AV(X)`$ such that $`B{}_{}{}^{}P_{}^{\text{fin}}(A)`$.
The example above, $`\{n{}_{}{}^{}:n<m\}`$, is an example of an infinite, hyperfinite, internal set. It is easy to see that all hyperfinite sets are internal.
For the examples 3.9 and 4.3, we use a nonstandard extension such that the standard reals are a subset of a hyperfinite set.
More general, we call a nonstandard extension $`V(Y)`$ an enlargement if for every $`AV(X)X`$ there is a hyperfinite $`BV(Y)`$ such that $`\{{}_{}{}^{}a:aA\}B`$.
It is provable (using the Axiom of Choice, of course) that for any base–set $`X`$ there is a index set $`I`$ and an ultrafilter $`𝒰`$ over $`I`$ such that the bounded ultrapower with respect to $`𝒰`$ is an enlargement. (For details, see e.g. \[CK90\]. The outline of the proof is as follows:
1. for every $`\kappa `$, there are $`\kappa `$–good filters
2. $`\kappa `$–good filters result in $`\kappa `$–saturated extensions
3. Given $`X`$, let $`\kappa _X`$ be the cardinality of $`V(X)`$. Then $`\kappa _X`$–saturated extensions are enlargements.)
The only consequence of the concept of enlargement that we are going to use is the following:
###### Lemma 2.4.
Let $`V(^{}X)`$ be a nonstandard enlargement of the universe. Then there is a (infinite) nonstandard natural number $`k`$ such that $`\mathrm{sin}(kt)`$ is infinitesimal for all standard reals $`t`$.
###### Proof.
The standard reals $``$ are a subset of some hyperfinite set $`A`$ of nonstandard reals.
It is a well–known (classical) fact that for every finite set of reals $`A`$, and every positive $`\epsilon `$ there is a natural number $`k`$ such that for all $`tA`$, $`\mathrm{sin}(kt)<\epsilon `$. So by transfer of the classical fact, setting $`\epsilon `$ infinitesimal, we get a $`k{}_{}{}^{}`$ such that for all $`t`$, $`\mathrm{sin}(kt)`$ is infinitesimal. ∎
When we compare nonstandard constructions with classical constructions, we will assume that our nonstandard extension is a bounded ultrapower (generated by the same ultrafilter) — otherwise it isn’t clear what the classical analogon should be. It should be noted that there are nonstandard universes which are not bounded ultrapowers: An ultrapower is always $`\mathrm{}_1`$–saturated (since all filters are $`\mathrm{}_1`$–good), but the union of a countable elementary chain of ultrapowers is not (see \[CK90, page 290, 4.4.29\]). However, every extension can be seen as limit of ultrapowers (see \[CK90\]), and in practice, authors often restrict their attention to the case of ultrapowers (see e.g. \[HL85, page 88\]).
## 3 Ultrapowers and Nonstandard–Hulls
### 3.1 Constructions without Semigroups
Let $`E`$ be a normed vectorspace, $`I`$ an arbitrary set. We define $`l^{\mathrm{}}(E)`$ to be the space of all bounded $`E`$-valued $`I`$-sequences with the sup–norm (denoted by $`_{\mathrm{}}`$). Then $`l^{\mathrm{}}(E)`$ is a normed vectorspace. Note that $`I`$ does not have to be countable, it can be of any cardinality.
Assume $`𝒰`$ is a filter over $`I`$ (not necessarily an ultrafilter). As in \[Der80\], we define $`c_𝒰`$ to be the set of sequences $`f_i`$ such that for all $`\epsilon `$ there is a set $`JI`$ such that $`J𝒰`$ and $`|f_i|<\epsilon `$ for all $`iJ`$.
###### Lemma 3.1.
$`|y+c_𝒰|:=inf\{|x|:xy+c_𝒰\}`$ is a norm on $`l^{\mathrm{}}\left(E\right)/c_𝒰`$, and $`E`$ is a Banachspace if and only if $`l^{\mathrm{}}(E)`$ is a Banachspace. If $`E`$ is a Banachspace, then $`l^{\mathrm{}}\left(E\right)/c_𝒰`$ is a Banachspace as well.
###### Proof.
If $`E`$ is a Banachspace, then so is $`l^{\mathrm{}}(E)`$ (since Cauchy–sequences converge pointwise). Clearly, $`c_𝒰`$ is a closed subspace of $`l^{\mathrm{}}(E)`$. And any normed (and complete) vectorspace can be factored by a closed subspace, resulting in a normed (and complete, resp.) vectorspace. If $`l^{\mathrm{}}(E)`$ is complete, then clearly so is $`E`$ (otherwise take any non–converging Cauchy–sequence $`f_i`$ in $`E`$, and map it to the sequence $`g_i=(f_i,0,0,\mathrm{})`$ in $`l^{\mathrm{}}(E)`$). ∎
Note that $`l^{\mathrm{}}\left(E\right)/c_𝒰`$ can be a Banachspace although $`E`$ (and therefore $`l^{\mathrm{}}(E)`$) is not, see the remark after theorem 3.3.
The corresponding nonstandard construction is the following: Let $`V({}_{}{}^{}X)`$ be a nonstandard universe, and $`EX`$ a normed vectorspace. Then $`{}_{}{}^{}E`$ is a nonstandard normed vectorspace, and a standard vectorspace (without canonical norm).
We define the finite part of $`{}_{}{}^{}E`$, denoted by $`\mathrm{fin}{}_{}{}^{}E`$, to consist of all elements $`f`$ of $`{}_{}{}^{}E`$ such that $`|f|`$ is finite. $`fg`$ means that $`|fg|`$ is infinitesimal. The infinitesimal part of $`{}_{}{}^{}E`$, denoted by $`E_0`$, consists of all $`f`$ such that $`|f|`$ is infinitesimal, i.e. $`f0`$. Clearly, $`E_0`$ is a (standard) sub–vectorspace of $`\mathrm{fin}{}_{}{}^{}E`$.
$`\widehat{E}`$ is the quotient of $`\mathrm{fin}{}_{}{}^{}E`$ and $`E_0`$, and the canonical quotient map $`\mathrm{fin}{}_{}{}^{}E\widehat{E}`$ is denoted by $`\widehat{}`$, i.e. the quotient–class of $`f`$ is $`\widehat{f}`$, and $`f`$ is called representant of $`\widehat{f}`$.
$`\widehat{E}`$ is a normed vectorspace with the norm $`|\widehat{f}|:=\mathrm{std}(|f|)`$, where $`f`$ is any representant of $`\widehat{f}`$ and $`\mathrm{std}(r)`$ is the standard part of a finite nonstandard real $`r`$. Of course, the vectorspace–operations are defined by $`\widehat{f}+\widehat{g}:=\widehat{f+g}`$ and $`\alpha \widehat{f}:=\widehat{\alpha f}`$.
###### Lemma 3.2.
If $`V({}_{}{}^{}X)`$ is a bounded ultrapower, then $`\widehat{E}`$ is a Banachspace.
A proof can be found e.g. in \[Lin88, page 59\], noting that a bounded ultrapower is $`\mathrm{}_1`$–good and therefore the extension will be $`\mathrm{}_1`$–saturated.
A straightforward calculation proves that for an ultrafilter $`𝒰`$, the F–product and the nonstandard hull are the same:
###### Theorem 3.3.
Let $`𝒰`$ be an ultrafilter, $`V(^{}X)`$ the corresponding bounded ultrapower. Then $`\iota :l^{\mathrm{}}\left(E\right)/c_𝒰\widehat{E}`$ defined by $`\iota (f+c_𝒰)=\widehat{f}`$ is an isomorphism.
So if $`𝒰`$ is a countably incomplete ultrafilter, then $`l^{\mathrm{}}\left(E\right)/c_𝒰`$ is a Banachspace (according to lemma 3.2), regardless of whether $`E`$ was a Banachspace or not.
### 3.2 Classical Constructions for Semigroups
#### The Maximal Continuous Subspace:
Assume, $`T(t)`$ is a semigroup on $`E`$ (not necessarily continuous), and assume $`M,\omega `$ are such that $`T(t)Me^{\omega t}`$ (for continuous semigroups, such $`M,\omega `$ always exist, according to lemma 2.1).
Define $`\left(E\right)^{T\text{–max}}:=\{fE:lim_{t0}|T(t)ff|=0\}`$.
A subspace $`F`$ of $`E`$ is called $`T`$–invariant, if $`T(t)(f)F`$ for all $`fF`$ and $`t^+`$.
###### Lemma 3.4.
$`\left(E\right)^{T\text{–max}}`$ is a closed $`T`$–invariant subspace of $`E`$, and it is maximal in the family of subspaces $`F`$ of $`E`$ such that $`T(t)`$ restricted to $`F`$ is continuous.
###### Proof.
* closed: Assume $`f_i,fE`$ such that $`f_nf`$, and choose an $`\epsilon >0`$. Let $`n`$ be such that $`|f_nf|<\mathrm{min}(\frac{\epsilon }{6M},\frac{\epsilon }{3})`$. $`|T(t)ff||T(t)fT(t)f_n|+|T(t)f_nf_n|+|f_nf|`$. If $`t<\delta `$ (for a suitable $`\delta `$) then $`T(t)<2M`$ and $`|T(t)f_nf_n|<\frac{\epsilon }{3}`$ (since $`f_n\left(E\right)^{T\text{–max}}`$), so $`|T(t)ff|2M\frac{\epsilon }{6M}+\frac{\epsilon }{3}+\frac{\epsilon }{3}\epsilon `$.
* invariant: Assume $`f\left(E\right)^{T\text{–max}},s_0^+`$. Then $`|T(t)(T(s)(f))T(s)(f)|=|T(s)(T(t)(f)f)|T(s)|T(t)ff|`$.
* maximal: If $`T`$ is continuous on $`F`$ and $`fF`$, then by definition $`f`$ is an element of $`\left(E\right)^{T\text{–max}}`$.
Remark: This definition and lemma just isolate a part of the proof that is given e.g. in \[Der80\] for corollary 3.5. This way we don’t have to repeat the same argument again (e.g. for lemma 3.8).
#### $`𝐦^𝐓`$:
Let $`T(t)`$ be continuous, $`T(t)Me^{\omega t}`$, and let $`f_i`$ be an element of $`l^{\mathrm{}}(E)`$. Define $`\stackrel{~}{T}(t)(f_i):=T(t)f_i`$.
Since $`T(t)`$ is bounded, $`\stackrel{~}{T}(t)(f_i)l^{\mathrm{}}(E)`$ for all $`f_il^{\mathrm{}}(E)`$ and $`t0`$. It is clear that $`\stackrel{~}{T}`$ defines a semigroup on $`l^{\mathrm{}}(E)`$. Also, $`\stackrel{~}{T}(t)Me^{\omega t}`$.
But in general $`\stackrel{~}{T}`$ will not be continuous on $`l^{\mathrm{}}(E)`$: Let $`T`$ be the translation semigroup on $`C^\text{b}`$ as in 2.2, $`I=\omega `$, $`𝒰`$ a free ultrafilter over $`\omega `$, and $`f_k(x)=sin(kx)C^\text{b}`$. Then $`f=f_kl^{\mathrm{}}(E)`$, but $`\stackrel{~}{T}(t)ff`$ does not converge.
On can show rather easily that $`\stackrel{~}{T}`$ is continuous on $`l^{\mathrm{}}(E)`$ if and only if $`T`$ is uniformly continuous on $`E`$ (assuming of course that the index set $`I`$ is infinite).
We define $`m^T`$ to be the maximal continuous subspace of $`l^{\mathrm{}}(E)`$ with respect to $`T`$.
###### Corollary 3.5.
$`m^T`$ is a closed, $`\stackrel{~}{T}`$–invariant subspace of $`l^{\mathrm{}}(E)`$.
#### F-Product of a semigroup:
As we have seen, if $`T`$ is a semigroup defined on $`E`$, under certain assumptions $`T`$ can be extended to a (not necessarily continuous) semigroup $`\stackrel{~}{T}`$ on $`l^{\mathrm{}}(E)`$. Assume $`f_ic_𝒰`$ (i.e. for all $`\epsilon ^+`$, $`|f_i|<\epsilon `$ holds on a filter–set). $`\stackrel{~}{T}(t)(f_i)c_𝒰`$ for all $`t0`$, since $`\stackrel{~}{T}(t)`$ is a bounded operator. This means that $`c_𝒰`$ is $`\stackrel{~}{T}`$–invariant, and that we can define $`\stackrel{~}{T}`$ on $`l^{\mathrm{}}\left(E\right)/c_𝒰`$. Of course, $`\stackrel{~}{T}`$ is a semigroup on $`l^{\mathrm{}}\left(E\right)/c_𝒰`$ (and generally not continuous).
$`c_𝒰`$ is not a subspace of $`m^T`$. (For example, any $`f_i`$ such that $`f_i=0`$ on a filter–set is in $`c_𝒰`$, but not necessarily in $`m^T`$.) However, it is easy to see that $`c_𝒰m^T`$ is a closed subset of $`m^T`$, and $`m^T/c_𝒰m^T`$ is a subspace of $`l^{\mathrm{}}\left(E\right)/c_𝒰`$.
To be more exact:
###### Lemma 3.6.
$`\varphi :m^T/c_𝒰m^Tl^{\mathrm{}}\left(E\right)/c_𝒰`$ defined by $`\varphi (f_i+c_𝒰m^T)=f_i+c_𝒰`$ is a (well–defined) injective isometry.
###### Proof.
See \[Kru90, page 163\]. ∎
###### Lemma 3.7.
$`m^T/c_𝒰m^T`$ is a closed, $`\stackrel{~}{T}`$–invariant subspace of $`l^{\mathrm{}}\left(E\right)/c_𝒰`$.
$`\stackrel{~}{T}`$ restricted to $`m^T/c_𝒰m^T`$ is continuous.
###### Proof.
The quotient $`m^T/c_𝒰m^T`$ of the Banachspace $`m^T`$ is a Banachspace, and so is its isometric image. The rest is clear. ∎
In general, $`m^T/c_𝒰m^T`$ is not the maximal continuous subspace of $`l^{\mathrm{}}\left(E\right)/c_𝒰`$. This is shown in example 3.9 (using that $`m^T/c_𝒰m^T`$ corresponds to $`\widehat{E_T}`$, see lemma 3.11).
### 3.3 Nonstandard Constructions for Semigroups
Let $`V(^{}X)`$ be a nonstandard universe, $`EX`$ a Banachspace.
Assume $`A:EE`$ is a continuous linear operator, and $`f`$ an element of $`\mathrm{fin}{}_{}{}^{}E`$, $`g`$ an element of $`E_0`$ (the finite and infinitesimal parts of the nonstandard vector space $`{}_{}{}^{}E`$, resp.). Then $`{}_{}{}^{}Af`$ (and $`{}_{}{}^{}Ag`$) are again finite (and infinitesimal, resp.). Therefore $`A`$ defines a bounded linear operator $`\widehat{A}:\widehat{E}\widehat{E}`$ by $`\widehat{A}(\widehat{f})=\widehat{Af}`$. (If $`A`$ is not bounded, then generally $`\widehat{A}`$ cannot be defined in a canonical way.)
If $`T`$ is a semigroup, then for all positive reals $`t`$, $`T(t)`$ is continuous, and $`\widehat{T}(t):=\widehat{T(t)}`$ is a semigroup on $`\widehat{E}`$.
If $`T(t)<Me^{\omega t}`$ (this is always the case if $`T`$ is continuous), then an alternative way to define $`\widehat{T}`$ is the following: $`{}_{}{}^{}T`$ is a nonstandard semigroup on $`{}_{}{}^{}E`$. If $`t`$ is a positive, finite nonstandard real, then $`{}_{}{}^{}T(t)<Me^{\omega t}`$, so we get: If $`f{}_{}{}^{}E`$ is finite (or infinitesimal), then so is $`{}_{}{}^{}T(t)f`$. Therefore $`\widehat{T}(t):\widehat{E}\widehat{E}`$ is well–defined by $`\widehat{T}(t)(\widehat{f})=\widehat{{}_{}{}^{}T(t)(f)}`$, even for finite nonstandard $`t`$. If $`t`$ is standard, then the two definitions are equivalent.
Assume $`T`$ is continuous, i.e. for all $`f`$ in $`E`$ and all $`\epsilon >0`$ there is a $`\delta >0`$ such that $`|T(t)ff|<\epsilon `$ for all $`t<\delta `$. We apply the transfer principle to this sentence, and get: For all $`f{}_{}{}^{}E`$ and all $`\epsilon {}_{}{}^{}_{}^{+}`$, there is a $`\delta {}_{}{}^{}_{}^{+}`$ such that $`|{}_{}{}^{}T(t)ff|<\epsilon `$ for all positive reals $`t<\delta `$: However, if $`\epsilon >0`$ is a standard real, then it is not guaranteed that $`\delta >0`$ can be chosen to be standard as well ($`\delta `$ could be infinitesimal). We define two subspaces of $`\mathrm{fin}{}_{}{}^{}E`$ (the finite part of $`{}_{}{}^{}E`$): $`E_T`$ consists of all the vectors $`f`$ with the following property: If $`\epsilon >0`$ is standard, then there is a standard $`\delta >0`$ such that $`|{}_{}{}^{}T(t)ff|<\epsilon `$ for all nonstandard $`t<\delta `$. $`E_{\text{max}}`$ consists of all $`f`$ satisfying the same condition for standard $`t<\delta `$ only. I.e. a finite $`f{}_{}{}^{}E`$ is in $`E_{\text{max}}`$ iff for all $`ϵ^+`$ there is a $`\delta ^+`$ such that $`|{}_{}{}^{}T(t)ff|<\epsilon `$ for all positive standard reals $`t<\delta `$.
It is easy to see that a finite $`f{}_{}{}^{}E`$ is in $`E_T`$ iff for all positive infinitesimal t, $`|{}_{}{}^{}T(t)ff|`$ is infinitesimal.
Clearly, $`E_0`$ is a subspace of $`E_T`$, which is in turn a subspace of $`E_{\text{max}}`$. So, if $`fE_T`$ (or $`E_{\text{max}}`$) and $`fg`$, then $`gE_T`$ (or $`E_{\text{max}}`$, resp.), and we can define the quotient spaces $`\widehat{E_T}=E_T/E_0`$ and $`\widehat{E}_{\text{max}}=E_{\text{max}}/E_0`$. Remember that (other then $`\mathrm{fin}{}_{}{}^{}E`$), $`E_T`$ is a (canonically) normed vectorspace.
###### Lemma 3.8.
1. $`\widehat{E}_{\text{max}}`$ is the maximal continuous subspace of $`\widehat{E}`$ with respect to $`\widehat{T}`$.
2. $`\widehat{E}_{\text{max}}`$ is a closed, $`\widehat{T}`$–invariant subspace of $`\widehat{E}`$
3. $`\widehat{E_T}`$ is a closed, $`\widehat{T}`$–invariant subspace of $`\widehat{E}_{\text{max}}`$.
###### Proof.
1. Assume, $`lim_{t0}|\widehat{T}(t)\widehat{f}\widehat{f}|=0`$, and fix $`\epsilon ^+`$. Then there is a $`\delta ^+`$ such that for all positive reals $`t<\delta `$, $`|\widehat{T}(t)\widehat{f}\widehat{f}|<\frac{\epsilon }{2}`$. Assume that $`f`$ is a representant of the quotient $`\widehat{f}`$, and that $`t<\delta `$. Then $`\mathrm{std}(|{}_{}{}^{}T(t)ff|)<\frac{\epsilon }{2}`$, and therefore $`|{}_{}{}^{}T(t)ff|<\epsilon `$ in $`{}_{}{}^{}`$, so $`fE_{\text{max}}`$ and $`\widehat{f}\widehat{E}_{\text{max}}`$.
2. this follows from 1. and lemma 3.4.
3. Assume that $`\widehat{f}_n\widehat{E_T}`$, and $`\widehat{f}_n\widehat{f}`$ in $`E_T`$. Let $`f_n`$ and $`f`$ be representants of $`\widehat{f}_n`$ and $`\widehat{f}`$, resp. Assume $`h`$ is a positive infinitesimal. Then $`|T(h)ff||T(h)fT(h)f_n|+|T(h)f_nf_n|+|f_nf|`$, which is smaller than every positive standard $`ϵ`$, since $`|ff_n|`$ gets arbitrary small, $`|T(h)ff|T(h)|ff_n|`$, and $`|T(h)f_nf_n|`$ is infinitesimal.
To show the invariance, assume that $`\widehat{f}\widehat{E_T}`$, and let $`fE_T`$ be a representant. Let $`t`$ be a positive standard real, and $`h`$ is a positive infinitesimal. $`|T(h)T(t)fT(t)f|T(t)|T(h)ff|`$, which is infinitesimal. So $`T(t)fE_T`$, and therefore $`\widehat{T}(t)\widehat{f}\widehat{E_T}`$.
Remark: In \[Wol00\], an alternative way to construct $`E_T`$ is presented. Theorem 4.6.1 there corresponds to our 4.4.
In general, $`\widehat{E_T}`$ is a proper subset of $`\widehat{E}_{\text{max}}`$ (or equivalently: $`E_T`$ is a proper subset of $`E_{\text{max}}`$). To see this, we bring an example similar to \[Kru90, page 165\] or \[Wol84, page 209\]:
###### Example 3.9.
Let $`T`$ be the translation semigroup on $`C^\text{b}`$ (as in example 2.2), $`V(^{}X)`$ a enlargement, $`k`$ such that $`\mathrm{sin}(kt)`$ is infinitesimal for all $`t`$ (see lemma 2.4). Define the element $`f`$ of $`{}_{}{}^{}C_{}^{\text{b}}`$ by $`f(x)=\mathrm{sin}(kx)`$. Then $`|f|=1`$, and $`f`$ is an element of $`E_{\text{max}}`$, but not of $`E_T`$.
###### Proof.
To see that $`fE_{\text{max}}`$, it is enough to show that for all positive (standard) reals $`t`$, $`|T(t)ff|`$ is infinitesimal. For any standard reals $`t`$, $`\mathrm{sin}(k(x+t))\mathrm{sin}(kx)=\mathrm{sin}(kx)(\mathrm{cos}(kt)1)+\mathrm{sin}(kt)\mathrm{cos}(kx)`$. $`\mathrm{sin}(kx)`$ and $`\mathrm{cos}(kx)`$ are finite, and $`\mathrm{sin}(kt)`$ and $`(\mathrm{cos}(kt)1)=(\sqrt{1\mathrm{sin}^2(kt)}1)`$ are infinitesimal. Therefore $`|T(t)ff|=\mathrm{sin}(k(x+t))\mathrm{sin}(kx)_{\mathrm{}}`$ is infinitesimal.
To see that $`fE_T`$, it is enough to note that $`h:=\frac{1}{k}`$ is infinitesimal, but $`|T(h)ff|=sin(k(x+h))sin(kx)_{\mathrm{}}=sin(x+1)sin(x)_{\mathrm{}}`$ is not. ∎
### 3.4 Relation of Classical and Nonstandard Constructions
Assume $`V(^{}X)`$ is the bounded ultrapower of an ultrafilter $`𝒰`$ over $`I`$. Then the isomorphism $`\iota `$ of theorem 3.3 allows us to identify $`l^{\mathrm{}}\left(E\right)/c_𝒰`$ and $`\widehat{E}`$. We can also consider $`m^T/c_𝒰m^T`$ to be a subspace of $`l^{\mathrm{}}\left(E\right)/c_𝒰`$ (via the injective isometry $`\varphi `$ of lemma 3.6).
###### Lemma 3.10.
$`\widehat{E}_{\text{max}}`$ is the maximal continuous subspace of $`l^{\mathrm{}}\left(E\right)/c_𝒰`$ with respect to $`\stackrel{~}{T}`$
(More formally, one would have to write $`\widehat{E}_{\text{max}}=\iota (\left(l^{\mathrm{}}\left(E\right)/c_𝒰\right)^{\stackrel{~}{T}\text{–max}})`$.)
###### Proof.
This follows from lemma 3.8 , and the fact that $`\iota (\stackrel{~}{T}(t)f)=\widehat{T}(t)(\iota (f))`$ for all $`fl^{\mathrm{}}\left(E\right)/c_𝒰`$. ∎
###### Lemma 3.11.
$`m^T/c_𝒰m^T\widehat{E_T}`$
(Again, more formally this should be written as $`\iota (\varphi (m^T/c_𝒰m^T))\widehat{E_T}`$.)
###### Proof.
By definition, $`f_im^T`$ iff for all reals $`ϵ>0`$ there is a real $`\delta >0`$ such that $`T(t)f_if_i_{\mathrm{}}<ϵ`$ for all reals $`t<\delta `$. Let $`\widehat{f}\widehat{E}`$ correspond to $`f_i`$. Assume that $`h`$ is an infinitesimal nonstandard real. We want to show that for an arbitrary fixed positive real $`ϵ`$, $`|{}_{}{}^{}T(h)\widehat{f}\widehat{f}|<ϵ`$. Let $`\delta `$ be the real corresponding to $`ϵ`$ as above. $`h`$ corresponds to a sequence of reals $`h_i`$ ($`iI`$) such that $`U=\{iI:h_i<\delta \}`$ is in the ultrafilter. $`{}_{}{}^{}T(h)\widehat{f}\widehat{f}`$ corresponds to the sequence $`T(h_i)f_if_i`$. If $`i`$ is in $`U`$, then $`|T(h_i)f_if_i|<ϵ`$, and $`U`$ is a filter–set, therefore $`|{}_{}{}^{}T(h)\widehat{f}\widehat{f}|<ϵ`$. ∎
It is not immediately clear whether $`m^T/c_𝒰m^T`$ is always identical to $`\widehat{E_T}`$.
One could construct a counterexample under the following assumption:
###### Assumption 3.12.
Fix $`\omega ,M,\eta ,\epsilon _1,\epsilon _2,\mathrm{},\delta _1,\delta _2,\mathrm{}`$. Assume that for all $`n,m^+`$ there is a Banachspace $`E_m^n`$, a continuous semigroup $`T_m^n`$ with $`T_m^n(t)<Me^{\omega t}`$ and a $`f_m^nE_m^n`$ such that
* $`|f_m^n|=1`$
* $`|T_m^n(t)f_m^nf_m^n|<\frac{1}{n}`$ for all $`1in`$, $`t<\delta _i`$
* For all $`\overline{f}E_m^n`$ such that $`|\overline{f}f_m^n|<\eta `$ exists a $`t<\frac{1}{m}`$ such that $`|T_m^n(t)\overline{f}\overline{f}|>\epsilon _n`$
If this assumption holds, then let the index set $`I`$ be $`^+\times ^+`$, and define $`\overline{E}=l^{\mathrm{}}(E_j^i)`$ with the sup–norm. (I.e. $`f_{i,j}\overline{E}`$ iff $`f_{i,j}E_j^i`$ and $`|f_{i,j}|`$ is bounded.) $`T(t)g_{i,j}:=T_j^i(t)g_{i,j}`$ is a semigroup on $`\overline{E}`$. Let $`E`$ be the maximal continuous subspace of $`\overline{E}`$, and define $`h_m^n=h_m^n(i,j)E`$ by
$`h_m^n(i,j)=\{\begin{array}{cc}f_m^n\hfill & \text{if }i=n,j=m\text{,}\hfill \\ 0\hfill & \text{otherwise}.\hfill \end{array}`$.
Assume that $`𝒰`$ is an ultrafilter over $`I=^+\times ^+`$, and that there is a partition $`\{\sigma _{i,j}\}`$ $`(i,j)`$ of $`I`$ such that for all $`n`$ and all functions $`f:`$, $`_{i>n,j>f(i)}\sigma _{i,j}𝒰`$ (it is easy to see that such a countably incomplete filter exists, e.g. apply a suitable filter–basis to \[Lin88, theorem A.4\]).
Then define $`k=k_{i,j}\overline{E}`$ by $`k_{i,j}=h_m^n`$, where $`n,m`$ is the (unique) pair of natural numbers such that $`(i,j)`$ is an element of $`\sigma _{n,m}`$.
Now consider the nonstandard extension defined by $`𝒰`$, and interpret $`k`$ as nonstandard element of $`{}_{}{}^{}E`$.
Clearly $`|k|=1`$, since $`|k_{i,j}|=1`$ for all $`(i,j)`$. For all natural numbers $`n`$, $`0<t<\epsilon _n`$ and $`(i,j)_{r>n,s\omega }\sigma _{r,s}`$: $`|T(t)k_{i,j}k_{i,j}|<\frac{1}{n}`$. So if $`r_{i,j}`$ is infinitesimal, then so is $`|{}_{}{}^{}T(r)kk|`$, i.e. $`kE_T`$.
Assume that the equivalence class of $`k`$ is in $`m^T/c_𝒰m^T`$. Pick a representant $`\overline{k}_{i,j}m^T`$, and $`U𝒰`$ such that $`|\overline{k}_{i,j}k_{i,j}|<\eta `$ for all $`(i,j)U`$. Since $`U𝒰`$, there is an $`n_0`$ such that for all $`m`$ there is a $`m^{}`$ such that $`\sigma _{n_0,m^{}}U\mathrm{}`$. Since $`\overline{k}_im^T`$, there is a $`m_0`$ such that $`|T(t)\overline{k}_i\overline{k}_i|<\delta _{n_0}`$ for all $`(i,j)I`$ and $`0<t<\frac{1}{m_0}`$. If $`m_{0}^{}{}_{}{}^{}>m_0`$ and $`(i_0,j_0)\sigma _{n_0,m_{0}^{}{}_{}{}^{}}U`$, then there is a $`t^+`$ such that $`t<\frac{1}{m_{0}^{}{}_{}{}^{}}<\frac{1}{m_0}`$ and $`|T(t)\overline{k}_i\overline{k}_i|\delta _{n_0}`$, a contradiction.
However, our investigation of the nonstandard generator of a semigroup will show that $`m^T/c_𝒰m^T=\widehat{E_T}`$.
## 4 The Generator
### 4.1 The Generator in $`{}_{}{}^{}E`$
If $`T`$ is a continuous semigroup, then $`{}_{}{}^{}T(t)`$ maps the finite (or infinitesimal) elements of $`{}_{}{}^{}E`$ to finite (or infinitesimal, resp.) elements for all finite $`t0`$. That is not true for the (unbounded) linear map $`A`$, as the following trivial example shows:
###### Example 4.1.
Let $`T`$ be the translation semigroup on $`C^\text{b}`$ (as in example 2.2), $`V(^{}X)`$ a nonstandard extension, $`k`$ an infinite natural number. Then there is an infinitesimal $`f_1{}_{}{}^{}E`$ such that $`f_1{}_{}{}^{}D(A)`$ and $`{}_{}{}^{}A(f)`$ is infinite.
###### Proof.
Define $`f_1`$ by case distinction:
$`f_1=\{\begin{array}{cc}0\hfill & x<0\hfill \\ k^3x^2\hfill & 0x<\frac{1}{k^2}\hfill \\ k^3x^2+4kx\frac{2}{k}\hfill & \frac{1}{k^2}x<\frac{2}{k^2}\hfill \\ \frac{2}{k}\hfill & \frac{2}{k^2}x\hfill \end{array}`$.
Then $`{}_{}{}^{}Af_1=\{\begin{array}{cc}0\hfill & x<0\hfill \\ 2k^3x\hfill & 0x<\frac{1}{k^2}\hfill \\ 2k^3x+4k\hfill & \frac{1}{k^2}x<\frac{2}{k^2}\hfill \\ 0\hfill & \frac{2}{k^2}x\hfill \end{array}`$. ∎
Let $`V(^{}X)`$ be a nonstandard universe. Then the subspace $`D(A)`$ of $`E`$ is mapped to a subspace $`{}_{}{}^{}D(A)`$ of $`{}_{}{}^{}E`$.
###### Lemma 4.2.
Assume, $`f`$ and $`g`$ are elements of $`{}_{}{}^{}D(A)`$
1. If $`f`$ and $`{}_{}{}^{}A(f)`$ are both finite, then $`fE_T`$.
2. If $`f`$ is finite and $`{}_{}{}^{}A(f)E_T`$, then $`\frac{{}_{}{}^{}T(h)ff}{h}{}_{}{}^{}A(f)`$ for all positive infinitesimal reals $`h`$.
3. If $`k{}_{}{}^{}E`$ and $`\frac{{}_{}{}^{}T(h)ff}{h}k`$ for all positive infinitesimal reals $`h`$, then $`k{}_{}{}^{}A(f)`$
4. If $`f`$, $`g`$, $`{}_{}{}^{}A(f)`$ and $`{}_{}{}^{}A(g)`$ all are in $`E_T`$, and $`fg`$, then $`{}_{}{}^{}A(f){}_{}{}^{}A(g)`$
###### Proof.
1. By transfer of 2.1.4 we get: $`|{}_{}{}^{}T(t)ff|tsup_{st}{}_{}{}^{}T(s)|{}_{}{}^{}A(f)|tM(e^{\omega t}+1)|{}_{}{}^{}A(f)|`$, which is infinitesimal for infinitesimal $`t`$.
2. For every standard $`\epsilon >0`$ and infinitesimal $`s`$, we have $`|{}_{}{}^{}T(s){}_{}{}^{}A(f){}_{}{}^{}A(f)|<\epsilon `$. Therefore transfer of 2.1.5 implies that $`|\frac{{}_{}{}^{}T(h)ff}{h}{}_{}{}^{}A(f)|\epsilon `$ for infinitesimal $`h`$. Since $`\epsilon `$ was arbitrary, $`|\frac{{}_{}{}^{}T(h)ff}{h}{}_{}{}^{}A(f)|`$ is infinitesimal.
3. By transfer of the definitions of $`A`$ and $`D(A)`$, we get the following: For all positive nonstandard reals $`\epsilon `$ there is a positive nonstandard real $`\delta `$ such that $`|\frac{{}_{}{}^{}T(h)ff}{h}{}_{}{}^{}Af|<ϵ`$ for all $`h\delta `$. Let $`\epsilon `$ be a fixed positive infinitesimal, and choose an infinitesimal $`h\delta `$ for the appropriate $`\delta `$. Then $`\frac{{}_{}{}^{}T(h)ff}{h}{}_{}{}^{}Af`$, and by assumption $`\frac{{}_{}{}^{}T(h)ff}{h}k`$, therefore $`k{}_{}{}^{}A`$.
4. Since $`{}_{}{}^{}AfE_T`$, $`\frac{{}_{}{}^{}T(h)ff}{h}{}_{}{}^{}Af`$ for all infinitesimal $`h`$, and the same applies to $`g`$ and $`{}_{}{}^{}Ag`$. So we have $`|{}_{}{}^{}Af{}_{}{}^{}Ag||{}_{}{}^{}Af\frac{{}_{}{}^{}T(t)ff}{t}|+|\frac{{}_{}{}^{}T(t)ff}{t}\frac{{}_{}{}^{}T(t)gg}{t}|+|\frac{{}_{}{}^{}T(t)gg}{t}{}_{}{}^{}Ag|\epsilon `$. So we just have to show that $`|\frac{{}_{}{}^{}T(t)ff}{t}\frac{{}_{}{}^{}T(t)gg}{t}|=|\frac{{}_{}{}^{}T(t)(fg)(fg)}{h}|`$ is infinitesimal for infinitesimal $`h`$. Transfer of 2.1.4 shows that $`|\frac{{}_{}{}^{}T(t)(k)(k)}{h}|sup_{sh}{}_{}{}^{}T(s)|{}_{}{}^{}A(k)|`$ for all positive $`h`$ and $`k{}_{}{}^{}D(A)`$. Now apply this to $`k=fg`$.
As we have already seen in example 4.1, $`fg`$ does not imply $`{}_{}{}^{}A(f){}_{}{}^{}A(g)`$ (set $`f:=f_1`$, $`g:=0`$). Also, it is not enough to assume that $`f`$, $`g`$, $`{}_{}{}^{}A(f)`$ and $`{}_{}{}^{}A(g)`$ are all in $`E_{\text{max}}{}_{}{}^{}D(A)`$ for the implication to hold:
###### Example 4.3.
Assume $`T`$, $`k`$ and $`f`$ are as in example 3.9. Define $`g(x):=\frac{\mathrm{cos}kx}{k}`$. Then $`g0`$. Clearly $`gE_{\text{max}}`$ (it is even in $`E_T`$), and $`{}_{}{}^{}A(g)=f`$, i.e. $`{}_{}{}^{}A(g)E_{\text{max}}`$, but $`{}_{}{}^{}A(g)`$ is not infinitesimal.
The same example shows that $`gE_T`$ does not imply $`{}_{}{}^{}A(g)E_T`$.
### 4.2 The Generator in $`\widehat{E_T}`$
As we have already seen in lemma 3.8, $`\widehat{T}`$ is a continuous semigroup on $`\widehat{E_T}`$ (since $`\widehat{E_T}`$ is a subspace of $`\widehat{E}_{\text{max}}`$). The generator is called $`\widehat{A}`$.
###### Theorem 4.4.
$`\widehat{f}`$ is element of $`D(\widehat{A})`$, iff there is a representant $`f`$ of $`\widehat{f}`$ such that $`f`$ is in $`D{}_{}{}^{}A`$ and $`{}_{}{}^{}Af`$ is in $`E_T`$. In this case, $`\widehat{A}(\widehat{f})=\widehat{{}_{}{}^{}A(f)}`$.
###### Proof.
Assume $`\widehat{f}`$ has a representant $`f`$ as in the lemma. We have to show that $`\widehat{f}D(\widehat{A})`$ and $`\widehat{A}(\widehat{f})=\widehat{{}_{}{}^{}A(f)}`$, i.e. we fix a (standard) real $`\epsilon >0`$ and have to find a $`\delta >0`$ such that $`|\frac{\widehat{T}(h)\widehat{f}\widehat{f}}{h}\widehat{{}_{}{}^{}A(f)}|<\epsilon `$ for all $`h<\delta `$. By lemma 4.2.2, $`n(h):=|\frac{{}_{}{}^{}T(h)ff}{h}{}_{}{}^{}A(f)|`$ is infinitesimal for infinitesimal $`h`$. Note that $`n`$ is an internal function, since it is element of $`{}_{}{}^{}\{h|\frac{1}{h}(T(h)ff)A(f)|:fD(A)\}`$. Therefore the set $`X=\{h:n(h)\epsilon \}`$ is internal as well, and we can apply the spillover principle 2.3: Assume toward a contradiction that for all standard $`\delta >0`$ there was a standard $`h<\delta `$ such that $`n(h)\epsilon `$, i.e. there are arbitrary small standard reals in $`X`$. Then there is a infinitesimal real in $`X`$ as well, a contradiction. Therefore $`\widehat{f}D(\widehat{A})`$.
Assume on the other hand that $`\widehat{f}`$, $`\widehat{g}`$ are elements of $`\widehat{E_T}`$ such that $`\widehat{A}(\widehat{f})=\widehat{g}`$. We want to show that $`fD{}_{}{}^{}A`$ and $`{}_{}{}^{}AfE_T`$ for some representant $`f`$ of $`\widehat{f}`$. Let $`\overline{f},\overline{g}`$ be arbitrary representants of $`f`$ and $`g`$, resp. (According to example 4.1, we cannot hope that $`\overline{f}`$ already has the required properties.) For all nonstandard natural numbers $`m`$ define $`c_m`$ to be the set if nonstandard natural numbers $`n>m`$ such that $`|\frac{{}_{}{}^{}T(\frac{1}{n})f^{}f^{}}{\frac{1}{n}}g^{}|<\frac{1}{m}`$. For all standard natural numbers $`m`$, the set $`c_m`$ is not empty, since $`\widehat{A}(\widehat{f})=\widehat{g}`$. The set $`X=\{m:c_M\mathrm{}\}`$ is internal, and contains all standard natural number, therefore it contains an infinite natural number as well (according to the spillover principle). Let $`m`$ be such an infinite number, and $`n`$ an element of $`c_m`$. Now we apply the transfer of 2.1.6, setting $`h:=\frac{1}{n}`$. So we get a $`f`$ in $`D({}_{}{}^{}A)`$ such that $`|\overline{f}f|sup_{sh}(|{}_{}{}^{}T(s)\overline{f}\overline{f}|)`$, which is infinitesimal, since $`\overline{f}E_T`$ and $`s`$ is infinitesimal. Therefore $`f`$ is a representant of $`\widehat{f}`$ as well. Also, $`{}_{}{}^{}A(f)=\frac{{}_{}{}^{}T(h)\overline{f}\overline{f}}{h}`$, so $`|{}_{}{}^{}A(f){}_{}{}^{}A(\overline{g})|\frac{1}{m}`$ is infinitesimal, therefore $`{}_{}{}^{}A(f)E_T`$ as well.
Let $`V(^{}X)`$ be a bounded ultrapower, $`T`$ continuous with generator $`A`$.
### 4.3 The equivalence of F–product and nonstandard hull
###### Lemma 4.5.
$`\widehat{E_T}`$=$`m^T/c_𝒰m^T`$
###### Proof.
Assume $`m^T/c_𝒰m^T\widehat{E_T}`$. Since $`m^T/c_𝒰m^T`$ is a closed subspace of $`\widehat{E_T}`$ and $`D(\widehat{A})`$ dense in $`\widehat{E_T}`$, there must be a $`\widehat{f}`$ which is element of $`D(\widehat{A})`$ but not $`m^T/c_𝒰m^T`$. We chose a representant $`f`$ of $`\widehat{f}`$ such that $`fD({}_{}{}^{}A)`$, $`{}_{}{}^{}AfE_T`$. Since $`f`$ is an element of the ultraproduct, $`f`$ is an equivalence class of a sequence $`f_i`$ ($`iI`$) such that for all $`i`$ $`f_iD(A)`$, $`A(f_i)l^{\mathrm{}}(E)`$. But then $`|\stackrel{~}{T}(h)ff|=sup_{iI}(|T(h)f_if_i|)hsup_{iI}(|A(f_i)|)sup_{sh}(T(s))`$ (according to 2.1.4), so $`fm^T`$, a contradiction. ∎
Combining this with 3.11 and 3.8, using the isometry $`\iota `$ of lemma 3.3, we get:
###### Theorem 4.6.
$`\widehat{E}`$ is isomorphic to $`l^{\mathrm{}}\left(E\right)/c_𝒰`$, and $`\widehat{E_T}`$ is isomorphic to $`m^T/c_𝒰m^T`$.
So we get the following picture ($`p`$ are the canonical projections from a space to its quotient, and $`\widehat{}`$ maps a finite element $`f`$ of $`{}_{}{}^{}E`$ to the equivalence class $`\widehat{f}`$. The maps labeled with $`p`$ and $`\widehat{}`$ are surjective, the ones labeled with $``$ isometries. $`>`$ denotes the subspace–relation, for $`>^1`$ the isometry $`\varphi `$ of lemma 3.6 is used):
(Remark: The entries of the first three rows are Banachspaces, the last row consists of vectorspaces.)
So we know that assumption 3.12 cannot hold. This proves
###### Corollary 4.7.
Assume that $`T_m^n`$ are contraction semigroups on $`E`$ (i.e. for all $`t^+`$, $`T(t)1`$), and and $`f_m^nE`$ are such that $`|f_m^n|=1`$ and $`|T_m^n(t)f_m^nf_m^n|<\frac{1}{n}`$ for all $`t<\frac{1}{n}`$. Then there is a $`f^{}E`$, and $`n,m`$, s.t. $`|f^{}f_m^n|<\frac{1}{2}`$ and $`|T_m^n(t)f^{}f^{}|<\frac{1}{n}`$ for all $`t<\frac{1}{m}`$.
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# Particle masses and the fifth dimension, final version
First we argue in an informal, qualitative way that it is natural to enlarge space-time to five dimensions to be able to solve the problem of elementary particle masses. Several criteria are developed for the success of this program. Extending the Poincaré group to the group C of all angle-preserving transformations of space-time is one such scheme which satisfies these criteria. Then we show that the field equation for spin 1/2 fermions coupled to a self-force gauge field predicts mass spectra of the desired type: for a certain range of a key parameter (Casimir invariant) a three-point mass spectrum which fits the “down” quarks $`d,s,`$ and $`b`$ to within their experimental bounds is obtained. Reasonable values of the coupling constant (of QCD magnitude) and the range of the spatial wave function (a few fermis) also result. Compatibility with the electroweak theory is also discussed.
1. INTRODUCTION
A theory of elementary particle masses which predicts the masses that we see in nature is lacking in present day particle physics. The Standard Model appeals to the Higgs mechanism. But even granting that the Higgs particle exists, successful fits must wait on the measurement of various unknown parameters . String theories claim to be able to predict these masses in principle, but they are still far from delivering quantitative numbers at their present stage .
First, some informal, qualitative remarks may be helpful to motivate the main idea of this paper. The idea that predicting particle masses should involve enlarging 4-D (“four-dimensional”) space-time (coordinates $`x^^\mu =\{x,y,z,x^oct\}`$ by a single new dimension, call it $`\lambda `$, seems very natural. The equal status of momentum $`𝐩`$, energy $`E`$, and mass $`m`$ in the free particle relation
$$𝐩^2E^2+m^2=0$$
(1)
suggests that in 5-D position space $`\lambda `$ should be conjugate to $`m`$, just as $`𝐫`$ is conjugate to $`𝐩`$ and $`t`$ is conjugate to $`E`$. (We shall use units $`c=1`$ in this paper.) And further, that the field equation for the field $`\varphi (x^^\mu ,\lambda )`$ of a free scalar boson, say, should be something like
$$(^2^2/t^2+^2/\lambda ^2)\varphi (x^^\mu ,\lambda )=0\text{ ,}$$
(2)
with the solution
$$\varphi (x^^\mu ,\lambda )=const\times \mathrm{exp}[i(𝐩𝐫Et\pm m\lambda )]$$
(3)
with the constraint (1) on the constants $`𝐩`$, $`E`$, and $`m`$.
However, this first try is too naive for several reasons. First, the new dimension $`\lambda `$ is simply grafted onto space-time, uncritically assuming that the enlarged space is still flat (cf. Eq. (2)). The symmetry group of Eq. (2) and of the corresponding 5-D metric
$$dS^2=d𝐫^2dt^2+d\lambda ^2$$
(4)
is the set of 5-D rotations and translations. But this group preserves nothing significant in space-time. One would like the new symmetry group to be related to some structure defined in space-time alone, to preserve some geometric entity of space-time.
The second reason that Eq. (2) is too naive is that the mass spectrum is continuous: $`0<m<\mathrm{}`$. But the whole mystery of particle mass spectra is that they consist of a few points with non-uniform spacing! Clearly a perfectly free particle field equation like (2) can never predict mass spectra of this type. We suggest that there should always be a self-force acting on the particle, whether or not it is acted on by external forces. The self-force must certainly involve the new coordinate $`\lambda `$, conjugate to mass.
The third reason that Eq. (2) is too naive is that it was simply written down ad hoc without any regard for the symmetry group of the new 5-D space. But as Bargmann and Wigner showed many years ago , the particle field equations now accepted — the scalar boson equation, the Dirac equation for spin 1/2 fermions, the photon field equation, etc. — correspond to the irreducible unitary representations of the Poincaré group $`𝐏`$, labelled by its two Casimir invariants spin $`j`$ and mass $`m`$, which uniquely fix these equations. Therefore the new symmetry group should have been chosen first, in accordance with the first criterion above, and then the field equations of the various particle species determined by its IUR’s.
Back to the first criterion: the present kinematical symmetry group of space-time is the Poincaré group $`𝐏`$, which preserves the space-time length element $`ds^2=g_{\mu \nu }dx^^\mu dx^^\nu d𝐫^2dt^2`$. One is thus motivated to search for the simplest and smallest extension of $`𝐏`$ which preserves something geometrical in space-time and has $`𝐏`$ as a subgroup. An immediate candidate is the group $`𝐂`$ which preserves space-time angle. By Liouville’s Theorem $`𝐂`$ is a 15-parameter Lie group composed of the 10-parameter subgroup $`𝐏`$, which preserves space-time length (and therefore space-time angle) augmented by a 5-parameter set of transformations which preserve space-time angle but not length.
To answer an expected immediate objection: of course $`𝐂`$’s transformations cannot act just on the 4-D space-time, with its length metric $`ds^2=g_{\mu \nu }`$ $`dx^^\mu dx^^\nu `$, because angle-preserving transformations of space-time do not in general preserve the length, and thus $`𝐂`$ would not be the symmetry group of this metric. (This was Einstein’s reason for rejecting the group $`𝐂`$, see .) The way to introduce the group $`𝐂^n`$ of conformal $`(`$ angle-preserving) transformations of $`n`$-dimensional euclidean space $`E^n`$ of coordinates $`x^^\mu ,\mu =1,2\mathrm{}n`$, was well-known to the great geometers of the nineteenth century (F. Klein, Liouville, Möbius, Lie et al.) some 150 years ago, but seems unknown today, at least to modern theoretical physicists. In brief, one introduces the $`(n+1)`$–dimensional space of spheres in $`E^n`$ characterized by their centers $`x^^\mu `$ and radii $`x^{n+1}`$. The group $`𝐂^n`$ is then that group of transformations $`x^^\alpha =f^^\alpha (x^1,x^2,\mathrm{}x^n,x^{n+1}),`$ $`\alpha =1,2,\mathrm{}n,n+1,`$ which preserve the angle $`\theta `$ under which two spheres $`x^^\alpha `$ and $`y^^\alpha `$ intersect, see Fig. 1. For infinitesimally close spheres $`y^^\alpha =x^^\alpha +dx^^\alpha `$ one gets
$$d\theta ^2=(x^{n+1})^2[(dx^1)^2+(dx^2)^2+\mathrm{}(dx^n)^2(dx^{n+1})^2]\text{ .}$$
(5)
(This is nothing but the Law of Cosines, familiar from plane geometry class in high school.) The expression (5) defines the metric (dimensionless angle metric) of the appropriate $`(n+1)`$-dimensional Riemannian space which has the conformal group $`𝐂^n`$ as its symmetry group. It turns out that this space is not flat but is of constant curvature. All of this is explained in exhaustive detail elsewhere .
Thus for the pseudo-euclidean space-time with $`n=4`$ we get the 5-dimensional space with metric
$$d\theta ^2=\lambda ^2(d𝐫^2dt^2+\sigma d\lambda ^2),\sigma =\pm 1\text{ ,}$$
(6)
with $`\{x^1,x^2,x^3,x^4\}`$ and $`x^5`$ renamed $`\{x^1,x^2,x^3,x^0\}`$ and $`\lambda `$ respectively. Of course the “sphere” $`x^^\alpha `$ is the hyperboloid $`g_{\mu \nu }(\xi ^^\mu x^^\mu )(\xi ^^\nu x^^\nu )+\sigma \lambda ^2=0`$ as a real locus. The sign $`\sigma `$, that is, whether the fifth dimension is spacelike $`(\sigma =+)`$ or timelike $`(\sigma =)`$ is left open for the moment.
This concludes the informal, qualitative part of this Introduction.
We show here how the field equation for spin 1/2 fermions in five dimensions coupled to a self-force dependent on the fifth coordinate predicts point mass spectra of just a few points and non-uniform spacing. If the Casimir invariant of this particular irreducible unitary representation has a certain range, it is a 3-point spectrum for isospin up or down. The spectrum is consistent with the experimental bounds on the isospin-down quarks $`d`$, $`s`$, and $`b`$ for values of the coupling constant $`\alpha `$ of order unity and range $`\kappa ^1`$ of the spatial wave functions of a few fermis.
To avoid a possible confusion at the outset: this 5-D theory has nothing to do with the Kaluza or Kaluza-Klein theories. The enlargement of space-time to a five-dimensional manifold is forced, not arbitrary, if the conformal group is demanded as the basic kinematical symmetry group . This fifth coordinate $`\lambda `$ is conjugate to mass just as position and time are conjugate to momentum and energy. Partial derivatives with respect to $`\lambda `$ replace mass terms in fermion and boson field equations. In solutions of gauge boson field equations $`\lambda `$ plays the role of a microscopic length “parameter” which modifies the usual space-time causality of point particles. It gives point particles a structure or extension in a certain sense .
We argue in this paper that this five-dimensional extension of special relativity (“conformal relativity”) is the natural framework for a theory of elementary particle mass. The results obtained here are promising but are only a first step; the main problem is the exact form of the quantum-mechanical self-force. Some extra points, including a puzzle, are made in the concluding remarks. These also include an argument that the 5-D theory gives a theoretical basis for some features of the electroweak theory which were postulated on the basis of experiment alone.
2. SOME BACKGROUND
As explained in the Introduction, the metric of conformal relativity is
$$d\theta ^2=\lambda ^2(dx^2+\sigma \text{ }d\lambda ^2)\text{ ,}$$
$$dx^2g_{\mu \nu }dx^^\mu dx^{\nu \text{ }},\mu ,\nu =0,1,2,3;\text{ }x^5\lambda ;\sigma =\pm \text{ ,}$$
(2.1)
where $`d\theta `$ is the infinitesimal angle under which spheres $`(x^^\mu ,\lambda )`$ and $`(x^^\mu +dx^^\mu ,`$ $`\lambda +d\lambda )`$ intersect. We use the metric $`g_{00}=g_{11}=g_{22}=g_{33}=+1`$. Whether the extra dimension is spacelike $`(\sigma =+)`$ or timelike $`(\sigma =)`$ is not yet clear, or maybe both occur. The ranges of the coordinates are $`\mathrm{}<x^^\mu <+\mathrm{}`$ as usual, and $`0<\lambda <\mathrm{}`$ (or possibly $`0<\lambda <\mathrm{})`$. The metric is singular if $`\lambda =0`$, so $`\lambda =0`$ is excluded from physical space, which is of course consistent with the action of the conformal group C . We call these two 5-D Riemannian spaces (2.1) K<sub>+</sub> and K<sub>-</sub> (after Felix Klein).
The field equation for spin 1/2 fermions in the C-covariant theory is<sup>1</sup><sup>1</sup>1Eq. (2.2) here is Eq. (4.3) of the second article of Ref. , where $`\nu (4/9)q_3`$. Note that these articles considered only the case $`\sigma =+`$. Much of the physical discussion there is dated.
$$(\gamma ^^\alpha __\alpha +\gamma \beta _7\nu )\text{ }\psi =0\text{ , }__\alpha \underset{_\alpha }{\overset{\gamma }{}}ig\text{ }A__\alpha \text{ .}$$
(2.2)
Here the six anticommuting $`\gamma `$-matrices obey
$$\gamma __\alpha \gamma __\beta +\gamma __\beta \gamma __\alpha =2\gamma _{\alpha \beta }\text{ }\mathrm{𝟏}\text{}\gamma __\alpha \gamma +\gamma \gamma __\alpha =0\text{}\gamma ^2=\mathrm{𝟏}\text{ ,}$$
(2.3a)
$$\beta _7i\lambda ^5\gamma _1\gamma _2\gamma _3\gamma _0\gamma _5\gamma \text{ ; }\alpha ,\beta =0,1,2,3,5\text{ ,}$$
(2.3b)
where $`\gamma _{\alpha \beta }`$ is the angle metric (2.1). Indices are raised and lowered with this metric. $`\underset{_\alpha }{\overset{\gamma }{}}`$ is the covariant derivative on spinors $`\psi `$ which fixes the spin algebra $`\gamma __\alpha ,\gamma `$. (Note that the spaces $`K__\sigma `$ are not flat, so that covariant derivatives occur in field equations.) We consider here only a $`U(1)`$ internal symmetry with gauge boson $`A__\alpha `$. The equation (2.2) is uniquely fixed by requiring that the solutions $`\psi `$ span an irreducible unitary representation $`(IUR)`$ of C. The parameter $`\nu `$ is a Casimir invariant for this $`IUR`$, and Eq. (2.2) is the sole independent condition for spin 1/2 . The six $`\gamma __\alpha ,\gamma `$ are $`8\times 8`$ and $`\psi `$ is an 8-spinor because eight is the minimum dimension allowed for a matrix representation of the algebra (2.3a). When the spin connection is inserted, Eq. (2.2) reduces to
$$(\stackrel{~}{\gamma }D+\stackrel{~}{\gamma }^5D_5+2\sigma \lambda \stackrel{~}{\gamma }_5+\nu \beta _7)\psi =0\text{ ,}$$
$$\stackrel{~}{\gamma }__\alpha \gamma \gamma __\alpha \text{ , }D__\alpha __\alpha igA__\alpha \text{ ,}$$
(2.4)
where the will always mean the 4-D scalar product $`\stackrel{~}{\gamma }D\stackrel{~}{\gamma }^^\mu D__\mu `$. Note that $`D__\alpha `$ involves the ordinary partial derivative $`__\alpha `$ ; the third term in Eq. (2.4) comes from the spin connection.
To be able to calculate with Eq. (2.4) a representation of the six $`8\times 8`$ matrices $`\gamma __\alpha ,\gamma `$ must of course be chosen. We choose $`\gamma __\alpha =\gamma \stackrel{}{\gamma __\alpha }`$ and
$$\underset{_\mu }{\overset{}{\gamma }}=\lambda ^1\left(\begin{array}{cc}\gamma __\mu \hfill & 0\hfill \\ 0\hfill & \gamma __\mu \hfill \end{array}\right),\underset{5}{\overset{}{\gamma _\text{ }}}\text{ }=\lambda ^1\left(\begin{array}{cc}h\hfill & 0\hfill \\ 0\hfill & h\hfill \end{array}\right),\gamma =\left(\begin{array}{cc}0\hfill & 1\hfill \\ 1\hfill & 0\hfill \end{array}\right).$$
(2.5a)
The $`\stackrel{\text{ }\alpha }{\stackrel{}{\gamma }}`$are obtained by raising the indices with the metric (2.1). For the $`4\times 4`$ $`\gamma __\mu ,`$ $`h,`$ and $`1`$ in these matrices, see Eq. (2.6). Then $`\beta _7`$, Eq. (2.3b), is
$$\beta _7=\left(\begin{array}{cc}1\hfill & 0\hfill \\ 0\hfill & 1\hfill \end{array}\right)\text{ .}$$
(2.5b)
It can be shown (unpublished) that by comparing a Lagrangian for the spin 1/2 field equation (2.4) with the Lagrangian for the electroweak theory (, Chap. 7) that we can identify the upper and lower 4-spinors in the 8-spinor $`\psi `$ as the $`T_3=+1/2`$ and $`1/2`$ components of the isodoublets of the electroweak theory in this representation. In fact, the whole electroweak theory can be reproduced. More on this in Sec. 4. Therefore we call the representation (2.5) the EW (electroweak) representation. The field equation (2.4) written in the $`EW`$ representation splits cleanly into wave equations for the $`T_3=+1/2`$ and $`1/2`$ components (there is no coupling between these fields) and further, these wave equations are identical.
This common wave equation for the case $`\sigma =`$ is
$$\{\gamma (ig\text{ }A)ih(_5ig\text{ }A_5)+(\nu +2ih)/\lambda \}\psi =0\text{ ,}$$
$$\gamma __\mu \gamma __\nu +\gamma __\nu \gamma __\mu =2g_{\mu \nu }\mathrm{𝟏}\text{ , }hi\gamma _1\gamma _2\gamma _3\gamma _0\text{ .}$$
(2.6)
Here the $`\gamma __\mu `$ are the usual $`4\times 4`$ constant $`\gamma `$-matrices, $`\psi `$ is now a $`4`$-spinor, and $`h`$ is the handedness operator (usually called $`\gamma _5`$ in the literature): $`h\psi _L=\psi _L`$, $`h\psi _R=+\psi _R`$ for left and right-handed spinors.
The field equation for the gauge boson $`A__\alpha `$ is
$$\underset{_\alpha }{\overset{\gamma }{}}F__\beta ^^\alpha =0\text{ , }F_{\alpha \beta }__\alpha A__\beta __\beta A_{{}_{\alpha }{}^{}\text{ }}\text{.}$$
(2.7)
These are reduced to a set of partial differential equations for the $`5`$-vector $`A__\alpha `$ in Ref. .
3. FERMION MASS SPECTRUM FOR A TIMELIKE FIFTH DIMENSION
We look at stationary states: $`\psi (t,𝐫,\lambda )=e^{iEt}g(𝐫,\lambda )`$ of Eq. (2.6). If we insert a self-force $`A__\alpha ^{SF}`$ and solve for a resting spin 1/2 fermion, the energy spectrum should be the mass spectrum: $`E=M`$. The self-force should certainly involve the fifth coordinate $`\lambda `$, so we adopt provisionally
$$A_0^{SF}=g^{}/\lambda \text{, other }A__\alpha ^{SF}0\text{ .}$$
(3.1)
More on this in Sec. 4. Then the equation becomes
$$\left\{\gamma ^0(M\alpha /\lambda )+i\gamma +h__\lambda +(i\nu 2h)/\lambda \right\}\text{ }g(𝐫,\lambda )=0\text{}$$
(3.2)
Here $`\alpha g^{}g`$ ($`g^{}=g`$ is natural for a self-force, but we leave this open for generality.) Consider $`s`$-states $`g(r,\lambda )`$ only; then $`i\gamma `$ becomes $`i\gamma _r_r`$ where $`\gamma _r\gamma 𝐧`$, $`𝐧`$ a unit 3-vector. We seek a separable solution in $`r`$ and $`\lambda `$, so take $`g(r,\lambda )=e^{\kappa r}g(\lambda )`$ with $`\kappa `$ real and positive. Eq. (3.2) then reduces to the ordinary differential equation in $`\lambda `$
$$\left\{\gamma ^0(M\alpha /\lambda )i\kappa \gamma _r+h__\lambda +(i\nu 2h)/\lambda \right\}g(\lambda )=0.$$
(3.3)
The solution is given in the Appendix. It is formally very similar to the solution of the Dirac equation for the relativistic hydrogen atom with $`\lambda `$ and the mass levels of the particle playing the roles of $`r`$ and the hydrogenic energy levels, respectively. (The spectrum is very different however.) The mass spectrum is
$$M_{(n^{},\tau )}/\kappa =S__\tau +n^{}\mathrm{}[\alpha ^2(S__\tau +n^{})^2]^{1/2}$$
$$n^{}=0,1,2,3,\mathrm{}\text{ },\text{ }\tau =\pm \text{ },$$
(3.4a)
$$S__\tau \tau (\alpha ^2\nu ^2)^{1/2}\text{ ,}$$
(3.4b)
$$S__\tau +n^{}\text{ has the sign of }\alpha g^{}g\text{ ,}$$
(3.4c)
norm restriction<sup>2</sup><sup>2</sup>2 For the 4-spinor $`\psi `$ the norm is $`\psi ^2d^3r_0^{\mathrm{}}𝑑\lambda \lambda ^4\stackrel{}{\psi }\gamma _0\psi `$, $`t=const`$., where $`\gamma _0`$ is the constant $`4\times 4`$ matrix. For this “bound” solution we require $`\psi ^2<\mathrm{}`$. The bound (3.4d) on $`\gamma `$ comes from requiring the $`\lambda `$-integral to converge at its lower limit $`\lambda =0`$.:
$$(\alpha ^2\nu ^2)^{1/2}<1/2\text{ for }\tau =\text{ .}$$
(3.4d)
One can see first in a general sort of way that this is a finite point spectrum: when the radicand in the denominator of Eq. (3.4a) goes negative, the spectrum ends. In fact, if we choose $`\gamma (\alpha ^2\nu ^2)^{1/2}`$ as a convenient independent variable (do not confuse this $`\gamma `$ with the matrix $`\gamma `$ in Eq. (2.3)!) and set $`F(\gamma ;n^{},\tau )\alpha ^2(S__\tau +n^{})^2`$, we get, on expanding and cancelling etc.
$$F(\gamma ;n^{},\tau )=2n^{}\tau \gamma +\nu ^2n^2\text{ .}$$
(3.5)
Now choose $`g^{}=g`$, or $`\alpha g^2>0,`$ as seems natural. Then the necessary and sufficient conditions for a spectral point $`(n^{},\tau )`$ are
$$\gamma <(\nu ^2n^2)\mathrm{}2n^{}\text{ , }\tau =+\text{ ; }\gamma >(n^2\nu ^2)\mathrm{}2n^{},\text{ }\tau =\text{ , }$$
(3.6a)
$$\gamma <n^{}\text{ for }\tau =\text{ },$$
(3.6b)
$$\gamma <1/2\text{ for }\tau =.$$
(3.6c)
These are respectively from $`F(\gamma ;n^{}\tau )>0`$, Eq. (3.4c) for $`\alpha >0`$, and Eq. (3.4d).
In modern particle theory there are three families (isodoublets) of quarks and three of leptons. Relevant to this, the following theorem can be proved from the conditions (3.6a, b, c):
Theorem. There are three and only three mass levels if and only if $`1<\nu ^2<2`$. These levels are $`(n^{},\tau )=(0,+)`$, $`(1,)`$, and $`(1,+)`$.
The mass spectrum written in terms of $`\gamma `$ is
$$M_{(n^{},\tau )}\mathrm{}\kappa =(\tau \gamma +n^{})\mathrm{}(2n^{}\tau \gamma n^2+\nu ^2)^{1/2}$$
(3.7)
from just above. Thus for the three levels $`(0,+)`$, $`(1,)`$, and $`(1,+)`$ we get
$$\begin{array}{ccc}M_{(0,+)}\mathrm{}\kappa \hfill & =\hfill & \gamma \mathrm{}\nu \text{ ,}\hfill \\ & & \\ M_{(1,)}\mathrm{}\kappa \hfill & =\hfill & (1\gamma )\mathrm{}(2\gamma 1+\nu ^2)^{1/2}\text{ ,}\hfill \\ & & \\ M_{(1,+)}\mathrm{}\kappa \hfill & =\hfill & (1+\gamma )\mathrm{}(2\gamma 1+\nu ^2)^{1/2}\text{ ,}\hfill \\ & & \\ 1<\nu ^2<2\hfill & ,\hfill & 0<\gamma <1/2\text{ }.\hfill \end{array}$$
(3.8)
Then from these expressions one can deduce that the only possibility that one mass is much greater than the other two is
$$\nu ^2=1+2\gamma +\epsilon \text{ , }0<\epsilon <<1\text{ ,}$$
(3.9)
in which case $`M_{(1,+)}`$ is the large one. (This assumes $`\epsilon <<\gamma .`$)
Fitting the quarks. We try to fit the $`T_3=1/2`$ set of quarks $`d`$, $`s`$, and $`b`$. The experimental mass limits in $`MeV`$ are
$$M_d=39\text{}M_s=60170\text{}M_b=41004400.$$
(3.10)
So we adopt the value (3.9) for $`\nu ^2`$ and identify $`(1,+)b`$. Next, inserting $`\nu ^2`$ (3.9) into the mass formulae (3.8) and neglecting $`\epsilon `$ in $`(0,+)`$ and $`(1,)`$, we get the ratio
$$M_{(1,)}\mathrm{}M_{(0,+)}=(1\gamma )(1+2\gamma )^{1/2}\mathrm{}2\gamma ^{3/2}\text{ .}$$
(3.11)
It can be checked that this ratio is always $`>1`$ for $`0<\gamma <1/2`$, so we choose $`(1,)s`$ and $`(0,+)d`$. Now equate the ratio (3.11) to $`M_s/M_d`$, using the average values $`M_d=6`$ $`MeV`$ and $`M_s=115`$ $`MeV`$. The resulting equation
$$(1\gamma )(1+2\gamma )^{1/2}=38.4\text{ }\gamma ^{3/2}$$
(3.12)
has the solution $`\gamma .088`$. Finally, to determine $`\epsilon `$, set the theoretical and experimental ratios $`M_b/M_d`$ equal. This gives
$$(1+\gamma )\nu \mathrm{}\gamma \epsilon ^{1/2}=(M_b/M_d)\text{ }_{\mathrm{exp}tl\text{ }}\text{.}$$
(3.13)
Insert $`\gamma =.088`$ and $`\nu =1.088`$ and use the minimum value $`4100/9455`$ for the ratio on the right to get the maximum size of $`\epsilon `$. This gives $`\epsilon _{\mathrm{max}}8.7\times 10^4`$, and verifies our assumption $`\epsilon <<\gamma `$.
The values of the coupling constant $`\alpha `$ and the range $`\kappa ^1`$ of the spatial wave functions are also of interest. We can evaluate $`\kappa `$ from $`\kappa (M_{(n^{},\tau )}/\kappa )=(M_q)_{\mathrm{exp}tl}`$. If we use the same average values for $`M_d`$ and $`M_s`$ as used above to determine $`\gamma `$, we will get the same $`\kappa `$ for either $`(1,)`$ or $`(0,+)`$. Choose $`(0,+).`$
$$\kappa \gamma /\nu \text{ }=.081\kappa =6\text{ }MeV\kappa =74.2\text{ }MeV\text{}$$
which gives $`\kappa ^1200/74.22.7`$ $`f`$. Also $`\alpha ^2=\gamma ^2+\nu ^21.18`$, or $`\alpha 1.09`$, which suggests a self-force of QCD origin.
In summary, a fit to the three isospin-down quarks $`d`$, $`s`$, and $`b`$ has been obtained as the levels
$$(0,+)d\text{}(1,)s\text{ ,}(1,+)b$$
(3.14a)
for the Casimir invariant $`\nu ^21.176`$ and the reasonable values of the physical parameters
$$\alpha 1.09\text{ and }\kappa ^12.7\text{ }f\text{ .}$$
(3.14b)
Of course nearby values of these parameters will also give a fit owing to the wide latitude (3.10) in the experimental masses.
4. CONCLUDING REMARKS
A further characteristic of this theory necessary in any theory of mass should be mentioned. In inelastic scattering of elementary particles, energy and momentum are conserved but mass is not. Thus in any theory which unifies these quantities in some sense mass must be qualitatively different from energy and momentum and so must the conjugate quantities. Now note that the fifth coordinate $`\lambda `$ is qualitatively different from the other four $`x^^\mu `$; look for example at the metric (2.1). Further, the symmetry group C includes translation groups on $`𝐫`$and $`t`$, hence momentum and energy are conserved in particle scattering . But there is no translation group on $`\lambda `$ , so the conjugate quantity mass need not be conserved.
The mass spectrum analyzed in Sec. 3 does fit the experimental numbers for the quarks, at least to within their (very loose) bounds. However, this spectrum is not intended to be final and quantitative at this stage. We only meant to show here that this particular 5-D theory required by conformal symmetry is capable of predicting few-point mass spectra of the right order of magnitude. The main problem is the crudity of the self-force (3.1) adopted. This field does not in fact satisfy the boson field equations (2.7) (see Ref. ) and must therefore be thought of as an approximation to an actual solution<sup>3</sup><sup>3</sup>3The boson field equations (2.7) have the Coulombic solution $`A_0=g^{}\mathrm{}\sqrt{\lambda ^2r^2},`$ $`0r<\lambda ;`$ $`=g^{}\mathrm{}\sqrt{r^2\lambda ^2},`$ $`0<\lambda <r<\mathrm{}`$, other $`A__\alpha 0`$. or simply as a model. A quantitative theory needs a realistic self-force, perhaps one involving also $`SU(2)`$ gauge bosons.
A few other points, including some puzzles, will be mentioned.
1) The signature $`\sigma =`$ was needed for an interesting mass spectrum. We can show that for $`\sigma =+`$ a one-point spectrum results for $`\alpha >0`$ (unpublished). The puzzle here is that $`\sigma =+`$ is definitely indicated in the classical self-force theory , which successfully resolves the anomalies due to classical point particles.
2) Notice that if the lepton self-force is electromagnetic: $`\alpha 1/137,`$ the mass spectrum (3.4) cannot fit the $`T_3=1/2`$ leptons $`e`$, $`\mu `$, and $`\tau `$ since then $`\gamma (\alpha ^2\nu ^2)^{1/2}`$ is pure imaginary for $`1<\nu ^2<2`$. This is a puzzle. But we add that for $`\sigma =+,`$ $`(\nu ^2\alpha ^2)^{1/2}`$ occurs where $`(\alpha ^2\nu ^2)^{1/2}`$ occurs for $`\sigma =`$, hence the equation (3.2) written for $`\sigma =+`$ with a better self-force than (3.1) might work.
3) For perfectly free spin 1/2 fermions (no external force and no self-force) the field equation (2.4) with $`A__\alpha 0`$, $`\sigma =+`$ or $``$, space-time dependence in $`e^{ipx}`$ with $`p^2+m^2=0`$, and $`\gamma __\alpha `$ and $`\gamma `$ in the EW representation is easily solved. The $`\lambda `$-dependence is in factors $`\lambda ^{5/2}Z_{\mu _L}(m\lambda )`$ and $`\lambda ^{5/2}Z_{\mu _R}(m\lambda )`$ for the $`L`$\- and $`R`$-handed components of $`\psi `$, with $`\mu _L\mu _R`$. The $`Z__\mu `$ are cylinder functions of order $`\mu `$. The mass spectrum is continuous, $`0m<\mathrm{}`$. In the case $`\sigma =+`$ if $`\nu =1/2`$ is chosen for the Casimir invariant, then in the limit $`m0`$ (neutrino solution) only a left-handed neutrino survives. This makes the value $`\nu =1/2`$ very attractive theoretically for leptons. Perfectly free fermions are unphysical because of the continuous mass spectrum. But this also supports the idea that the mass problem for leptons should be phrased in the space $`\sigma =+`$ (cf. point (2) above) with $`\nu =1/2`$.
4) As indicated briefly above, this theory based on C instead of P as the kinematical symmetry group of particle physics is compatible with the EW theory. Further, it furnishes a theoretical foundation for some of the features of that theory adopted on the basis of experiment. Consider the following points. (a) The six basic anticommuting $`\gamma `$-matrices (2.3a) demand an 8-dimensional spinspace, thus allowing the upper and lower 4-spinors to be identified with the $`T_3=\pm 1/2`$ isodoublets. (b) But more than this, in the differential operator involving the primary gauge bosons B$`__\alpha `$ and W$`__\alpha `$ <sup>i</sup>
$`(i=1,2,3)`$, the spin algebra of the $`SU(2)\times U(1)`$ internal symmetry group is formed entirely from the $`8\times 8`$ $`\gamma `$-matrices (2.3a,b). Define the matrices
$$\tau _i\gamma \text{ , }\tau _2i\gamma \beta _7\text{ , }\tau _3\beta _7\text{ .}$$
(4.1)
Then these have the same commutation relations as the Pauli matrices. Further, in the EW representation (2.5) they take exactly the standard form, where the 1’s and 0’s are $`4\times 4`$. Contrast this with the situation in the present day EW theory where generators of the internal symmetry group $`SU(2)`$, unrelated to the $`\gamma __\mu `$, are imported from the outside. The handedness projections $`P_h^{}`$, $`h^{}=\pm `$, are built from the $`8\times 8`$ $`H\lambda \beta _7\stackrel{~}{\gamma }_5`$, which takes the form
$$H=\left(\begin{array}{cc}h\hfill & 0\hfill \\ 0\hfill & h\hfill \end{array}\right)\text{ ,}$$
(4.2)
where $`h`$ is the $`4\times 4`$ handedness operator (see below Eq. (2.6)), in the EW representation. (c) If the Lagrangian
$$=\overline{\psi }\left[\stackrel{~}{\gamma }D+\stackrel{~}{\gamma }\text{ }^5\text{ }D_5+2\sigma \lambda \stackrel{~}{\gamma }_5+\nu \beta _7\right]\psi \text{ ,}$$
(4.3)
which yields the field equation (2.4), is equipped with the gauge bosons $`B__\alpha `$ and $`W__\alpha `$ <sup>i</sup>, it exactly reproduces the Lagrangian of the EW theory (, Chap. 7) plus some extra terms coming from the fifth components $`B_5`$ and $`W_5`$ <sup>i</sup>, presumably small corrections to the $`4D`$ theory. Then the standard mixing produces the photon and $`Z`$ fields. d) However, the aspect in which this theory is not compatible with the EW theory (or the whole Standard Model) is the main point of this paper. In this theory the fermions may be massive, like the quarks considered in this paper. The fifth dimension plus an appropriate self force provides the masses. The Higgs mechanism is unnecessary.
APPENDIX. SOLUTION FOR THE MASS EIGENSTATES AND SPECTRUM
Insert the formally $`2\times 2`$ representation
$$\gamma ^0=i\left(\begin{array}{cc}0\hfill & 1\hfill \\ 1\hfill & 0\hfill \end{array}\right)\text{}\gamma _r=i\left(\begin{array}{cc}0\hfill & 1\hfill \\ 1\hfill & 0\hfill \end{array}\right)\text{}h=\left(\begin{array}{cc}1\hfill & 0\hfill \\ 0\hfill & 1\hfill \end{array}\right)$$
(A1)
and $`g(\lambda )=\left(\begin{array}{c}F\hfill \\ G\hfill \end{array}\right)`$ into Eq. (3.3). $`F`$ and $`G`$ are thus 2-spinors; in fact $`F=g_R`$ and $`G=g_L`$ in view of the form (A1) of the handedness operator $`h`$. Multiplying by $`i`$ we get
$$\begin{array}{c}(Mi\kappa \alpha /\lambda )Gi(__\lambda +(i\nu 2)/\lambda )F=0\text{ ,}\hfill \\ \\ (M+i\kappa \alpha /\lambda )F+i(__\lambda (i\nu +2)/\lambda )G=0\text{ .}\hfill \end{array}$$
(A2)
Rephase: $`iFF`$, $`GG`$. Define
$$\beta _1M+i\kappa \text{}\beta _2Mi\kappa \text{}\beta ^2\beta _1\beta _2=M^2+\kappa ^2\text{.}$$
(A3)
Divide equations $`(A2)`$ by $`\beta \sqrt{\beta ^2}`$ and put $`\beta \lambda \tau `$.
$$\begin{array}{c}(\beta _2/\beta \alpha /\tau )G(__\tau +(i\nu 2)/\tau )F=0\text{ ,}\hfill \\ \\ (\beta _1/\beta \alpha /\tau )F(__\tau (i\nu +2)/\tau )G=0\text{ .}\hfill \end{array}$$
(A4)
Set $`F,Ge^\tau (f,g)`$. Then $`__\tau F=(\dot{f}f)e^\tau `$ etc. where $`{}_{}{}^{}/\tau `$. Solve the equations in terms of $`f`$ and $`g`$ by the power series
$$f=\tau ^s\underset{n=0}{\overset{\mathrm{}}{}}a_n\tau ^n\text{ , }g=\tau ^s\underset{n=0}{\overset{\mathrm{}}{}}b_n\tau ^n\text{}a_0\text{ and }b_00\text{ .}$$
(A5)
When these power series are inserted into the equations for $`f`$ and $`g`$ and coefficients of $`\tau ^{s+n1}`$ equated to $`0`$, we obtain
$$\begin{array}{c}(\beta _2/\beta )b_{n1}\alpha b_n(s+n)a_n+a_{n1}(i\nu 2)a_n=0\text{ ,}\hfill \\ \\ (\beta _1/\beta )a_{n1}\alpha a_n(s+n)b_n+b_{n1}+(i\nu +2)b_n=0\text{ .}\hfill \end{array}$$
(A6)
Multiply the top equation (A6) by $`\beta _1/\beta `$ and subtract the bottom equation. The terms $`a_{n1}`$ and $`b_{n1}`$ go out since $`\beta _1\beta _2/\beta ^2=1`$. After rearrangement this gives
$$\left[(\beta _1/\beta )(s+n2+i\nu )\alpha \right]a_n=\left[s+n2i\nu \beta _1\alpha /\beta _2\right]b_n\text{.}$$
(A7)
To get the indicial equation choose $`n=0`$ in Eq. (A6) and ignore the terms $`a_1`$ and $`b_1`$. The determinant must vanish so that nonzero $`a_0`$ and $`b_0`$ result; the result is
$$S__\eta s__\eta 2=\eta (\alpha ^2\nu ^2)^{1/2}\text{}\eta =\pm \text{ .}$$
(A8)
(We have changed the subscript $`\tau `$ on $`S__\tau `$, Eq. (3.4b), to $`\eta `$ so as not to confuse it with the $`\tau \beta \lambda `$ of Eq. (A4) et seq.) This is Eq. (3.4b). By letting $`n\mathrm{}`$ in Eq. (A7) we get $`b_n=(\beta _1/\beta )a_n`$ in this limit; substituting this into both equations (A6) for $`n\mathrm{}`$, we find $`a_n/a_{n1}=2/n`$ and the same for the $`b`$$`s`$ in this limit. Thus both series (A5) diverge like $`e^{2\tau }`$, which is not allowed by the assumed finiteness of the norm. Hence both series must terminate:
$$a_{n^{}+1}=b_{n^{}+1}=0\text{ , }n^{}=0,1,2,\mathrm{}\text{ .}$$
(A9)
Set $`n=n^{}+1`$ in Eq. (A6); we get $`b_n^{}=(\beta _1/\beta )a_n^{}`$ . Put this result into Eq. (A7) for $`n=n^{}`$. After cancellation of some terms and rearrangement
$$2(\beta _1/\beta )(s+n^{}2)\alpha (1+(\beta _1/\beta )^2)=0$$
(A10)
results. Divide this by $`2\beta _1/\beta `$ and use $`\beta _1/\beta =\beta /\beta _2`$. After some algebra we obtain
$$S__\eta +n^{}=\alpha M/\beta \text{ .}$$
(A11)
(This implies Eq. (3.4c).) Finally, do some algebra on Eq. (A11), using $`\beta \sqrt{M^2+\kappa ^2}`$, to solve for $`M`$. This gives the mass spectrum (3.4a).
The mass eigenstates. From Sec. 3 and this Appendix, the mass eigenstates are $`\psi =\left(\begin{array}{c}F\hfill \\ G\hfill \end{array}\right)`$, where the $`2`$-spinors $`F`$ and $`G`$ are
$$F(t,r,\lambda )=e^{iMt}\text{ }e^{\kappa r}\text{ }F(\lambda )\text{ , }G(t,r,\lambda )=e^{iMt}\text{ }e^{\kappa r}G(\lambda )\text{ ,}$$
(A12)
$$\begin{array}{c}F(\lambda )=(i)e^\tau \tau ^{s__\eta }\underset{n=0}{\overset{n^{}}{}}a_n\tau ^n\times u_+\text{ ,}\hfill \\ \\ G(\lambda )=e^\tau \tau ^{s__\eta }\underset{n=0}{\overset{n^{}}{}}b_n\tau ^n\times u_{}\text{ ,}\hfill \end{array}$$
(A13)
where
$$\tau \beta \lambda =\left[(M_q/\kappa )^2+1\right]^{1/2}\kappa \lambda \text{ .}$$
(A14)
Here the quantum number of the eigenstate $`q(n^{},\eta )`$ and $`M_q/\kappa `$ is given by Eq. (3.4a) with the sign $`\tau `$ changed to $`\eta `$. The relation of the $`b_n`$ to $`a_n`$ and the $`a_n`$ to the $`a_{n1}`$ are given by Eqs. (A7) and (A6). The constant $`2`$-spinors $`u_+`$ and $`u_{}`$ are normalized in some way; the overall normalization of the $`4`$-spinor $`\psi `$ is secured by the free parameter $`a_0`$.
REFERENCES
G. Kane, Modern Elementary Particle Physics (Addison-Wesley, New York, 1993), Chap. 8.
P.C.W. Davies and J. Brown, Superstrings, a Theory of Everything? (Cambridge University Press, Cambridge, 1988).
B. Greene, The Elegant Universe (Norton, New York, 1999).
V. Bargmann and E.P. Wigner, Proc. Nat. Acad. Sci. 34, 211 (1948).
F. Klein, Math. Ann. 5, 257 (1872). \_______, Vorlesungen über Höhere Geometrie 3 Aufl, (Springer, Berlin 1926).
A. Einstein, The Meaning of Relativity, 4th ed. (Princeton University Press, Princeton 1953).
R. L. Ingraham, Int. J. Mod. Phys. 7, 603 (1998).
\[8) R. L. Ingraham, Nuovo Cimento 68B, 203, 218 (1982).
L.I. Schiff, Quantum Mechanics, 2nd ed. (McGraw-Hill, New York 1955), Sec. 44.
Review of Particle Physics, Euro. Phys. J. C3 (1998), p. 24.
J.M. Jauch and F. Rohrlich, The Theory of Photons and Electrons (Addison-Wesley, Cambridge U.S.A. 1955), Sec. 1-11.
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# Scale Invariance in the Nonstationarity of Physiological Signals
## Abstract
We introduce a segmentation algorithm to probe temporal organization of heterogeneities in human heartbeat interval time series. We find that the lengths of segments with different local values of heart rates follow a power-law distribution. This scale-invariant structure is not a simple consequence of the long-range correlations present in the data. We also find that the differences in mean heart rates between consecutive segments display a common functional form, but with different parameters for healthy individuals and for patients with heart failure. This finding may provide information into the way heart rate variability is reduced in cardiac disease.
A time series is stationary if the mean, standard deviation and all higher moments, as well as the correlation functions, are invariant under time translation . Signals that do not obey these conditions are nonstationary. Nonstationarity is a prominent feature of biological variability that can be associated with regimes (segments) of different statistical properties. The borders between different segments can be gradual or abrupt (Fig. 1).
A major problem in contemporary physiology is the presence of nonstationarity in time series generated under free-running conditions . Physiological signals obtained under widely-varying conditions raise serious challenges to both technical and fundamental aspects of time series analysis. By filtering out effects of nonstationarity, much work has focused on “intrinsic properties” of physiological signals . This approach is based on the implicit assumption that the nonstationarity arises simply from changes in environmental conditions — e.g., different daily activities — so environmental “noise” could be treated as a “trend” and distinguished from the more subtle fluctuations that may reveal intrinsic correlation properties of the dynamics. Indeed, important scale-invariant features in physiological processes were recently revealed after filtering out masking effects of nonstationarity . However, nonstationarity itself is also an important feature of physiological time series and is known to change from healthy to pathological conditions , suggesting more than only enviromental conditions are reflected in the phenomena. Thus one would expect that there is a non-trivial structure associated with the nonstationarity in physiological signals, which may change with disease. To test this hypothesis we focus on one statistical property, the mean heart rate, which is related to physiologic responses and is commonly used for medical evaluation.
The problem is to partition a nonstationary time series, which is composed of many segments with different mean value, in such a way as to maximize the difference in the mean values between adjacent segments. We apply the following procedure: we move a sliding pointer from left to right along the signal. At each position of the pointer, we compute the mean of the subset of the signal to the left of the pointer ($`\mu _{\mathrm{left}}`$) and to the right ($`\mu _{\mathrm{right}}`$). To measure the difference between $`\mu _{\mathrm{left}}`$ and $`\mu _{\mathrm{right}}`$, we compute the t-statistic :
$$t\left|\frac{\mu _{\mathrm{left}}\mu _{\mathrm{right}}}{s_D}\right|$$
(1)
where $`s_D`$ is the pooled variance .
We next determine the position of the pointer for which $`t`$ reaches its maximum value, $`t_{\mathrm{max}}`$, and compute the statistical significance of $`t_{\mathrm{max}}`$ . We check if this significance exceeds a given threshold $`𝒫_0`$. If so, then the signal is cut at this point into two subsequences; otherwise the signal remains undivided. If the sequence is cut, the procedure continues recursively for each of the two resulting subsequences created by each cut. Before a new cut is accepted, we also compute $`t`$ between the right-hand new segment and its right neighbor (obtained by a previous cut) and the $`t`$ between the left-hand new segment and its left neighbor (also obtained by a previous cut) and check if both values of $`t`$ have a statistical significance exceeding $`𝒫_0`$. If so, we proceed with the new cut; otherwise we do not cut. This ensures that all resulting segments have a statistically significant difference in their means. The process stops when none of the possible cutting points has a significance exceeding $`𝒫_0`$, and we say that the signal has been segmented at the “significance level $`𝒫_0`$” (Fig. 2).
Our method leads to partitioning of a time series into segments with well-defined means, each significantly different from the mean of the adjacent segments (Fig. 1). This allows us to probe the nonstationarity in a signal through the statistical analysis of the properties of the segments.
Here we consider 47 datasets from 18 healthy subjects, 17 records of cosmonauts during orbital flight and 12 patients with congestive heart failure . We separately analyze 6–hour long subsets of each dataset, corresponding to the periods when the subject is awake or sleeping. Figure 1 shows a representative dataset of a healthy subject, and a subject with heart failure. Superposed on the interbeat interval series, we also plot the segments obtained by means of our segmentation algorithm.
To quantify the nonstationarity in heart rate variability, we study the statistical properties of the segments corresponding to parts of the signal with significantly different mean values. To characterize the segments, we analyze two quantities: (i) the length of the segments; (ii) the absolute values of the differences between the mean values of consecutive segments, which we call jumps.
(i) Distribution of segment lengths — Healthy subjects typically exhibit nonstationary behavior associated with large variability, trends, and segments with large differences in their mean values, while data from heart failure subjects are characterized by reduced variability and appear to be more homogeneous (Fig. 1) . Thus, one might naively expect that signals from healthy subjects will be characterized by a large number of segments, while signals from heart failure subjects will exhibit a smaller number of segments (i.e., the average length of the segments for healthy subjects could be expected to be smaller than for heart failure subjects).
To test this hypothesis, we apply the segmentation algorithm to 6–hour records of interbeat intervals during daily activity, and find that for each healthy subject the distribution of segment lengths is well described by a power law with an identical exponent, indicating absence of a characteristic length for the segments. Surprisingly, we find that this power law remains unchanged for records obtained from cosmonauts during orbital flight (under conditions of microgravity) and for patients with heart failure (Fig. 3). A similar common type of behavior is also observed from 6–hour records during sleep for all three groups .
To verify the results of the segmentation procedure, we perform several tests. First, we check the validity of the observed power law in the distribution of segment lengths. We generate a surrogate signal formed by joining segments of white noise with standard deviation $`\sigma =0.5`$, and mean values chosen randomly from the interval $`[0,1]`$. We choose the lengths of these segments from a power-law distribution with a given exponent. Even when the difference between the mean values of adjacent segments is smaller than the standard deviation of the noise inside the segments, we find that our procedure partitions the surrogate signal into segments with lengths that reproduce the original power-law distribution \[Fig. 4(a)\]. This test shows that the distributions obtained after segmenting surrogate data with similar values of their exponents, appear clearly different from each other, making more plausible that the distributions obtained for the lengths of the segments for the healthy, cosmonauts and congestive heart failure subjects (Fig.3) follow indeed an identical distribution.
Second, we test if the observed power-law distribution for the segment lengths is simply due to the known presence of long-range correlations in the heartbeat interval series . For that, we generate correlated linear noise with the same correlation exponent as the heartbeat data. We find that the distribution of segment lengths obtained for the linear noise differs from the distribution obtained for the heartbeat data \[Fig. 4(b)\]. For the noise, the distribution decays faster, which means that these signals are more segmented than the heart data. In fact, for different linear noises with a broad range of correlation exponents, we do not find power-law behavior in the distribution of the segments. Thus we conclude that the linear correlations are not sufficient to explain the power-law distribution of segment lengths in the heartbeat data.
(ii) Differences between the mean values of consecutive segments (jumps) — Different healthy records can be characterized by different overall variance, depending on the activity and the individual characteristics of the subjects. Moreover, subjects with heart failure exhibit interbeat intervals with lower mean and reduced beat-to-beat variability (lower standard deviation). Thus one can trivially assume that these larger jumps in healthy records are due only to the fact that their average standard deviation is larger \[Fig. 1(a)(b)\]. In order to systematically compare the statistical properties of the jumps between different individuals and different groups, we normalize each time series by subtracting the global average (over 6 hours) and dividing by the global standard deviation. In this way, all individual time series have zero mean and unit standard deviation \[Fig. 1(c)(d)\]. Such a normalization does not affect the results of our segmentation procedure.
We find that both the healthy subjects and the cosmonauts follow identical distributions, but the distribution of the jumps obtained from the heart failure group are markedly different — centered around lower values — indicating that, even after normalization, there is a higher probability for smaller jumps compared to the healthy subjects \[Fig. 5(a)\]. Note that the distributions for all groups appear to follow an identical homogeneous functional form, so we can collapse these distributions on top of each other by means of a homogeneous transformation \[Fig. 5(b)\]. The ratio between the scaling parameters used in this transformation gives us a factor by which this feature of the heartrate variability is reduced for the subjects with heart failure as compared to the healthy subjects. This finding indicates that, although the heartrate variability is reduced with disease, there may be a common structure to this variability, reflected in the identical functional form. These observations agree with previously reported results for the distribution of heartbeat fluctuations obtained by means of wavelet and Hilbert transforms .
In summary, we present a new method to probe the nonstationarity of a signal by partitioning it into segments with different mean values. We find a scale-invariant structure in the nonstationarity of a time series representative of a complex dynamics, namely the human heartbeat. This structure is characterized by a power-law distribution of the lengths of segments with a scaling exponent which does not change under certain pathological conditions and cannot be explained by the presence of correlations in the data. We find also a common structure to the jumps between consecutive segments, with a change in the scaling parameters with disease.
We thank Y. Ashkenazy, V. Schulte-Frohlinde, I. Grosse, S. Havlin, S. Mossa, C.-K. Peng, and Z. Struzik for helpful discussions and suggestions, grants BIO99-0651-CO2-01 (from the Spanish Government) and NIH/NCRR (P41RR13622), NASA, and the Mathers Charitable Foundation for support.
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# KUNS-1663hep-th/0005101 Noncommutative Monopole at the Second Order in 𝜃
## 1 Introduction
An explicit relation between the noncommutative fields and the commutative ones has been presented in , called “the Seiberg-Witten (SW) map.” The noncommutative Dirac-Born-Infeld (DBI) theory and the ordinary one appear as low-energy effective theories of D-brane in a constant NSNS $`B`$-field. They differ by the choice of the regularization for the worldsheet theory; the Pauli-Villars regularization for the commutative description and the point-splitting regularization for the noncommutative one. This means that these two descriptions are connected by some field redefinition and this is the SW map.
The relation between the commutative and noncommutative DBI theories has been examined in various aspects. Now in particular, let us concentrate on the BPS solutions and compare them in both the descriptions. The reason is that the BPS solutions are considered as powerful tools beyond the perturbative understanding. Noncommutative BPS monopoles describe, by the brane interpretation of , the configurations of tilted D-strings ending on parallel D3-branes in a constant NSNS $`B`$-field and have been investigated in various papers . In , noncommutative $`U(2)`$ monopole was considered at the first order in the noncommutativity parameter $`\theta ^{ij}`$. The analysis using the noncommutative eigenvalue equation for the scalar field successfully reproduced the tilted D-string picture. In , the similar analysis was carried out for the string junction and the anticipated result was obtained. Study of the noncommutative monopoles using the SW map was carried out in (see also ). There, the noncommutative BPS solutions were transformed into the commutative description via the SW map, and then the brane interpretation was done for the eigenvalues of the mapped scalar field to give the expected tilted D-string picture.
The purpose of this paper is to extend the analysis of the noncommutative $`U(2)`$ monopole using the SW map to second order in $`\theta `$. The motivating fact is as follows: the SW map possesses some ambiguities in higher orders in $`\theta `$ . This map is derived from the requirement of the gauge equivalence of the two descriptions. Since this is a very weak requirement, arbitrary parameters appear in the map. There are two types of ambiguities in it. One is in the form of the gauge transformation and has no physical effect. However, the other type of ambiguity consists of gauge covariant quantities and can cause physical differences.
We apply the SW map to the noncommutative monopole solution at the second order in $`\theta `$ and examine the effects of the ambiguities. Concretely, we compare the eigenvalues of the scalar field obtained by the SW map with that in the commutative Yang-Mills theory in a background magnetic field. Note that the ambiguities in the SW map can change the scalar eigenvalues (which are gauge invariant quantities) and hence change their brane interpretation. It is found that we can make these two eigenvalues coincide with each other by tuning the free parameters in the SW map. This gives an example of how the ambiguities in the SW map are fixed in concrete physical situations.
The rest of this paper is organized as follows. In section 2, we solve the noncommutative version of the BPS equation to second order in $`\theta `$. In section 3, we apply the SW map to the solution and evaluate the eigenvalues of the scalar field in the commutative description. In section 4, we compare the scalar eigenvalues of section 3 with those in the commutative Yang-Mills theory in a constant magnetic field, and examine the effect of the ambiguities in the SW map. In section 4, we summarize the paper and give some discussions. The SW map to second order in the change of $`\theta `$ is presented in Appendix A.
## 2 Noncommutative BPS monopole solution at $`\theta ^2`$
We shall consider the $`𝒩=4`$ $`U(2)`$ noncommutative super Yang-Mills theory in $`1+3`$ dimensions with the metric $`G_{\mu \nu }=\mathrm{diag}(1,1,1,1)`$, and construct the BPS monopole solution to second order in the noncommutativity parameter. The BPS equation reads
$$\widehat{D}_i\widehat{\mathrm{\Phi }}+\frac{1}{2}ϵ_{ijk}\widehat{F}_{jk}=0,$$
(2.1)
where the quantities with a hat denote those in the noncommutative description. In particular, we have
$`\widehat{F}_{ij}`$ $`_i\widehat{A}_j_j\widehat{A}_ii\widehat{A}_i\widehat{A}_j+i\widehat{A}_j\widehat{A}_i,`$ (2.2)
$`\widehat{D}_i\widehat{\mathrm{\Phi }}`$ $`_i\widehat{\mathrm{\Phi }}i\widehat{A}_i\widehat{\mathrm{\Phi }}+i\widehat{\mathrm{\Phi }}\widehat{A}_i,`$ (2.3)
where the $``$ product is defined by
$`(fg)(x)`$ $`f(x)\mathrm{exp}\left({\displaystyle \frac{i}{2}}\theta ^{ij}\stackrel{}{_i}\stackrel{}{_j}\right)g(x)`$
$`=f(x)g(x)+{\displaystyle \frac{i}{2}}\theta ^{ij}_if(x)_jg(x){\displaystyle \frac{1}{8}}\theta ^{ij}\theta ^{kl}_i_kf(x)_j_lg(x)+𝒪(\theta ^3).`$ (2.4)
In order to solve the BPS equation (2.1), we expand the fields in powers of $`\theta `$:
$`\widehat{\mathrm{\Phi }}`$ $`\left(\widehat{\mathrm{\Phi }}^{a(0)}+\widehat{\mathrm{\Phi }}^{a(1)}+\widehat{\mathrm{\Phi }}^{a(2)}\right){\displaystyle \frac{1}{2}}\sigma _a+\left(\widehat{\mathrm{\Phi }}^{0(1)}+\widehat{\mathrm{\Phi }}^{0(2)}\right){\displaystyle \frac{1}{2}}\text{ }\text{ ̵},`$
$`\widehat{A}_i`$ $`\left(\widehat{A}_i^{a(0)}+\widehat{A}_i^{a(1)}+\widehat{A}_i^{a(2)}\right){\displaystyle \frac{1}{2}}\sigma _a+\left(\widehat{A}_i^{0(1)}+\widehat{A}_i^{0(2)}\right){\displaystyle \frac{1}{2}}\text{ }\text{ ̵},`$ (2.5)
where the superscript $`(n)`$ denotes the order of $`\theta `$. As the solution at $`𝒪(\theta ^0)`$, we adopt the BPS monopole with vanishing $`U(1)`$ components $`\widehat{\mathrm{\Phi }}^{0(0)}=\widehat{A}_i^{0(0)}=0`$:
$$\widehat{\mathrm{\Phi }}^{a(0)}=\frac{\widehat{x}_a}{r}H(\xi ),\widehat{A}_i^{a(0)}=ϵ_{aij}\frac{\widehat{x}_j}{r}\left(1K(\xi )\right),$$
(2.6)
with $`\widehat{x}_ix_i/r`$ and $`\xi Cr`$ ($`C`$ is the parameter characterizing the mass of the monopole). The functions $`H`$ and $`K`$ are defined by
$$H(\xi )=\frac{\xi }{\mathrm{tanh}\xi }1,K(\xi )=\frac{\xi }{\mathrm{sinh}\xi }.$$
(2.7)
These two functions behave asymptotically as
$$H(\xi )=\xi 1+𝒪(e^\xi ),K(\xi )=0+𝒪(e^\xi ).$$
(2.8)
Note that the solution (2.6) has an invariance under the rotation by the diagonal subgroup $`SO(3)`$ of $`SO(3)_{\mathrm{space}}SO(3)_{\mathrm{gauge}}`$.
The solution at $`𝒪(\theta )`$ was constructed in and it is given by
$`\widehat{A}_i^{a(1)}`$ $`=0,`$ $`\widehat{A}_i^{0(1)}`$ $`=\theta ^{ij}\widehat{x}_j{\displaystyle \frac{1}{4r^3}}(1K)(1K+2H),`$
$`\widehat{\mathrm{\Phi }}^{a(1)}`$ $`=0,`$ $`\widehat{\mathrm{\Phi }}^{0(1)}`$ $`=0.`$ (2.9)
This solution is invariant under the generalized rotation, namely, the simultaneous rotation of the diagonal $`SO(3)`$ and the indices of the noncommutativity parameter $`\theta _{ij}`$. Note that the noncommutativity has no influence on the scalar solution at $`𝒪(\theta )`$.
Now let us consider the components at the second order in $`\theta `$ in the expansion (2.5). The $`𝒪(\theta ^2)`$ part of the BPS equation reads
$`_i\widehat{\mathrm{\Phi }}^{0(2)}+ϵ_{ijk}_j\widehat{A}_k^{0(2)}=0,`$ (2.10)
$`_i\widehat{\mathrm{\Phi }}^{a(2)}+ϵ_{ijk}_j\widehat{A}_k^{a(2)}+ϵ_{abc}\left(\widehat{A}_i^{b(0)}\widehat{\mathrm{\Phi }}^{c(2)}\widehat{\mathrm{\Phi }}^{b(0)}\widehat{A}_i^{c(2)}+ϵ_{ijk}\widehat{A}_j^{b(0)}\widehat{A}_k^{c(2)}\right)`$
$`={\displaystyle \frac{1}{2}}\theta ^{kl}_k\widehat{\mathrm{\Phi }}^{a(0)}_l\widehat{A}_i^{0(1)}{\displaystyle \frac{1}{2}}ϵ_{ijk}\theta ^{lm}_l\widehat{A}_j^{a(0)}_m\widehat{A}_k^{0(1)}`$
$`+{\displaystyle \frac{1}{8}}ϵ_{abc}\theta ^{lm}\theta ^{pq}\left(_l_p\widehat{A}_i^{b(0)}_m_q\widehat{\mathrm{\Phi }}^{c(0)}+{\displaystyle \frac{1}{2}}ϵ_{ijk}_l_p\widehat{A}_j^{b(0)}_m_q\widehat{A}_k^{c(0)}\right),`$ (2.11)
where the first equation (2.10) is the $`U(1)`$ part of (2.1), while the second equation (2.11) is the $`SU(2)`$ part. The $`U(1)`$ part has no regular solutions and we shall concentrate on the $`SU(2)`$ part (2.11).
In order to solve eq. (2.11), we adopt the generalized rotational invariance used in the construction of the $`𝒪(\theta )`$ part (2.9), and expand $`\widehat{\mathrm{\Phi }}^{a(2)}`$ and $`\widehat{A}_i^{a(2)}`$ as
$`\widehat{\mathrm{\Phi }}^{a(2)}`$ $`={\displaystyle \frac{1}{r^5}}\left[\varphi _1(\xi )(\theta \widehat{x})\theta _a+\varphi _2(\xi )\theta ^2\widehat{x}_a+\varphi _3(\xi )(\theta \widehat{x})^2\widehat{x}_a\right],`$
$`\widehat{A}_i^{a(2)}`$ $`={\displaystyle \frac{1}{r^5}}[a_1(\xi )\theta ^2ϵ_{aij}\widehat{x}_j+a_2(\xi )(\theta \widehat{x})ϵ_{aij}\theta _j+a_3(\xi )ϵ_{ajk}\theta _i\theta _j\widehat{x}_k`$
$`+a_4(\xi )(\theta \widehat{x})^2ϵ_{aij}\widehat{x}_j+a_5(\xi )(\theta \widehat{x})ϵ_{ajk}\widehat{x}_i\theta _j\widehat{x}_k],`$ (2.12)
where we have used $`\theta _i(1/2)ϵ_{ijk}\theta ^{jk}`$, $`\theta ^2\theta _i\theta _i`$ and $`(\theta \widehat{x})\theta _i\widehat{x}_i`$. One can check that this is the most general expression satisfying the generalized rotational invariance, by using the identities,
$$ϵ_{aij}x^2=\left(ϵ_{kij}x_a+ϵ_{akj}x_i+ϵ_{aik}x_j\right)x_k,$$
(2.13)
and the same one with $`x_i`$ replaced by $`\theta _i`$. Putting the expansion (2.12) into (2.11), we obtain the following linear differential equations with inhomogeneous terms for the unknown functions $`\varphi _k(\xi )`$ and $`a_k(\xi )`$:
$`𝒟(a_1a_3)+\varphi _2+4a_1a_2+5a_3+a_5(1K)\varphi _2Ha_1=I_1,`$
$`𝒟(\varphi _2+a_1+a_3)6\varphi _26a_16a_3a_5+(1K)(\varphi _2+2a_1+a_3)+Ha_1=I_2,`$
$`𝒟a_3+\varphi _1+a_25a_3a_5Ha_3=I_3,`$
$`𝒟(a_2a_3)+2\varphi _36a_2+6a_3+2a_5+(1K)(\varphi _1+a_2a_3a_5)+Ha_3=I_4,`$
$`𝒟(\varphi _1a_3)6\varphi _1+6a_3+2a_4+a_5+(1K)a_2+H(a_2a_5)=I_5,`$
$`𝒟(a_2+a_3a_4)+\varphi _3+6a_26a_3+4a_42a_5(1K)(\varphi _1+\varphi _3)H(a_2+a_4)=I_6,`$
$`𝒟(\varphi _3+a_4)8\varphi _38a_4+(1K)(\varphi _3+2a_4+a_5)+H(a_4+a_5)=I_7,`$ (2.14)
where the differential operator $`𝒟\xi (d/d\xi )`$ has been introduced. The seven equations (2.14) correspond to seven independent structures of eq. (2.11); $`\theta ^2\delta _{ai}`$, $`\theta ^2\widehat{x}_a\widehat{x}_i`$, $`\theta _a\theta _i`$, $`(\theta \widehat{x})\widehat{x}_a\theta _i`$, $`(\theta \widehat{x})\theta _a\widehat{x}_i`$, $`(\theta \widehat{x})^2\delta _{ai}`$ and $`(\theta \widehat{x})^2\widehat{x}_a\widehat{x}_i`$, respectively. The functions $`I_k`$ ($`k=1,\mathrm{},7`$) are the inhomogeneous terms which are polynomials of $`H`$ and $`K`$:
$`I_1(\xi )`$ $`={\displaystyle \frac{13}{8}}+{\displaystyle \frac{3H}{8}}{\displaystyle \frac{H^2}{2}}+{\displaystyle \frac{9K}{4}}+{\displaystyle \frac{5HK}{4}}+{\displaystyle \frac{3H^2K}{8}}+{\displaystyle \frac{H^3K}{8}}+{\displaystyle \frac{K^2}{8}}{\displaystyle \frac{5HK^2}{8}}`$
$`{\displaystyle \frac{H^2K^2}{4}}{\displaystyle \frac{5K^3}{8}}HK^3{\displaystyle \frac{3H^2K^3}{8}}{\displaystyle \frac{K^5}{8}},`$
$`I_2(\xi )`$ $`={\displaystyle \frac{3}{8}}+{\displaystyle \frac{7H}{8}}+{\displaystyle \frac{H^2}{2}}{\displaystyle \frac{9K}{8}}{\displaystyle \frac{3HK}{2}}{\displaystyle \frac{3H^2K}{8}}{\displaystyle \frac{H^3K}{8}}+K^2+{\displaystyle \frac{HK^2}{8}}`$
$`{\displaystyle \frac{H^2K^2}{2}}{\displaystyle \frac{H^3K^2}{4}}+{\displaystyle \frac{3HK^3}{4}}+{\displaystyle \frac{3H^2K^3}{8}}{\displaystyle \frac{3K^4}{8}}{\displaystyle \frac{HK^4}{4}}+{\displaystyle \frac{K^5}{8}},`$
$`I_3(\xi )`$ $`={\displaystyle \frac{5}{8}}+{\displaystyle \frac{9H}{8}}+{\displaystyle \frac{H^2}{2}}{\displaystyle \frac{17K}{8}}3HK{\displaystyle \frac{5H^2K}{4}}{\displaystyle \frac{H^3K}{4}}+{\displaystyle \frac{21K^2}{8}}+{\displaystyle \frac{21HK^2}{8}}`$
$`+{\displaystyle \frac{3H^2K^2}{4}}{\displaystyle \frac{11K^3}{8}}{\displaystyle \frac{3HK^3}{4}}+{\displaystyle \frac{K^4}{4}},`$
$`I_4(\xi )`$ $`={\displaystyle \frac{11}{8}}{\displaystyle \frac{17H}{8}}{\displaystyle \frac{H^2}{2}}+{\displaystyle \frac{17K}{4}}+5HK+{\displaystyle \frac{3H^2K}{2}}+{\displaystyle \frac{H^3K}{4}}{\displaystyle \frac{9K^2}{2}}{\displaystyle \frac{29HK^2}{8}}`$
$`H^2K^2+{\displaystyle \frac{7K^3}{4}}+{\displaystyle \frac{3HK^3}{4}}{\displaystyle \frac{K^4}{8}},`$
$`I_5(\xi )`$ $`=1{\displaystyle \frac{3H}{2}}{\displaystyle \frac{3H^2}{4}}+{\displaystyle \frac{13K}{4}}+{\displaystyle \frac{31HK}{8}}+{\displaystyle \frac{3H^2K}{2}}+{\displaystyle \frac{H^3K}{4}}{\displaystyle \frac{15K^2}{4}}{\displaystyle \frac{13HK^2}{4}}`$
$`{\displaystyle \frac{3H^2K^2}{4}}+{\displaystyle \frac{7K^3}{4}}+{\displaystyle \frac{7HK^3}{8}}{\displaystyle \frac{K^4}{4}},`$
$`I_6(\xi )`$ $`=2H+{\displaystyle \frac{3H^2}{4}}{\displaystyle \frac{19K}{8}}{\displaystyle \frac{3HK}{4}}{\displaystyle \frac{3H^2K}{8}}{\displaystyle \frac{H^3K}{8}}K^2+{\displaystyle \frac{HK^2}{2}}`$
$`+{\displaystyle \frac{H^2K^2}{4}}+{\displaystyle \frac{5K^3}{4}}+{\displaystyle \frac{5HK^3}{4}}+{\displaystyle \frac{3H^2K^3}{8}}+{\displaystyle \frac{K^5}{8}},`$
$`I_7(\xi )`$ $`={\displaystyle \frac{7}{4}}+{\displaystyle \frac{7H}{4}}{\displaystyle \frac{43K}{8}}{\displaystyle \frac{39HK}{8}}{\displaystyle \frac{11H^2K}{8}}{\displaystyle \frac{H^3K}{8}}+{\displaystyle \frac{23K^2}{4}}+{\displaystyle \frac{19HK^2}{4}}`$
$`+{\displaystyle \frac{7H^2K^2}{4}}+{\displaystyle \frac{H^3K^2}{4}}{\displaystyle \frac{5K^3}{2}}{\displaystyle \frac{15HK^3}{8}}{\displaystyle \frac{3H^2K^3}{8}}+{\displaystyle \frac{K^4}{2}}+{\displaystyle \frac{HK^4}{4}}{\displaystyle \frac{K^5}{8}}.`$ (2.15)
We can solve the differential equations (2.14) by the same polynomial assumption as used in the construction of the noncommutative $`1/4`$ BPS solution , that is, we assume that the functions $`=\varphi _k`$ and $`a_k`$ are given as polynomials of $`H`$ and $`K`$:
$$=\underset{n=0}{\overset{n_{\mathrm{max}}}{}}\underset{m=0}{\overset{m_{\mathrm{max}}}{}}_{nm}H^nK^m,$$
(2.16)
with suitably large $`n_{\mathrm{max}}`$ and $`m_{\mathrm{max}}`$. This assumption is owing to the property of $`H`$ and $`K`$,
$`𝒟K`$ $`=HK,`$
$`𝒟H`$ $`=1+HK^2,`$ (2.17)
which implies that the operation of $`𝒟`$ on a polynomial of $`H`$ and $`K`$ just reproduces another polynomial of them. With the assumption (2.16), the differential equations (2.14) are reduced to a set of linear algebraic equations for the coefficients $`_{nm}`$, which can be solved straightforwardly. Note that we originally had seven differential equations (2.14) for eight unknown functions, so the solution contains one undetermined function. This is the gauge freedom which preserves the generalized rotational invariant form (2.12), and the corresponding gauge transformation function is
$$\lambda ^a=ϵ_{aij}\theta _ix_j(\theta \widehat{x})\frac{1}{r^4}\lambda (\xi ).$$
(2.18)
Using this freedom to choose $`a_5(\xi )=0`$, the solution to the the $`𝒪(\theta ^2)`$ part (2.11) of the noncommutative BPS equation are given as follows:
$`\varphi _1(\xi )`$ $`={\displaystyle \frac{1}{4}}H+{\displaystyle \frac{1}{4}}H^2{\displaystyle \frac{1}{8}}H^3+{\displaystyle \frac{1}{4}}HK^2,`$
$`\varphi _2(\xi )`$ $`={\displaystyle \frac{1}{8}}{\displaystyle \frac{3}{8}}H+{\displaystyle \frac{1}{8}}H^2{\displaystyle \frac{1}{4}}K^2+{\displaystyle \frac{3}{8}}HK^2+{\displaystyle \frac{1}{8}}H^2K^2+{\displaystyle \frac{1}{8}}K^4,`$
$`\varphi _3(\xi )`$ $`={\displaystyle \frac{1}{8}}+{\displaystyle \frac{7}{8}}H{\displaystyle \frac{5}{8}}H^2+{\displaystyle \frac{1}{8}}H^3+{\displaystyle \frac{1}{4}}K^2{\displaystyle \frac{7}{8}}HK^2{\displaystyle \frac{1}{8}}H^2K^2{\displaystyle \frac{1}{8}}K^4.`$
$`a_1(\xi )`$ $`={\displaystyle \frac{1}{8}}+{\displaystyle \frac{1}{2}}H{\displaystyle \frac{1}{8}}K{\displaystyle \frac{1}{2}}HK{\displaystyle \frac{1}{4}}H^2K+{\displaystyle \frac{5}{8}}K^2+{\displaystyle \frac{1}{4}}HK^2{\displaystyle \frac{3}{8}}K^3{\displaystyle \frac{1}{4}}HK^3,`$
$`a_2(\xi )`$ $`={\displaystyle \frac{1}{4}}+{\displaystyle \frac{1}{2}}H{\displaystyle \frac{3}{8}}H^2{\displaystyle \frac{3}{4}}K{\displaystyle \frac{1}{2}}HK+{\displaystyle \frac{1}{8}}H^2K+{\displaystyle \frac{3}{4}}K^2{\displaystyle \frac{1}{4}}K^3,`$
$`a_3(\xi )`$ $`={\displaystyle \frac{1}{8}}{\displaystyle \frac{1}{4}}H{\displaystyle \frac{1}{8}}H^2+{\displaystyle \frac{3}{8}}K+{\displaystyle \frac{1}{2}}HK+{\displaystyle \frac{1}{8}}H^2K{\displaystyle \frac{3}{8}}K^2{\displaystyle \frac{1}{4}}HK^2+{\displaystyle \frac{1}{8}}K^3,`$
$`a_4(\xi )`$ $`={\displaystyle \frac{1}{4}}{\displaystyle \frac{3}{2}}H+{\displaystyle \frac{1}{2}}H^2+{\displaystyle \frac{5}{4}}K+{\displaystyle \frac{3}{2}}HK+{\displaystyle \frac{1}{4}}H^2K{\displaystyle \frac{7}{4}}K^2{\displaystyle \frac{1}{4}}HK^2+{\displaystyle \frac{3}{4}}K^3+{\displaystyle \frac{1}{4}}HK^3.`$ (2.19)
## 3 SW map and the eigenvalues of the scalar field
Having obtained the classical solution to the noncommutative BPS equation to $`𝒪(\theta ^2)`$, our next task is to transform it into the commutative description via the SW map to get the eigenvalues of the scalar field . For this purpose, we first have to establish the SW map to second order in the change $`\delta \theta `$ of the noncommutativity parameter.
It was pointed out in that the SW map has inherent ambiguities. There are two types of ambiguities in it. One is of the form identifiable as gauge transformations. The other type of ambiguity consists of gauge covariant quantities. The latter can cause physical differences and must be fixed by some physical requirements. This type of ambiguities comes from the path dependence of the map in the $`\theta `$-space. In other words, even if we perform the map to go round in the $`\theta `$-space, we do not come back to the original configuration. This means that the SW map at $`𝒪(\delta \theta ^2)`$ and higher has such a type of ambiguities.
The SW map for the gauge field to $`𝒪(\delta \theta ^2)`$, including the ambiguities, is presented in Appendix A. Here, we need the SW map for the scalar field from a noncommutative space with small $`\theta `$ to the commutative space. This map is obtained by performing the dimensional reduction of the map for the gauge field and taking
$$\delta \theta ^{ij}=\theta ^{ij}.$$
(3.1)
Then, the scalar field $`\mathrm{\Phi }`$ in the commutative description is expressed in terms of $`\widehat{\mathrm{\Phi }}`$ and $`\widehat{A}_i`$ in the noncommutative description as
$$\mathrm{\Phi }=\widehat{\mathrm{\Phi }}+\mathrm{\Delta }\widehat{\mathrm{\Phi }}^{(1)}+\mathrm{\Delta }\widehat{\mathrm{\Phi }}^{(2)},$$
(3.2)
with $`\mathrm{\Delta }\widehat{\mathrm{\Phi }}^{(1)}`$ and $`\mathrm{\Delta }\widehat{\mathrm{\Phi }}^{(2)}`$ given by
$`\mathrm{\Delta }\widehat{\mathrm{\Phi }}^{(1)}`$ $`={\displaystyle \frac{1}{4}}\delta \theta ^{kl}\{A_k,(_l+D_l)\mathrm{\Phi }\}i\alpha \delta \theta ^{kl}[\mathrm{\Phi },F_{kl}]2\beta \delta \theta ^{kl}[\mathrm{\Phi },A_kA_l],`$ (3.3)
$`\mathrm{\Delta }\widehat{\mathrm{\Phi }}^{(2)}`$ $`={\displaystyle \frac{1}{4}}\text{Im}(\begin{array}{c}\text{ }\text{ }\hfill \\ A\mathrm{\Phi }\hfill \end{array})+{\displaystyle \frac{1}{4}}\text{Re}(\begin{array}{c}\text{ }\text{ }\hfill \\ AA\mathrm{\Phi }\hfill \end{array}){\displaystyle \frac{1}{4}}\text{Re}(\begin{array}{c}\text{ }\text{ }\hfill \\ AA\mathrm{\Phi }\hfill \end{array})`$ (3.10)
14Re( ADΦA)+14Re( AADΦ)+14Re(
AF
D
Φ)14Refragments 𝐴𝐷Φ𝐴14Refragments 𝐴𝐴𝐷Φ14Re
AF
D
Φ\displaystyle\quad-\frac{1}{4}\mathop{\mbox{Re}}(\mathop{\vbox{\halign{#\cr\kern 1.72218pt\cr$\hbox{$\hskip 3.75pt\vrule height=3.65973pt,width=0.29999pt,depth=-0.0004pt\kern-0.29999pt\vrule height=3.65973pt,width=8.90446pt,depth=-3.35974pt\kern-0.29999pt\vrule height=3.65973pt,width=0.29999pt,depth=-0.11153pt$\hss}\hbox{$\hskip 18.75415pt\vrule height=3.65973pt,width=0.29999pt,depth=-0.0004pt\kern-0.29999pt\vrule height=3.65973pt,width=16.91733pt,depth=-3.35974pt\kern-0.29999pt\vrule height=3.65973pt,width=0.29999pt,depth=-0.0004pt$\hss}$\crcr\kern 1.72218pt\nointerlineskip\cr\hbox{$\displaystyle{}{A}{\,}{\partial}{}{D}{\Phi\,}{A}$}\crcr}}}\limits)+\frac{1}{4}\mathop{\mbox{Re}}(\mathop{\vbox{\halign{#\cr\kern 1.72218pt\cr$\hbox{$\hskip 3.75pt\vrule height=3.65973pt,width=0.29999pt,depth=-0.0004pt\kern-0.29999pt\vrule height=3.65973pt,width=18.0711pt,depth=-3.35974pt\kern-0.29999pt\vrule height=3.65973pt,width=0.29999pt,depth=-0.11153pt$\hss}\hbox{$\hskip 12.91664pt\vrule height=7.31946pt,width=0.29999pt,depth=-0.0004pt\kern-0.29999pt\vrule height=7.31946pt,width=15.00415pt,depth=-7.01947pt\kern-0.29999pt\vrule height=7.31946pt,width=0.29999pt,depth=-0.0004pt$\hss}$\crcr\kern 1.72218pt\nointerlineskip\cr\hbox{$\displaystyle{}{A}{\,}{A}{\,}{\partial}{}{D}\Phi$}\crcr}}}\limits)+\frac{1}{4}\mathop{\mbox{Re}}(\vbox{\halign{#\cr\hskip 4.0pt\vrule depth=4.0pt\hrulefill\vrule depth=4.0pt\hskip 5.0pt\crcr\kern-1.0pt\vskip 3.55658pt\nointerlineskip\cr$\hfil\displaystyle{AF}\hfil$\crcr}}\vbox{\halign{#\cr\hskip-1.0pt\vrule depth=4.0pt\hrulefill\vrule depth=4.0pt\hskip 3.0pt\crcr\kern-1.0pt\vskip 3.55658pt\nointerlineskip\cr$\hfil\displaystyle{D}\hfil$\crcr}}\Phi) (3.19)
$`+{\displaystyle \frac{1}{4}}\text{Re}(\begin{array}{c}A\hfill \end{array}\begin{array}{c}D\hfill \end{array}\mathrm{\Phi }F){\displaystyle \frac{1}{8}}\text{Re}(\begin{array}{c}\text{ }\text{ }\hfill \\ AA\mathrm{\Phi }\hfill \end{array})+{\displaystyle \frac{1}{8}}\text{Re}(\begin{array}{c}\text{ }\text{ }\hfill \\ A\mathrm{\Phi }A\hfill \end{array})`$ (3.28)
$`{\displaystyle \frac{1}{8}}\text{Re}(\begin{array}{c}\text{ }\text{ }\hfill \\ AAAA\mathrm{\Phi }\hfill \end{array})+{\displaystyle \frac{1}{4}}\text{Re}(\begin{array}{c}\text{ }\text{ }\hfill \\ AAA\mathrm{\Phi }A\hfill \end{array}){\displaystyle \frac{1}{8}}\text{Re}(\begin{array}{c}\text{ }\text{ }\hfill \\ AA\mathrm{\Phi }AA\hfill \end{array})`$ (3.35)
$`\left({\displaystyle \frac{1}{16}}+8\alpha \beta +4\beta ^2\right)\text{Re}(\begin{array}{c}\text{ }\text{ }\hfill \\ AA[\mathrm{\Phi },AA]\hfill \end{array})+8\alpha \beta \text{Im}(\begin{array}{c}\text{ }\text{ }\hfill \\ A[\mathrm{\Phi },AA]\hfill \end{array})`$ (3.40)
$`\gamma _1\text{Re}(\begin{array}{c}F\hfill \end{array}[\mathrm{\Phi },\begin{array}{c}F\hfill \end{array}])\gamma _2\text{Re}(\begin{array}{c}F\hfill \end{array}[\mathrm{\Phi },\begin{array}{c}F\hfill \end{array}])+\gamma _3\text{Im}(\begin{array}{c}\text{ }\hfill \\ D\mathrm{\Phi }D\hfill \end{array}\begin{array}{c}F\hfill \end{array})`$ (3.53)
$`+\mathrm{\Delta }\widehat{\mathrm{\Phi }}_{\mathrm{metric}}^{(2)}+(\text{gauge-type ambiguities}).`$ (3.54)
In (3.3) and (3.54), all the fields are defined at $`\theta `$ and all the products are the $``$ products (we have omitted the hats on the fields and the $``$ for the products). We have used the following simplified notation:
ABδθklAkBl, AFiδθklAkFli,
FδθklFkl,formulae-sequence AB𝛿superscript𝜃𝑘𝑙subscript𝐴𝑘subscript𝐵𝑙formulae-sequence AFi𝛿superscript𝜃𝑘𝑙subscript𝐴𝑘subscript𝐹𝑙𝑖
F𝛿superscript𝜃𝑘𝑙subscript𝐹𝑘𝑙\mathop{\vbox{\halign{#\cr\kern 1.72218pt\cr$\hbox{$\hskip 3.75pt\vrule height=3.65973pt,width=0.29999pt,depth=-0.0004pt\kern-0.29999pt\vrule height=3.65973pt,width=7.79341pt,depth=-3.35974pt\kern-0.29999pt\vrule height=3.65973pt,width=0.29999pt,depth=-0.0004pt$\hss}$\crcr\kern 1.72218pt\nointerlineskip\cr\hbox{$\displaystyle{}{A}{}{B}$}\crcr}}}\limits\equiv\delta\theta^{kl}A_{k}B_{l},\quad\mathop{\vbox{\halign{#\cr\kern 1.72218pt\cr$\hbox{$\hskip 3.75pt\vrule height=3.65973pt,width=0.29999pt,depth=-0.0004pt\kern-0.29999pt\vrule height=3.65973pt,width=7.65974pt,depth=-3.35974pt\kern-0.29999pt\vrule height=3.65973pt,width=0.29999pt,depth=-0.0004pt$\hss}$\crcr\kern 1.72218pt\nointerlineskip\cr\hbox{$\displaystyle{}{A}{}{F}_{i}$}\crcr}}}\limits\equiv\delta\theta^{kl}A_{k}F_{li},\quad\vbox{\halign{#\cr\hskip 2.0pt\vrule depth=4.0pt\hrulefill\vrule depth=4.0pt\hskip 2.0pt\crcr\kern-1.0pt\vskip 3.55658pt\nointerlineskip\cr$\hfil\displaystyle{F}\hfil$\crcr}}\equiv\delta\theta^{kl}F_{kl}, (3.55)
and
$$\text{Re}𝒪\frac{1}{2}\left(𝒪+𝒪^{}\right),\text{Im}𝒪\frac{1}{2i}\left(𝒪𝒪^{}\right).$$
(3.56)
Note that the contraction symbol has the property that $`(\begin{array}{c}\text{ }\hfill \\ AB\hfill \end{array})^{}=\begin{array}{c}\text{ }\hfill \\ B{}_{}{}^{}A_{}^{}\hfill \end{array}`$ due to the anti-symmetry of $`\delta \theta ^{kl}`$. We shall mention the $`\mathrm{\Delta }\widehat{\mathrm{\Phi }}_{\mathrm{metric}}^{(2)}`$ term in (3.54) soon below.
The SW map at $`𝒪(\delta \theta )`$ (3.3) contains two ambiguities parameterized by $`\alpha `$ and $`\beta `$. They are the gauge-type ambiguities. On the other hand, there exist three covariant-type ambiguities with coefficients $`\gamma _k`$ ($`k=1,2,3`$) in the SW map at $`𝒪(\delta \theta ^2)`$ (3.54). This type of ambiguities directly affects the eigenvalues of the scalar. Note that the gauge-type ambiguities at $`𝒪(\delta \theta )`$ have influence on the map at $`𝒪(\delta \theta ^2)`$ and may possibly change the eigenvalues.
The term $`\mathrm{\Delta }\widehat{\mathrm{\Phi }}_{\mathrm{metric}}^{(2)}`$ in (3.54) represents the covariant-type ambiguities using the metric.<sup>*</sup><sup>*</sup>* SW map containing the metric was also considered in in a different context. Note that, in all the terms written explicitly in (3.54), the upper indices of $`\theta ^{ij}`$ are contracted with the lower ones of $`A_i`$ and $`_i`$ without using the metric. On the other hand, the terms in $`\mathrm{\Delta }\widehat{\mathrm{\Phi }}_{\mathrm{metric}}^{(2)}`$ are constructed by using the metric. There are many terms belonging to $`\mathrm{\Delta }\widehat{\mathrm{\Phi }}_{\mathrm{metric}}^{(2)}`$. Examples are
$$\delta ^{km}\delta ^{ln}\text{Re}(F_{kl}[\mathrm{\Phi },F_{mn}])\theta ^2,\delta ^{km}\delta ^{ln}\text{Re}(D_k\mathrm{\Phi }[\mathrm{\Phi },D_l\mathrm{\Phi }])\theta _m\theta _n.$$
(3.57)
These are obtained by the dimensional reduction of the corresponding operators in the SW map for the gauge field given in eq. (A.91) of Appendix A.
Now we shall proceed to the evaluation of the eigenvalues of the scalar $`\mathrm{\Phi }`$ (3.2) as a $`2\times 2`$ matrix to $`𝒪(\theta ^2)`$. For this purpose, we expand $`\mathrm{\Phi }`$ in the commutative description in powers of $`\theta `$:
$$\mathrm{\Phi }=\mathrm{\Phi }^{(0)}+\mathrm{\Phi }^{(1)}+\mathrm{\Phi }^{(2)},$$
(3.58)
where $`\mathrm{\Phi }^{(n)}`$ is of order $`\theta ^n`$. Let us write explicitly the arguments of the SW map, (3.3) and (3.54), as $`\mathrm{\Delta }\widehat{\mathrm{\Phi }}^{(n)}[\widehat{\mathrm{\Phi }},\widehat{A}_i,\theta ]`$ with the last argument representing the $`\theta `$-dependence only through the $``$ product (2.4). Then, using the noncommutative classical solution (2.5), we have
$`\mathrm{\Phi }^{(0)}`$ $`=\widehat{\mathrm{\Phi }}^{a(0)}{\displaystyle \frac{1}{2}}\sigma _a,`$ (3.59)
$`\mathrm{\Phi }^{(1)}`$ $`=\mathrm{\Delta }\widehat{\mathrm{\Phi }}^{(1)}[\widehat{\mathrm{\Phi }}^{(0)},\widehat{A}_i^{(0)},\theta =0],`$ (3.60)
$`\mathrm{\Phi }^{(2)}`$ $`=\widehat{\mathrm{\Phi }}^{a(2)}{\displaystyle \frac{1}{2}}\sigma _a+\mathrm{\Delta }\widehat{\mathrm{\Phi }}^{(1)}[\widehat{\mathrm{\Phi }},\widehat{A}_i,\theta ]|_{\theta ^2}+\mathrm{\Delta }\widehat{\mathrm{\Phi }}^{(2)}[\widehat{\mathrm{\Phi }}^{(0)},\widehat{A}_i^{(0)},\theta =0].`$ (3.61)
We shall add some explanations about (3.59) – (3.61). First, we have set $`\theta =0`$ in (3.60) and the last term of (3.61). This implies that we take the commutative products among the fields. Next, the second term on the RHS of (3.61) means the sum of all the terms quadratic in $`\theta `$ in $`\mathrm{\Delta }\widehat{\mathrm{\Phi }}^{(1)}[\widehat{\mathrm{\Phi }},\widehat{A}_i,\theta ]`$. There are three sources of $`\theta `$; $`\delta \theta ^{kl}=\theta ^{kl}`$ in (3.3), $`\theta `$ in the $``$ product, and $`\theta `$ in the noncommutative classical solution (2.9).
Then, the two eigenvalues $`\lambda _\pm `$ of the scalar $`\mathrm{\Phi }`$ (3.58) are given using the well-known perturbation theory formula as
$$\lambda _\pm =\lambda _\pm ^{(0)}+\lambda _\pm ^{(1)}+\lambda _\pm ^{(2)},$$
(3.62)
with
$`\lambda _\pm ^{(0)}`$ $`=\pm {\displaystyle \frac{H(\xi )}{2r}},`$ (3.63)
$`\lambda _\pm ^{(1)}`$ $`=\pm \left|\mathrm{\Phi }^{(1)}\right|\pm ,`$ (3.64)
$`\lambda _\pm ^{(2)}`$ $`=\pm \left|\mathrm{\Phi }^{(2)}\right|\pm +{\displaystyle \frac{\pm \left|\mathrm{\Phi }^{(1)}\right|\left|\mathrm{\Phi }^{(1)}\right|\pm }{\lambda _\pm ^{(0)}\lambda _{}^{(0)}}},`$ (3.65)
where the kets $`|\pm `$ are the eigenvectors of $`\mathrm{\Phi }^{(0)}`$ satisfying
$$\widehat{x}\sigma |\pm =\pm |\pm .$$
(3.66)
Plugging the noncommutative classical solution obtained in section 2 into (3.64) and (3.65), we get after a tedious but straightforward calculation
$`\lambda _+^{(1)}`$ $`={\displaystyle \frac{1}{4r^3}}H(1K^2)(\theta \widehat{x}),`$ (3.67)
$`\lambda _+^{(2)}`$ $`={\displaystyle \frac{1}{16r^5}}\left(H^2HK^2+{\displaystyle \frac{1}{2}}H^3K^2+HK^4\right)\theta ^2`$
$`+{\displaystyle \frac{1}{16r^5}}\left(2H3H^2HK^2{\displaystyle \frac{1}{2}}H^3K^2HK^4\right)(\theta \widehat{x})^2`$
$`+c_1f_1(x,\theta )+c_2f_2(x,\theta )++\left|\mathrm{\Delta }\widehat{\mathrm{\Phi }}^{(2)}[\widehat{\mathrm{\Phi }}^{(0)},\widehat{A}_i^{(0)},\theta =0]_{\mathrm{metric}}\right|+,`$ (3.68)
with
$$c_1=\frac{1}{2}\gamma _1+\gamma _2\frac{1}{2}\gamma _3+\alpha ^2,c_2=\frac{1}{2}\gamma _3,$$
(3.69)
and
$`f_1(x,\theta )`$ $`={\displaystyle \frac{1}{r^5}}H^3K^2\left(\theta ^2(\theta \widehat{x})^2\right),`$
$`f_2(x,\theta )`$ $`={\displaystyle \frac{1}{r^5}}\left(HK^2H^2K^2HK^4\right)\left(\theta ^23(\theta \widehat{x})^2\right).`$ (3.70)
The other eigenvalue $`\lambda _{}`$ is given by $`\lambda _{}^{(1)}=\lambda _+^{(1)}`$ and $`\lambda _{}^{(2)}=\lambda _+^{(2)}`$. The first order eigenvalue (3.67) is already obtained in . The origin of the $`\alpha ^2`$ term in $`c_1`$ (3.69) is the last term of (3.65). There are no other contributions to $`\lambda ^{(2)}`$ from the last term of (3.65) since we have $`\mathrm{\Phi }^{(1)}|_{\alpha =\beta =0}\text{ }\text{ }\text{ }\text{ }\text{ ̵}`$. The terms in the SW map (3.54) quadratic in $`\alpha `$ and $`\beta `$ do not contribute to $`\lambda ^{(2)}`$ owing to the property $`[\mathrm{\Phi }^{(0)},\begin{array}{c}\text{ }\hfill \\ A{}_{}{}^{(0)}A_{}^{(0)}\hfill \end{array}]=0`$ for the zero-th order solution. All the constituents of $`\lambda _+^{(2)}`$ (3.68), which are polynomials of $`H`$ and $`K`$ divided by $`r^5`$, vanish at the origin $`r=0`$.
## 4 Tilted D-string picture
We would like to compare the the scalar eigenvalues obtained in the previous section with those which are obtained by different ways and are expected to describe the same physical situation of the tilted D-string between two parallel D3-branes. In the $`U(1)`$ case, there are three ways giving the same result ; the SW map of the noncommutative BPS solution, the nonlinear BPS solution in the commutative space, and the target space rotation of the linear BPS solution in the commutative space. In particular, the linear BPS solution (under a constant magnetic field) gives the tilted D3-brane picture, which is related to the tilted D-string picture by the target space rotation (see figure 1). In the nonabelian case, the nonlinearly realized supertransformation of the DBI theory is not well-understood. Therefore, we shall take the target space rotation of the linear BPS solution in the commutative space as the object to be compared with the eigenvalues of section 3.
Let us consider the $`U(2)`$ super Yang-Mills theory in the commutative space with a constant $`U(1)`$ magnetic field $`B_i`$. The BPS equation of this system, which we regard as describing the tilted D3-brane in a constant NSNS $`B`$-field $`B_{ij}=ϵ_{ijk}B_k`$, is
$$D_i\mathrm{\Phi }+\frac{1}{2}ϵ_{ijk}\left(F_{jk}+B_{jk}\frac{1}{2}\text{ }\text{ }\text{ }\text{ }\text{ ̵}\right)=0.$$
(4.1)
For our present purpose of comparing with the previous section, we should in fact consider the Yang-Mills theory in the commutative space with the metric $`g_{ij}`$ related to the metric $`G_{ij}=\delta _{ij}`$ in the noncommutative theory of section 2 by
$$G^{ij}+\frac{\theta ^{ij}}{2\pi \alpha ^{}}=\left(\frac{1}{g+2\pi \alpha ^{}B}\right)^{ij}.$$
(4.2)
However, the desired scalar eigenvalue is obtained by considering the BPS equation (4.1) with $`g_{ij}=\delta _{ij}`$ and coordinate transforming back to the original $`g_{ij}`$ afterwards .
The $`U(1)`$ part of this equation
$$_i\mathrm{\Phi }^0+\frac{1}{2}ϵ_{ijk}B_{jk}=0,$$
(4.3)
is easily solved to give
$$\mathrm{\Phi }^0=\frac{1}{2}ϵ_{ijk}B_{jk}x^i=\frac{1}{(2\pi \alpha ^{})^2}(\theta x),$$
(4.4)
where the relation $`2\pi \alpha ^{}B_i=\theta _i/2\pi \alpha ^{}+𝒪(\theta ^3)`$ has been used. As a solution to the nonabelian part, we adopt the ordinary BPS monopole solution (2.6). We shall attach tilde to the space coordinates in the present system for distinguishing them from those in the rotated system to be discussed below. Then, the (larger) eigenvalue of the scalar field is
$$\stackrel{~}{\mathrm{\Lambda }}=\frac{1}{2\stackrel{~}{r}}H(\stackrel{~}{\xi })+\frac{1}{(2\pi \alpha ^{})^2}(\theta \stackrel{~}{x}),$$
(4.5)
with $`\stackrel{~}{\xi }\stackrel{~}{C}\stackrel{~}{r}`$ and $`\stackrel{~}{r}\stackrel{~}{x}_i\stackrel{~}{x}_i`$ ($`\stackrel{~}{C}`$ is the mass scale of the present monopole).
Now let us carry out the target space rotation and turn to the tilted D-string picture (see figure 2):
$$\left(\begin{array}{c}2\pi \alpha ^{}\mathrm{\Lambda }\\ (\widehat{\theta }x)\end{array}\right)=\left(\begin{array}{cc}\mathrm{cos}\varphi & \mathrm{sin}\varphi \\ \mathrm{sin}\varphi & \mathrm{cos}\varphi \end{array}\right)\left(\begin{array}{c}2\pi \alpha ^{}\stackrel{~}{\mathrm{\Lambda }}\\ (\widehat{\theta }\stackrel{~}{x})\end{array}\right),$$
(4.6)
where $`\widehat{\theta }_i`$ is the unit vector $`\widehat{\theta }_i\theta _i/|\theta |`$, and the rotation angle $`\varphi `$ is given as $`\mathrm{tan}\varphi =|\theta |/2\pi \alpha ^{}`$. The components perpendicular to $`\theta _i`$ are common between $`x_i`$ and $`\stackrel{~}{x}_i`$. Expressing the new eigenvalue $`\mathrm{\Lambda }`$ in the tilted D-string picture as a function of the coordinate $`x^i`$, we have
$`\mathrm{\Lambda }`$ $`={\displaystyle \frac{H}{2r}}{\displaystyle \frac{1}{4r^3}}H(1K^2)(\theta \widehat{x})+{\displaystyle \frac{1}{16r^5}}\left(H^2H^2K^2\right)\theta ^2`$
$`+{\displaystyle \frac{1}{16r^5}}\left(2H3H^24HK^2+3H^2K^2+2H^3K^2+2HK^4\right)(\theta \widehat{x})^2`$
$`{\displaystyle \frac{1}{(2\pi \alpha ^{})^2}}{\displaystyle \frac{1}{4r}}\left\{H\theta ^2+(1K^2)(\theta \widehat{x})^2\right\},`$ (4.7)
where the arguments of $`H`$ and $`K`$ are $`\stackrel{~}{C}r`$ with $`r^2=x^ix^i`$. Now we make the coordinate transformation in (4.7) from the metric $`\delta _{ij}`$ to $`g_{ij}=\delta _{ij}(\theta ^2\delta _{ij}\theta _i\theta _j)/(2\pi \alpha ^{})^2`$ corresponding to the open string metric $`G_{ij}=\delta _{ij}`$ adopted in section 2. This is accomplished by the replacement of $`r`$ with
$$\left(g_{ij}x^ix^j\right)^{1/2}=\left(1\frac{1}{(2\pi \alpha ^{})^2}\frac{1}{2}\left[\theta ^2(\theta \widehat{x})^2\right]\right)r.$$
(4.8)
Then, the eigenvalue $`\mathrm{\Lambda }`$ in the new coordinate system is given by (4.7) with the last $`1/(2\pi \alpha ^{})^2`$ term omitted and the arguments of $`H`$ and $`K`$ replaced by $`Cr`$ (cf. for the $`U(1)`$ case). Here, $`C`$ is the the D3-brane separation, $`C\stackrel{~}{C}\mathrm{cos}\varphi `$ (see figure 2). We can show that this eigenvalue is exactly the same as $`\mathrm{\Lambda }(x)`$ obtained by solving
$$\mathrm{\Lambda }(x)=\frac{1}{2r}H(|x_i\mathrm{\Lambda }(x)\theta _i|),$$
(4.9)
which implies the tilted D-string picture in figure 2(B). Namely, for a given value of $`\mathrm{\Lambda }`$, the corresponding $`x_i`$ lies on a sphere with its center at $`x_i=\mathrm{\Lambda }\theta _i`$ (cf. ).
Having finished the preparation of obtaining the eigenvalue from the target space rotation of the linear BPS solution, let us proceed to the comparison between this eigenvalue $`\mathrm{\Lambda }`$ (4.7) (without the $`1/(2\pi \alpha ^{})^2`$ term) and the eigenvalue (3.67) and (3.68) obtained from the noncommutative monopole via the SW map. First, the $`𝒪(\theta )`$ terms agree between them as was already shown in . Second, the $`𝒪(\theta ^2)`$ parts coincide perfectly in the asymptotic region $`r\mathrm{}`$ where we can drop the exponentially decaying terms (see eq. (2.8)). Note in particular that all the ambiguity terms in (3.68) disappear in the the asymptotic region. (The last term of (3.68) using the metric is also exponentially decaying as $`r\mathrm{}`$.)
Let us compare the $`𝒪(\theta ^2)`$ terms in the two eigenvalues for a general $`x_i`$ not restricted to the asymptotic region. Since the SW map is defined by the gauge equivalence relation independent of the metric, we shall consider first the simpler case of (3.68) without the last term $`+\left|\mathrm{\Delta }\widehat{\mathrm{\Phi }}_{\mathrm{metric}}^{(2)}\right|+`$ using the metric. In this case, by taking $`c_2=1/16`$, we can make (3.68) agree with the $`𝒪(\theta ^2)`$ part of (4.7) except only the $`H^3K^2`$ terms. However, for the complete agreement between the two eigenvalues, the introduction of the metric term $`+\left|\mathrm{\Delta }\widehat{\mathrm{\Phi }}_{\mathrm{metric}}^{(2)}\right|+`$ is inevitable.
As we mentioned in section 3, there are many contributions to the covariant-type ambiguity using the metric. A complete analysis shows that the term $`+\left|\mathrm{\Delta }\widehat{\mathrm{\Phi }}_{\mathrm{metric}}^{(2)}\right|+`$ is a sum of three functions, $`f_1`$ and $`f_2`$ of (3.70) and a new one
$$f_3(x,\theta )=\frac{1}{r^5}H^3K^2\theta ^2,$$
(4.10)
each multiplied by an arbitrary coefficient. In fact, we have $`+\left|\delta ^{km}\delta ^{ln}\text{Re}(F_{kl}[\mathrm{\Phi },F_{mn}])\theta ^2\right|+=f_3(x,\theta )`$. Then, expressing the RHS of $`\lambda _+^{(2)}`$ (3.68) as the sum of its first two terms and $`c_1f_1+c_2f_2+c_3f_3`$ with the redefined $`c_1`$ and $`c_2`$, the complete agreement between $`\lambda _+^{(2)}`$ and the $`𝒪(\theta ^2)`$ term of $`\mathrm{\Lambda }`$ (4.7) is achieved by taking the three parameters as $`c_1=5/32`$, $`c_2=1/16`$ and $`c_3=1/8`$. This is the unique choice for the coefficients $`c_k`$. Note that this agreement is a non-trivial one since we have to tune eight coefficients by using only three free parameters. Of course, the three coefficients $`c_k`$ do not completely fix the ambiguity in the SW map since there are many contributions to $`c_k`$ if we allow the $`\mathrm{\Delta }\widehat{\mathrm{\Phi }}_{\mathrm{metric}}^{(2)}`$ term using the metric. The use of the metric in the SW map seems not so unnatural if we recall that the noncommutative classical solution (cf. (2.9) and (2.12)) as well as the BPS equation (2.1) already contains the metric.
## 5 Summary and discussions
In this paper, we considered the noncommutative monopole solutions at the second order in $`\theta `$. We solved the noncommutative version of the BPS equation to $`𝒪(\theta ^2)`$, mapped the solution to the commutative side, and obtained the eigenvalues of the resulting scalar field. We saw that the ambiguities in the SW map have explicit influence on the scalar eigenvalues. We made the brane interpretation to the scalar eigenvalues and examined whether they can reproduce the configuration of a tilted D-string suspended between two parallel D3-branes. In the asymptotic region, the effect of the ambiguities in the SW map disappear and at the same time the scalar eigenvalue precisely give the expected D-string picture. Without the restriction to the asymptotic region, we found that we can tune the free parameters in the SW map so that the scalar eigenvalues reproduce the desired configuration. It is necessary to introduce the covariant-type ambiguity term using the metric. The number of free parameters $`c_k`$ in the eigenvalues is just enough to adjust them to the expected ones.
We would like to make a few comments. Our first comment is on the covariant type ambiguity in the SW map. In this paper we have constructed the SW map first in the pure Yang-Mills system without the scalar field and then obtained the map for the scalar by the dimensional reduction of the map for the gauge field. This is natural if we recall the origin of the present super Yang-Mills theory via the dimensional reduction. However, if we forget this origin, there are other covariant-type ambiguities treating the scalar field $`\mathrm{\Phi }`$ as a gauge covariant quantity from the start. For example, as an ambiguity for the scalar field at $`𝒪(\theta )`$, we have $`\delta \theta ^{ij}\{\mathrm{\Phi },F_{ij}\}`$. However, this term gives the same contribution (with an arbitrary coefficient) to the $`𝒪(\theta )`$ eigenvalue as the existing one (3.67), and hence even the tilt angle at $`𝒪(\theta )`$ becomes a free parameter.
Next we shall comment on the noncommutative eigenvalue equation for the scalar field proposed and examined in . At $`𝒪(\theta )`$, the eigenvalues of the noncommutative eigenvalue problem for the scalar gave the same asymptotic behavior as those obtained via the SW map . We have carried out the analysis of the noncommutative eigenvalue equation for the classical solution at $`𝒪(\theta ^2)`$ given in section 2. However, the resulting eigenvalues do not agree with those from the SW map even in the asymptotic region. Therefore, the noncommutative eigenvalue equation seems to work well only at the first order in $`\theta `$, though it is still an interesting subject to understand why it gives a good result at this order.
Finally, we would like to emphasize the usefulness of the analysis using the BPS solutions. The BPS solutions are expected to remain intact even if we include the $`\alpha ^{}`$ corrections. Thus, the BPS solutions would be helpful for giving a support for the equivalence between the noncommutative description and the commutative one independently of the $`\alpha ^{}`$ expansion. It is a very interesting subject to pursue the method which enables us to examine this equivalence to all orders in $`\theta `$.
## Acknowledgments
We would like to thank T. Asakawa, K. Hashimoto, I. Kishimoto and S. Moriyama for valuable discussions and useful comments. This work is supported in part by Grant-in-Aid for Scientific Research from Ministry of Education, Science, Sports and Culture of Japan (#03602 and #12640264). The work of S. G. is supported in part by the Japan Society for the Promotion of Science under the Predoctoral Research Program.
## Appendix A Seiberg-Witten map to $`𝒪(\delta \theta ^2)`$
In this appendix, we present the SW map for the gauge field to second order in the change $`\delta \theta `$ of the noncommutativity parameter $`\theta `$. The SW map is derived from the gauge equivalence relation
$$A_i(\widehat{A})+\delta _\lambda A_i(\widehat{A})=A_i(\widehat{A}+\widehat{\delta }_{\widehat{\lambda }}\widehat{A}),$$
(A.1)
where the quantities with a hat are defined at $`\theta `$ and those without hat at $`\theta +\delta \theta `$, a nearby point of $`\theta `$. This unconventional meaning of hat is for the convenience of the use in section 3.
We expand $`A_i`$ and $`\lambda `$ in powers of $`\delta \theta `$:
$`A_i`$ $`=\widehat{A}_i+\mathrm{\Delta }\widehat{A}_i^{(1)}+\mathrm{\Delta }\widehat{A}_i^{(2)}+𝒪(\delta \theta ^3),`$
$`\lambda `$ $`=\widehat{\lambda }+\mathrm{\Delta }\widehat{\lambda }^{(1)}+\mathrm{\Delta }\widehat{\lambda }^{(2)}+𝒪(\delta \theta ^3).`$ (A.2)
Substituting them into (A.1), the first order part is solved in the most general form as
$`\mathrm{\Delta }\widehat{A}_i^{(1)}`$ $`={\displaystyle \frac{1}{4}}\delta \theta ^{kl}\{\widehat{A}_k,_l\widehat{A}_i+\widehat{F}_{li}\}+\alpha \delta \theta ^{kl}\widehat{D}_i\widehat{F}_{kl}i\beta \delta \theta ^{kl}\widehat{D}_i[\widehat{A}_k,\widehat{A}_l],`$ (A.3)
$`\mathrm{\Delta }\widehat{\lambda }^{(1)}`$ $`={\displaystyle \frac{1}{4}}\delta \theta ^{kl}\{\widehat{A}_k,_l\widehat{\lambda }\}2i\beta \delta \theta ^{kl}[\widehat{A}_k,_l\widehat{\lambda }],`$ (A.4)
where $`\alpha `$ and $`\beta `$ are arbitrary real coefficients. Note that these two ambiguity terms are both gauge-type ones. Next, we shall solve the second order part of the equation (A.1),
$`\widehat{\delta }_{\widehat{\lambda }}\mathrm{\Delta }\widehat{A}_i^{(2)}+i[\mathrm{\Delta }\widehat{A}_i^{(2)},\widehat{\lambda }]\widehat{D}_i\mathrm{\Delta }\widehat{\lambda }^{(2)}={\displaystyle \frac{i}{8}}\delta \theta ^{kl}\delta \theta ^{mn}[_k_m\widehat{A}_i,_l_n\widehat{\lambda }]`$
$`+{\displaystyle \frac{1}{2}}\delta \theta ^{kl}\{_k\widehat{A},_l\mathrm{\Delta }\widehat{\lambda }^{(1)}\}+{\displaystyle \frac{1}{2}}\delta \theta ^{kl}\{_k\mathrm{\Delta }\widehat{A}_i^{(1)},_l\widehat{\lambda }\}i[\mathrm{\Delta }\widehat{A}_i^{(1)},\mathrm{\Delta }\widehat{\lambda }^{(1)}],`$ (A.5)
to obtain $`\mathrm{\Delta }\widehat{A}_i^{(2)}`$ and $`\mathrm{\Delta }\widehat{\lambda }^{(2)}`$. We solved this equation (A.5) by assuming the most general forms for $`\mathrm{\Delta }\widehat{A}_i^{(2)}`$ and $`\mathrm{\Delta }\widehat{\lambda }^{(2)}`$. The result is as follows: The terms in $`\mathrm{\Delta }\widehat{\mathrm{\Phi }}^{(2)}`$ of the form $`\text{Re}𝒪`$ ($`\text{Im}𝒪`$) with $`𝒪`$ containing odd (even) number of derivatives do not contribute to the scalar eigenvalue formula (3.65). Therefore, in eq. (A.62) we have omitted such kind of terms, which would appear as the covariant-type ambiguity terms and the terms quadratic in $`\alpha `$ and $`\beta `$.
$`\mathrm{\Delta }\widehat{A}_i^{(2)}`$ $`={\displaystyle \frac{1}{4}}\text{Im}(\begin{array}{c}\text{ }\text{ }\hfill \\ AA_i\hfill \end{array}){\displaystyle \frac{1}{8}}\text{Im}(\begin{array}{c}\text{ }\text{ }\hfill \\ A_iA\hfill \end{array})+{\displaystyle \frac{1}{4}}\text{Re}(\begin{array}{c}\text{ }\text{ }\hfill \\ AAA_i\hfill \end{array}){\displaystyle \frac{1}{4}}\text{Re}(\begin{array}{c}\text{ }\text{ }\hfill \\ AAA_i\hfill \end{array})`$ (A.14)
14Re( AFAi)+14Re( AAFi)+14Re(
AF
Fi
)+14Re(
A
Fi
F)14Refragments 𝐴𝐹subscript𝐴𝑖14Refragments 𝐴𝐴subscript𝐹𝑖14Re
AF
Fi
14Re
A
Fi
𝐹\displaystyle\quad-\frac{1}{4}\mathop{\mbox{Re}}(\mathop{\vbox{\halign{#\cr\kern 1.72218pt\cr$\hbox{$\hskip 3.75pt\vrule height=3.65973pt,width=0.29999pt,depth=-0.0004pt\kern-0.29999pt\vrule height=3.65973pt,width=8.90446pt,depth=-3.35974pt\kern-0.29999pt\vrule height=3.65973pt,width=0.29999pt,depth=-0.11153pt$\hss}\hbox{$\hskip 18.3854pt\vrule height=3.65973pt,width=0.29999pt,depth=-0.0004pt\kern-0.29999pt\vrule height=3.65973pt,width=11.25563pt,depth=-3.35974pt\kern-0.29999pt\vrule height=3.65973pt,width=0.29999pt,depth=-0.0004pt$\hss}$\crcr\kern 1.72218pt\nointerlineskip\cr\hbox{$\displaystyle{}{A}{\,}{\partial}{}{F}{{}_{i}\,}{A}$}\crcr}}}\limits)+\frac{1}{4}\mathop{\mbox{Re}}(\mathop{\vbox{\halign{#\cr\kern 1.72218pt\cr$\hbox{$\hskip 3.75pt\vrule height=3.65973pt,width=0.29999pt,depth=-0.0004pt\kern-0.29999pt\vrule height=3.65973pt,width=18.0711pt,depth=-3.35974pt\kern-0.29999pt\vrule height=3.65973pt,width=0.29999pt,depth=-0.11153pt$\hss}\hbox{$\hskip 12.91664pt\vrule height=7.31946pt,width=0.29999pt,depth=-0.0004pt\kern-0.29999pt\vrule height=7.31946pt,width=14.6354pt,depth=-7.01947pt\kern-0.29999pt\vrule height=7.31946pt,width=0.29999pt,depth=-0.0004pt$\hss}$\crcr\kern 1.72218pt\nointerlineskip\cr\hbox{$\displaystyle{}{A}{\,}{A}{\,}{\partial}{}{F}_{i}$}\crcr}}}\limits)+\frac{1}{4}\mathop{\mbox{Re}}(\vbox{\halign{#\cr\hskip 5.0pt\vrule depth=4.0pt\hrulefill\vrule depth=4.0pt\hskip 5.0pt\crcr\kern-1.0pt\vskip 3.55658pt\nointerlineskip\cr$\hfil\displaystyle{AF}\hfil$\crcr}}\vbox{\halign{#\cr\hskip-2.0pt\vrule depth=4.0pt\hrulefill\vrule depth=4.0pt\hskip 5.0pt\crcr\kern-1.0pt\vskip 3.55658pt\nointerlineskip\cr$\hfil\displaystyle{F_{i}}\hfil$\crcr}})+\frac{1}{4}\mathop{\mbox{Re}}(\vbox{\halign{#\cr\hskip 5.0pt\vrule depth=8.0pt\hrulefill\vrule depth=8.0pt\hskip-19.0pt\crcr\kern-1.0pt\vskip 3.55658pt\nointerlineskip\cr$\hfil\displaystyle{A}\hfil$\crcr}}\vbox{\halign{#\cr\hskip 5.0pt\vrule depth=4.0pt\hrulefill\vrule depth=4.0pt\hskip-4.0pt\crcr\kern-1.0pt\vskip 3.55658pt\nointerlineskip\cr$\hfil\displaystyle{F_{i}}\hfil$\crcr}}F) (A.27)
$`{\displaystyle \frac{1}{8}}\text{Re}(\begin{array}{c}\text{ }\text{ }\hfill \\ AAA_i\hfill \end{array})+{\displaystyle \frac{1}{8}}\text{Re}(\begin{array}{c}\text{ }\text{ }\hfill \\ AA_iA\hfill \end{array}){\displaystyle \frac{1}{8}}\text{Im}(\begin{array}{c}\text{ }\text{ }\hfill \\ AAA_iA\hfill \end{array})+{\displaystyle \frac{1}{8}}\text{Im}(\begin{array}{c}\text{ }\text{ }\hfill \\ AA_iAA\hfill \end{array})`$ (A.36)
$`{\displaystyle \frac{1}{8}}\text{Re}(\begin{array}{c}\text{ }\text{ }\hfill \\ AAAAA_i\hfill \end{array})+{\displaystyle \frac{1}{4}}\text{Re}(\begin{array}{c}\text{ }\text{ }\hfill \\ AAAA_iA\hfill \end{array}){\displaystyle \frac{1}{8}}\text{Re}(\begin{array}{c}\text{ }\text{ }\hfill \\ AAA_iAA\hfill \end{array})`$ (A.43)
$`+\left({\displaystyle \frac{1}{16}}+8\alpha \beta +4\beta ^2\right)\text{Im}(\begin{array}{c}\text{ }\text{ }\hfill \\ AAD_i(AA)\hfill \end{array})+8\alpha \beta \text{Re}(\begin{array}{c}\text{ }\text{ }\hfill \\ AD_i(AA)\hfill \end{array})`$ (A.48)
$`+\gamma _1\text{Im}(\begin{array}{c}F\hfill \end{array}D_i\begin{array}{c}F\hfill \end{array})+\gamma _2\text{Im}(\begin{array}{c}F\hfill \end{array}D_i\begin{array}{c}F\hfill \end{array})+\gamma _3\text{Im}(\begin{array}{c}\text{ }\hfill \\ F{}_{i}{}^{}D\hfill \end{array}\begin{array}{c}F\hfill \end{array})`$ (A.61)
$`+\mathrm{\Delta }\widehat{A}_{i\mathrm{metric}}^{(2)}+(\text{gauge-type ambiguities}),`$ (A.62)
$`\mathrm{\Delta }\widehat{\lambda }^{(2)}`$ $`={\displaystyle \frac{1}{8}}\text{Im}(\begin{array}{c}\text{ }\text{ }\hfill \\ A\lambda \hfill \end{array}){\displaystyle \frac{1}{8}}\text{Re}(\begin{array}{c}\text{ }\text{ }\hfill \\ AA\lambda \hfill \end{array}){\displaystyle \frac{1}{8}}\text{Re}(\begin{array}{c}\text{ }\text{ }\hfill \\ A\lambda A\hfill \end{array})+{\displaystyle \frac{1}{4}}\text{Re}(\begin{array}{c}\text{ }\text{ }\hfill \\ AA\lambda \hfill \end{array})`$ (A.71)
$`{\displaystyle \frac{1}{8}}\text{Im}(\begin{array}{c}\text{ }\text{ }\hfill \\ AAA\lambda \hfill \end{array})+{\displaystyle \frac{1}{8}}\text{Im}(\begin{array}{c}\text{ }\text{ }\hfill \\ AA\lambda A\hfill \end{array})`$ (A.76)
$`+\left({\displaystyle \frac{1}{16}}+8\alpha \beta +4\beta ^2\right)\text{Im}(\begin{array}{c}\text{ }\hfill \\ AA\hfill \end{array}(\begin{array}{c}\text{ }\hfill \\ A\lambda \hfill \end{array}+\begin{array}{c}\text{ }\hfill \\ \lambda A\hfill \end{array}))+8\alpha \beta \text{Re}(\begin{array}{c}\text{ }\hfill \\ A\hfill \end{array}(\begin{array}{c}\text{ }\hfill \\ A\lambda \hfill \end{array}+\begin{array}{c}\text{ }\hfill \\ \lambda A\hfill \end{array}))`$ (A.89)
$`+(\text{gauge-type ambiguities}),`$ (A.90)
where the meanings of the contraction, $`\text{Re}𝒪`$ and $`\text{Im}𝒪`$ are as given by eqs. (3.55) and (3.56). We have omitted hats on the RHS of (A.62) and (A.90). The ambiguities of the SW map at $`𝒪(\theta ^2)`$ are the homogeneous solutions to eq. (A.5). The terms in (A.62) multiplied by $`\gamma _k`$ ($`k=1,2,3`$) are the covariant-type ambiguities which cannot be identified as gauge transformation. All other covariant-type terms are reduced to the three $`\gamma _k`$ terms owing to the Bianchi identity. The term $`\mathrm{\Delta }\widehat{A}_{i\mathrm{metric}}^{(2)}`$ denotes the covariant-type ambiguity using the metric $`G_{ij}`$. There are many operators belonging to this type; for example,
$$G^{km}G^{ln}\text{Im}(F_{kl}D_iF_{mn})\theta ^2,G^{kp}G^{mq}G^{ln}\text{Im}(F_{kl}D_iF_{mn})\theta _p\theta _q.$$
(A.91)
The SW map for the scalar field $`\mathrm{\Phi }`$ used in section 3 is obtained from (A.62) by the dimensional reduction using $`A_\mathrm{\Phi }=\mathrm{\Phi }`$, $`F_{i\mathrm{\Phi }}=D_i\mathrm{\Phi }`$ and $`D_\mathrm{\Phi }𝒪=i[\mathrm{\Phi },𝒪]`$. The second quantity in (3.57) is obtained from that in (A.91) by setting $`i=\mathrm{\Phi }`$ and taking the $`G^{\mathrm{\Phi }\mathrm{\Phi }}`$ part.
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# Triangular dynamical 𝑟-matrices and quantization
## 1 Introduction
In the last two decades, the theory of quantum groups has undergone tremendous development. The classical counterparts of quantum groups are Lie bialgebras . Many interesting quantum groups were found and studied by various authors, but the proof of existence of quantization for arbitrary Lie bialgebras was obtained only recently by Etingof and Kazhdan . For triangular Lie bialgebras, however, an elementary proof of quantization was given by Drinfeld in 1983 . Drinfeld’s idea can be outlined as follows. A triangular $`r`$-matrix on a Lie algebra $`g`$ defines a left invariant Poisson structure on its corresponding Lie group $`G`$. By restricting to a Lie subalgebra if necessary, one may in fact assume that this is symplectic. One may then quantize the $`r`$-matrix by finding a $`G`$-invariant $``$-product on $`G`$, of which there may be several. In , Drinfeld identified the symplectic manifold with a coadjoint orbit of a central extension of $`g`$, and then applied Berezin quantization .
Recently, there has been growing interest in the so-called quantum dynamical Yang-Baxter equation (see Equation (13)). This equation arises naturally from various contexts in mathematical physics. It first appeared in the work of Gervais-Neveu in their study of quantum Liouville theory . Recently it reappeared in Felder’s work on the quantum Knizhnik-Zamolodchikov-Bernard equation. It also has been found to be connected with the quantum Caloger-Moser systems . Just like the quantum Yang-Baxter equation is connected with quantum groups, the quantum dynamical Yang-Baxter equation is known to be connected with elliptic quantum groups , as well as with Hopf algebroids or quantum groupoids .
The classical counterpart of the quantum dynamical Yang-Baxter equation was first considered by Felder , and then studied by Etingof and Varchenko . This is the so-called classical dynamical Yang-Baxter equation, and a solution to such an equation (plus some other reasonable conditions) is called a classical dynamical $`r`$-matrix. More precisely, given a Lie algebra $`g`$ over $`R`$ (or over $`C`$) with an Abelian Lie subalgebra $`h`$, a classical dynamical $`r`$-matrix is a smooth (or meromorphic) function $`r(\lambda ):h^{}gg`$ satisfying the following conditions:
1. (zero weight condition) $`[h1+1h,r(\lambda )]=0,hh`$;
2. (normal condition) $`r^{12}+r^{21}=\mathrm{\Omega }`$, where $`\mathrm{\Omega }(S^2g)^g`$ is a Casimir element;
3. (classical dynamical Yang-Baxter equation)
$$\text{Alt}(dr)+[r^{12},r^{13}]+[r^{12},r^{23}]+[r^{13},r^{23}]=\mathrm{\hspace{0.17em}0},$$
(1)
where $`\text{Alt}dr=(h_i^{(1)}\frac{r^{23}}{\lambda ^i}h_i^{(2)}\frac{r^{13}}{\lambda ^i}+h_i^{(3)}\frac{r^{12}}{\lambda ^i})`$.
A fundamental question is whether any classical dynamical $`r`$-matrix is quantizable. There have appeared many results in this direction. For the standard classical dynamical $`r`$-matrix for $`sl_2(C)`$, a quantization was obtained by Babelon in 1991. For general simple Lie algebras, quantizations were recently found independently by Arnaudon et al. and Jimbo et al. based on the approach of Fronsdal . Similar results were also found by Etingof and Varchenko using intertwining operators. Recently, using a method similar to , Etingof et al. obtained a quantization of all the classical dynamical $`r`$-matrices of semi-simple Lie algebras in Schiffmann’s classification list . However, the general quantization problem still remains open; a recipe has yet to be found. Moreover, the problem of classification of quantizations has not yet been touched.
In this paper, we study the quantization problem for general classical triangular dynamical $`r`$-matrices. Classical triangular dynamical $`r`$-matrices are those satisfying the skew-symmetric condition $`r^{12}(\lambda )+r^{21}(\lambda )=0`$. In this case, Equation (1) is equivalent to $`_ih_i\frac{r}{\lambda ^i}+\frac{1}{2}[r,r]=0`$. These $`r`$-matrices are in one-one correspondence with regular Poisson structures $`\pi =_i\stackrel{}{h_i}\frac{}{\lambda ^i}+\stackrel{}{r(\lambda )}`$ on the manifold $`h^{}\times G`$, which are invariant under the left $`G`$ and right $`H`$-actions. Thus one may expect to quantize a classical dynamical $`r`$-matrix by looking for a certain special type of star-products on the corresponding Poisson manifold. This is exactly the route we take in the present paper. In some sense, this is also a natural generalization of the quantization method used by Drinfeld in as outlined at the beginning of the introduction. In fact, in the present paper, we mainly deal with non-degenerate triangular classical dynamical $`r`$-matrices (i.e., the corresponding Poisson manifolds are in fact symplectic). Berezin quantization no longer works in this situation. However, one may use the Fedosov method to obtain the desired star-products as we will see later. It is well-known that star products on a symplectic manifold are classified by the second cohomology group of the manifold with coefficients in formal $`\mathrm{}`$-power series. In light of this result, we are able to classify the quantizations of a non-degenerate triangular classical dynamical $`r`$-matrix and prove that the quantizations are parameterized by the relative Lie algebra cohomology $`H^2(g,h)[[\mathrm{}]]`$.
For a general triangular classical dynamical $`r`$-matrix, it is natural to ask whether it is possible to reduce it to a non-degenerate one by restricting to a Lie subalgebra. This is always true in the non-dynamical case . Unfortunately, in general this fails in the dynamical case, and we will study the conditions under which this is possible. In this case, these $`r`$-matrices are called splittable. Splittable triangular classical dynamical $`r`$-matrices resemble in many ways non-degenerate ones. And in particular, they can be quantized by the Fedosov method.
The outline of this paper is as follows. After Section 1 (this introduction), in Section 2, we study general properties of triangular classical dynamical $`r`$-matrices. It is proved that triangular classical dynamical $`r`$-matrices correspond to some special Poisson structures on $`h^{}\times G`$, which are always regular. This may seem surprising at first glance since the rank of $`r(\lambda )`$ may depend on the point $`\lambda `$. The main tool in Section 2 is the method of Lie groupoids and Lie algebroids. In particular, we show how gauge transformations, first introduced by Etingof and Varchenko , enter naturally from the viewpoint of Lie algebroids. The study of the tangent space of the moduli space of dynamical $`r`$-matrices naturally leads to the notion of dynamical $`r`$-matrix cohomology, which is shown to be isomorphic to the relative Lie algebra cohomology when $`r`$ is non-degenerate. Section 3 is devoted to the proof of the equivalence between quantizations of triangular classical dynamical $`r`$-matrices and the so called compatible star products on their corresponding Poisson manifolds $`h^{}\times G`$. In Section 4, we study symplectic connections on such symplectic manifolds ($`M=h^{}\times G`$). In particular, we show that there always exists a $`G\times H`$-invariant (i.e. left $`G`$-invariant and right $`H`$-invariant) torsion-free symplectic connection on $`M`$ such that the left invariant vector fields $`\stackrel{}{h},hh`$ are all parallel. The main result of Section 5 is that the Fedosov quantization obtained via such a symplectic connection and some suitable choice of Weyl curvatures gives rise to compatible $``$-products on $`M=h^{}\times G`$. Therefore, as a consequence, we prove the existence of a quantization of non-degenerate triangular classical dynamical $`r`$-matrices. The presentation in Section 5, however, is made in a more general setting, which is of its own interest. Section 6 is devoted to the classification of quantizations. In particular, we show that the equivalence classes of quantizations of a non-degenerate triangular dynamical $`r`$-matrix $`r:h^{}^2g`$ are parameterized by the relative Lie algebra cohomology with coefficients in the formal $`\mathrm{}`$-power series $`H^2(g,h)[[\mathrm{}]]`$. Some speculation on the classification of quantizations of a general triangular classical dynamical $`r`$-matrix is given as a conjecture, which is consistent with Kontesvich’s formality theorem . In the appendix we recall some basic ingredients of the Fedosov quantization, which are used throughout the paper.
Finally, some remarks are in order. Quantization of dynamical $`r`$-matrices is related to quantization of Lie bialgebroids as shown in . However, for simplicity, we will avoid using quantum groupoids in the present paper even though many ideas are rooted from there. Also in this paper, we work in the smooth case. Namely, Lie algebras are finite dimensional Lie algebras over $`R`$, all manifolds and maps are smooth, but our approach works for the complex category as well. For simplicity, we assume that a dynamical $`r`$-matrix is always defined on $`h^{}`$. In reality, it may only be defined on an open submanifold $`Uh^{}`$, but our results hold in this situation as well.
Acknowledgments. The author would like to thank Martin Bordemann, Pavel Etingof and Boris Tysgan for useful discussions. Especially, he is grateful to Pavel Etingof for his suggestion of writing up this work. In addition to the funding sources mentioned in the first footnote, he would also like to thank the Max-Planck Institut for the hospitality and financial support while part of this project was being done.
## 2 Triangular dynamical $`r`$-matrices
In this section, we study some general aspects of triangular dynamical $`r`$-matrices. As a useful tool, we shall utilize the method of Lie algebroids and Lie groupoids. Let $`g`$ be a Lie algebra and $`hg`$ an Abelian Lie subalgebra of dimension $`l`$. By a triangular dynamical $`r`$-matrix, we mean a smooth function $`r:h^{}^2g`$ satisfying:
1. the zero weight condition: $`[h,r(\lambda )]=0,\lambda h^{},hh`$, and
2. the classical dynamical Yang-Baxter equation (CDYBE):
$$\underset{i}{}h_i\frac{r}{\lambda ^i}+\frac{1}{2}[r,r]=0,$$
(2)
where the bracket $`[,]`$ refers to the Schouten type bracket: $`^kg^lg^{k+l1}g`$ induced from the Lie algebra bracket on $`g`$. Here $`\{h_1,\mathrm{},h_l\}`$ is a basis in $`h`$, and $`(\lambda ^1,\mathrm{},\lambda ^l)`$ its induced coordinate system on $`h^{}`$. It is known that the CDYBE is closely related to Lie bialgebroids. Recall that a Lie bialgebroid is a pair of Lie algebroids ($`A`$, $`A^{}`$) satisfying the following compatibility condition (see ):
$$d_{}[X,Y]=[d_{}X,Y]+[X,d_{}Y],X,Y\mathrm{\Gamma }(A),$$
(3)
where the differential $`d_{}`$ on $`\mathrm{\Gamma }(^{}A)`$ comes from the Lie algebroid structure on $`A^{}`$.
Given a Lie algebroid $`A`$ over $`P`$ with anchor $`a`$, and a section $`\mathrm{\Lambda }`$ of $`\mathrm{\Gamma }(^2A)`$ satisfying the condition $`[\mathrm{\Lambda },\mathrm{\Lambda }]=0`$, one may define a Lie algebroid structure on $`A^{}`$ by simply requiring the differential $`d_{}:\mathrm{\Gamma }(^kA)\mathrm{\Gamma }(^{k+1}A)`$ to be $`d_{}=[\mathrm{\Lambda },]`$. More explicitly, denote by $`\mathrm{\Lambda }^\mathrm{\#}`$ the bundle map $`A^{}A`$ defined by $`\mathrm{\Lambda }^\mathrm{\#}(\xi )(\eta )=\mathrm{\Lambda }(\xi ,\eta ),\xi ,\eta \mathrm{\Gamma }(A^{})`$. Then the bracket on $`\mathrm{\Gamma }(A^{})`$ is defined by
$$[\xi ,\eta ]=L_{\mathrm{\Lambda }^\mathrm{\#}\xi }\eta L_{\mathrm{\Lambda }^\mathrm{\#}\eta }\xi d[\mathrm{\Lambda }(\xi ,\eta )],$$
(4)
and the anchor $`a_{}`$ is the composition $`a\mathrm{\Lambda }^\mathrm{\#}:A^{}TP`$. It is easy to show that $`(A,A^{})`$ is indeed a Lie bialgebroid, which is called a triangular Lie bialgebroid .
Now consider $`A=Th^{}\times g`$ and equip $`A`$ with the standard product Lie algebroid structure. Then the anchor $`a:Th^{}\times gTh^{}`$ is simply the projection. The relation between triangular dynamical $`r`$-matrices and triangular Lie bialgebroids are described by the following :
###### Proposition 2.1
Given a smooth function $`r:h^{}^2g`$, $`r`$ is a triangular dynamical $`r`$-matrix iff the Lie algebroid $`(A,a)`$ together with $`\mathrm{\Lambda }=_ih_i\frac{}{\lambda ^i}+r(\lambda )\mathrm{\Gamma }(^2A)`$ defines a triangular Lie bialgebroid.
Proof. By a straightforward computation, we have $`[\mathrm{\Lambda },\mathrm{\Lambda }]=2(_ih_i\frac{r}{\lambda ^i}+\frac{1}{2}[r,r]+_i[r,h_i]\frac{}{\lambda ^i})`$. It thus follows that $`[\mathrm{\Lambda },\mathrm{\Lambda }]=0`$ iff $`_ih_i\frac{r}{\lambda ^i}+\frac{1}{2}[r,r]=0`$ and $`[r,h_i]=0(i=1,\mathrm{}l`$), i.e., $`r`$ is a triangular dynamical $`r`$-matrix.
$`\mathrm{}`$
Let $`G`$ be a Lie group with Lie algebra $`g`$ and $`HG`$ an Abelian Lie subgroup with Lie algebra $`h`$. Consider $`M=h^{}\times G`$. Let $`G`$ act on $`M`$ from the left by left multiplication on $`G`$, and $`H`$ act from the right by right multiplication on $`G`$. An equivalent version of Proposition 2.1 is the following
###### Proposition 2.2
For a smooth function $`r:h^{}^2g`$, $`r`$ is a triangular dynamical $`r`$-matrix iff $`\pi =_i\stackrel{}{h_i}\frac{}{\lambda ^i}+\stackrel{}{r(\lambda )}`$ defines a $`G\times H`$-invariant Poisson structure on $`M=h^{}\times G`$, where $`\stackrel{}{h_i}X(M)`$ is the left invariant vector field on $`M`$ generated by $`h_i`$ and similarly $`\stackrel{}{r(\lambda )}\mathrm{\Gamma }(^2TM)`$ is the left invariant bivector field on $`M`$ corresponding to $`r(\lambda )`$.
###### Theorem 2.3
If $`r:h^{}^2g`$ is a triangular dynamical $`r`$-matrix, then $`h+r(\lambda )^\mathrm{\#}h^{}`$ is a Lie subalgebra of $`g`$. Moreover the Lie subalgebras $`h+r(\lambda )^\mathrm{\#}h^{}`$, $`\lambda h^{}`$, are all isomorphic, and the isomorphisms are given by the adjoint action of $`G`$.
Proof. For any $`\lambda h^{}`$, $`A_\lambda =T_\lambda h^{}gh^{}g`$ and $`A_\lambda ^{}hg^{}`$. Under these identifications, the bundle map $`\mathrm{\Lambda }_\lambda ^\mathrm{\#}:A_\lambda ^{}A_\lambda `$ is given by
$$(h,\xi )(i^{}\xi ,h+r(\lambda )^\mathrm{\#}\xi ),hh\text{ and }\xi g^{},$$
(5)
where $`i:hg`$ is the inclusion. Set $`B=\mathrm{\Lambda }^\mathrm{\#}(A^{})=_{\lambda h^{}}\mathrm{\Lambda }_\lambda ^\mathrm{\#}(A_\lambda ^{})A`$. Since $`(A,\mathrm{\Lambda })`$ defines a triangular Lie bialgebroid, $`B`$ is integrable. I.e., $`\mathrm{\Gamma }(B)`$ is closed under the Lie algebroid bracket on $`\mathrm{\Gamma }(A)`$. Hence $`\text{ker}a|_{B_\lambda }`$ is a Lie subalgebra of $`\text{ker}a|_{A_\lambda }`$. Now it is easy to see that $`\text{ker}a|_{B_\lambda }=h+r(\lambda )^\mathrm{\#}h^{}`$ and $`\text{ker}a|_{A_\lambda }=g`$. It thus follows that $`h+r(\lambda )^\mathrm{\#}h^{}`$ is a Lie subalgebra of $`g`$. On the other hand, from Equation (5), it is easy to see that $`a(B_\lambda )=T_\lambda h^{}`$. Hence $`a:BTh^{}`$ is surjective, which implies that $`B`$ is in fact a transitive Lie algebroid (also called a gauge Lie algebroid ). Thus it follows that the dimension of $`B_\lambda `$ is independent of $`\lambda `$, and therefore $`B`$ is a subbundle of $`A`$. Moreover the isotropic Lie algebras of $`B`$ at different points of $`h^{}`$ are all isomorphic, and the isomorphisms are given by the adjoint action of $`G`$. This implies that, for any $`\lambda ,\mu h^{}`$, $`h+r(\lambda )^\mathrm{\#}h^{}`$ is isomorphic to $`h+r(\mu )^\mathrm{\#}h^{}`$ by the adjoint action of a group element in $`G`$.
$`\mathrm{}`$
For the sake of simplicity, we denote by $`g_\lambda `$ the Lie subalgebra $`h+r(\lambda )^\mathrm{\#}h^{}`$. Define the rank of a triangular dynamical $`r`$-matrix $`r`$ to be $`\text{dim}g_\lambda \text{dim}h`$, which is denoted as rank$`r`$. We say a triangular dynamical $`r`$-matrix $`r`$ is non-degenerate if rank$`r=\text{dim}g\text{dim}h`$.
An immediate consequence of Theorem 2.3 is
###### Corollary 2.4
Under the same hypothesis as in Theorem 2.3, rank$`r`$ is independent of the point $`\lambda `$ and therefore is a well-defined even number. Moreover $`B=\mathrm{\Lambda }^\mathrm{\#}A^{}A`$ is a Lie subalgebroid of rank $`2\text{dim}h`$+rank$`r`$, and $`(M,\pi )`$ is a regular Poisson manifold of rank $`2\text{dim}h`$+rank$`r`$.
In particular, we have the following
###### Corollary 2.5
Given a triangular dynamical $`r`$-matrix $`r:h^{}^2g`$, the following statements are all equivalent:
1. $`r`$ is non-degenerate;
2. the bundle map $`\mathrm{\Lambda }^\mathrm{\#}:A^{}A`$ is nondegenerate;
3. $`g_\lambda =g,\lambda h^{}`$;
4. $`(M,\pi )`$ is a symplectic manifold.
If we choose a decomposition $`g=hm`$, where $`m`$ is a subspace of $`g`$, and choose a basis $`\{h_1,\mathrm{},h_l\}`$ for $`h`$ and a basis $`\{e_1,\mathrm{},e_m\}`$ for $`m`$, we may write
$$r(\lambda )=a^{ij}(\lambda )h_ih_j+b^{ij}(\lambda )h_ie_j+c^{ij}(\lambda )e_ie_j.$$
(6)
It is simple to see that $`g_\lambda =h\text{Span}\{_jc^{ij}(\lambda )e_j|i=1,\mathrm{},m\}`$, and rank$`r`$ is the rank of the matrix $`(c^{ij}(\lambda ))`$. Therefore, we immediately know that the rank of $`(c^{ij}(\lambda ))`$ is independent of $`\lambda `$. Clearly $`r`$ is non-degenerate iff the matrix $`(c^{ij}(\lambda ))`$ is non-degenerate.
A natural question arises as to whether it is possible to make an arbitrary triangular dynamical $`r`$-matrix non-degenerate by considering it to be valued in a Lie subalgebra of $`g`$. This is true in the non-dynamical case , for example. However, in the dynamical case, this is not always possible as we will see below. Nevertheless we will single out those $`r`$-matrices possessing this property, which will be called splittable. Splittable triangular dynamical $`r`$-matrices contain a large class of interesting dynamical $`r`$-matrices, which in fact include almost all examples we know, e.g., those as classified in when $`g`$ is a simple Lie algebra. More precisely,
###### Definition 2.6
A triangular dynamical $`r`$-matrix $`r:h^{}^2g`$ is said to be splittable if for any $`\lambda h^{}`$, $`i^{}(r(\lambda )^{\mathrm{\#}1}h)=h^{}`$, where $`i:hg`$ is the inclusion.
###### Proposition 2.7
Suppose that $`r`$ is a triangular dynamical $`r`$-matrix. Then the following statements are equivalent:
1. $`r`$ is splittable;
2. for any $`\lambda h^{}`$, $`r(\lambda )^\mathrm{\#}g^{}g_\lambda `$;
3. if $`r(\lambda )`$ is given as in Equation (6) under a decomposition $`g=hm`$, then for any $`i`$, $`_jb^{ij}(\lambda )e_j\text{Span}\{_jc^{ij}(\lambda )e_j|i=1,\mathrm{},m\}`$;
4. for any fixed $`\lambda h^{}`$, there exists a decomposition $`g=hm`$, under which
$$r(\lambda )=a^{ij}(\lambda )h_ih_j+c^{ij}(\lambda )e_ie_j;$$
(7)
5. $`Th^{}\times \{0\}B`$.
Let us first prove the following simple lemma from linear algebra.
###### Lemma 2.8
Let $`V=hm`$ be a decomposition of vector spaces, and let $`\{h_1,\mathrm{},h_l\}`$ be a basis of $`h`$, and $`\{e_1,\mathrm{},e_m\}`$ a basis of $`m`$. Let $`r^2V`$ be any element such that
$$r=a^{ij}h_ih_j+h_ix^i+c^{ij}e_ie_j,$$
where $`x^im`$, and $`a_{ij},c_{ij}`$ are skew-symmetric, i.e., $`a_{ij}=a_{ji}`$ and $`c_{ij}=c_{ji}`$. If $`I\{1,\mathrm{},l\}`$ is a subset of indexes such that for any $`i_0I`$, $`x^{i_0}Span\{_jc^{ij}e_j|i=1,\mathrm{},m\}`$. Then one can change the decomposition $`V=h\stackrel{~}{m}`$ so that under a suitable basis $`\{\stackrel{~}{e_1},\mathrm{},\stackrel{~}{e_m}\}`$ of $`\stackrel{~}{m}`$, $`r`$ can be written as
$$r=\stackrel{~}{a}^{ij}h_ih_j+\underset{iI}{}h_ix^i+c^{ij}\stackrel{~}{e}_i\stackrel{~}{e}_j.$$
Proof. $`i_0I`$, by assumption, there are constants $`\gamma _i^{i_0},i=1,\mathrm{},m`$, such that $`x^{i_0}=2_{ij}\gamma _i^{i_0}c^{ij}e_j`$. Let $`\stackrel{~}{e}_i=e_i+_{i_0I}\gamma _i^{i_0}h_{i_0},i=1,\mathrm{},m`$. Then
$`{\displaystyle c^{ij}\stackrel{~}{e}_i}\stackrel{~}{e}_j`$
$`=`$ $`{\displaystyle c^{ij}(e_i+\underset{i_0I}{}\gamma _i^{i_0}h_{i_0})}(e_j+{\displaystyle \underset{i_0I}{}}\gamma _j^{i_0}h_{i_0})`$
$`=`$ $`{\displaystyle c^{ij}e_i}e_j+2{\displaystyle c^{ij}\gamma _i^{i_0}h_{i_0}}e_j(\text{mod}^2h)`$
$`=`$ $`{\displaystyle c^{ij}e_i}e_j+{\displaystyle \underset{i_0I}{}}h_{i_0}x^{i_0}(\text{mod}^2h).`$
Hence $`r=c^{ij}\stackrel{~}{e}_i\stackrel{~}{e}_j+_{iI}h_ix^i(\text{mod}^2h)`$. This concludes the proof.
$`\mathrm{}`$
Proof of Proposition 2.7
(i)$``$(ii) Let us fix a basis $`\{h_1,\mathrm{},h_l\}`$ of $`h`$, and let $`\{h_{}^1,\mathrm{},h_{}^l\}`$ be its dual basis in $`h^{}`$. By assumption, for any $`1jl`$, there is a $`\xi ^jg^{}`$ such that $`i^{}\xi ^j=h_{}^j`$ and $`r(\lambda )^\mathrm{\#}\xi ^jh`$. Given any $`\xi g^{}`$, take $`a_j=<\xi ,h_j>`$ and $`\eta =\xi a_j\xi ^j`$. Then it is easy to see that $`\eta h^{}`$. Hence $`r(\lambda )^\mathrm{\#}\xi =a_jr(\lambda )^\mathrm{\#}\xi ^j+r(\lambda )^\mathrm{\#}\eta h+r(\lambda )^\mathrm{\#}h^{}=g_\lambda `$.
(ii)$``$(iii) Let $`\{h_1,\mathrm{},h_l\}`$ be a basis of $`h`$, $`\{e_1,\mathrm{},e_m\}`$ a basis of $`m`$, and $`\{h_{}^1,\mathrm{},h_{}^l,e_{}^1,\mathrm{},e_{}^m\}`$ the dual basis of $`\{h_1,\mathrm{},h_l,e_1,\mathrm{},e_m\}`$ in $`g^{}`$. It is trivial to see that $`r(\lambda )^\mathrm{\#}e_{}^i=_jb^{ji}(\lambda )h_j+2_jc^{ij}(\lambda )e_j`$. Hence we have
$$g_\lambda =h\text{Span}\{\underset{j}{}c^{ij}(\lambda )e_j|i=1,\mathrm{},m\}.$$
Now $`r(\lambda )^\mathrm{\#}h_{}^i=_j2a^{ij}(\lambda )h_j+_jb^{ij}(\lambda )e_j`$. Since $`r(\lambda )^\mathrm{\#}h_{}^ig_\lambda `$ by assumption, it follows that $`_jb^{ij}(\lambda )e_j\text{Span}\{_jc^{ij}(\lambda )e_j|i=1,\mathrm{},m\}`$.
(iii)$``$(vi) This follows from Lemma 2.8.
(vi)$``$(v) If $`r(\lambda )=a^{ij}(\lambda )h_ih_j+c^{ij}(\lambda )e_ie_j`$, then $`r(\lambda )^\mathrm{\#}h_{}^i=2_ja^{ij}(\lambda )h_j`$. Thus according to Equation (5), $`\mathrm{\Lambda }_\lambda ^\mathrm{\#}(2_ja^{ij}(\lambda )h_j,h_{}^i)=(h_{}^i,0)`$. Hence, $`(h_{}^i,0)B_\lambda `$. This implies that $`T_\lambda h^{}\times \{0\}B_\lambda `$.
(v)$``$(i) Given any $`\phi h^{}`$, we know that $`(\phi ,0)B_\lambda `$ by assumption. Therefore there exist $`hh`$ and $`\xi g^{}`$ such that $`\mathrm{\Lambda }_\lambda ^\mathrm{\#}(h,\xi )=(\phi ,0)`$, i.e., $`(i^{}\xi ,h+r(\lambda )^\mathrm{\#}\xi )=(\phi ,0)`$ according to Equation (5). This implies that $`\phi =i^{}\xi `$ and $`r(\lambda )^\mathrm{\#}\xi =h`$. Hence $`\phi i^{}(r(\lambda )^{\mathrm{\#}1}h)`$. Therefore, we conclude that $`h^{}i^{}(r(\lambda )^{\mathrm{\#}1}h)`$.
$`\mathrm{}`$
Remark In the proof above, the decomposition $`g=hm`$ and the choice of the basis $`\{e_1,\mathrm{},e_m\}`$ in (iv) depend on a particular point $`\lambda `$. It is not clear whether it is possible to find a decomposition so that Equation (7) holds uniformly for all points in $`h^{}`$. On the other hand, if there exists such a decomposition $`g=hm`$ so that a triangular dynamical $`r`$-matrix is of the form as in Equation (7), it is always splittable.
An immediate consequence of Proposition 2.7 is the following:
###### Corollary 2.9
If $`r:h^{}^2g`$ is a splittable triangular dynamical $`r`$-matrix, then
1. $`g_\lambda `$ is independent of $`\lambda `$, i.e., $`g_\lambda =g_\mu ,\lambda ,\mu h^{}`$. We will denote $`g_\lambda `$ by $`g_1`$.
2. $`r`$ can be considered as a non-degenerate triangular dynamical $`r`$-matrix valued in $`^2g_1`$.
Proof. By Proposition 2.7, $`Th^{}\times \{0\}`$ is a Lie subalgebroid of $`B`$. Hence for any $`XX(h^{}),(X,0)\mathrm{\Gamma }(B)`$. Let $`\phi _t`$ be the (local) flow on $`h^{}`$ generated by $`X`$. The bisection $`\mathrm{exp}t(X,0)`$ on the groupoid $`\mathrm{\Gamma }=h^{}\times h^{}\times G`$ generated by the section $`(X,0)\mathrm{\Gamma }(A)`$ is $`\{(\lambda ,\phi _t(\lambda ),1)|\lambda h^{}\}`$. Hence its induced isomorphism between $`\mathrm{\Gamma }_\lambda `$ and $`\mathrm{\Gamma }_{\phi _t(\lambda )}`$ is the identity map, when both of them are naturally identified with $`G`$. Here $`\mathrm{\Gamma }_\lambda `$ and $`\mathrm{\Gamma }_{\phi _t(\lambda )}`$ denote the isotropic groups of $`\mathrm{\Gamma }`$ at the points $`\lambda `$ and $`\phi _t(\lambda )`$, respectively. Therefore, $`Ad_{\mathrm{exp}t(X,0)}`$ is an identity map between their corresponding isotropic Lie algebras. On the other hand, since $`(X,0)\mathrm{\Gamma }(B)`$, hence $`Ad_{\mathrm{exp}t(X,0)}`$, when being restricted to $`B`$, is exactly the map which establishes the isomorphism between $`g_\lambda `$ and $`g_{\phi _t(\lambda )}`$. Hence, $`g_\lambda `$ and $`g_{\phi _t(\lambda )}`$ are equal as Lie subalgebras of $`g`$.
For the second part, since $`r`$ is splittable, we have $`r(\lambda )^\mathrm{\#}g^{}g_1`$ according to Proposition 2.7. Hence $`\lambda h^{},r(\lambda )^2(r(\lambda )^\mathrm{\#}g^{})^2g_1`$. By dimension counting, one easily sees that $`r`$ is non-degenerate when being considered as a dynamical $`r`$-matrix valued in $`^2g_1`$.
$`\mathrm{}`$
Let $`g:h^{}G^H`$ be a smooth map, where $`G^H`$ denotes the centralizer of $`H`$ in $`G`$ with its Lie algebra being denoted by $`g^H`$. Then $`g`$ can be naturally considered as a bisection of the groupoid $`\mathrm{\Gamma }=h^{}\times h^{}\times G`$, and hence we can talk about the induced automorphism $`Ad_g`$ of the corresponding Lie algebroid. In particular, we have a Gerstenhaber algebra automorphism $`Ad_g`$ on $`\mathrm{\Gamma }(^{}A)`$ .
Given a smooth function $`r:h^{}^2g`$, let $`\mathrm{\Lambda }_r=_ih_i\frac{}{\lambda ^i}+r(\lambda )\mathrm{\Gamma }(^2A)`$ as in Proposition 2.1. Then
$`Ad_g\mathrm{\Lambda }_r`$ $`=`$ $`Ad_g({\displaystyle \underset{i}{}}h_i{\displaystyle \frac{}{\lambda ^i}}+r)`$
$`=`$ $`{\displaystyle \underset{i}{}}Ad_gh_i({\displaystyle \frac{}{\lambda ^i}}{\displaystyle \frac{g}{\lambda ^i}}g^1)+Ad_gr`$
$`=`$ $`{\displaystyle \underset{i}{}}h_i({\displaystyle \frac{}{\lambda ^i}}{\displaystyle \frac{g}{\lambda ^i}}g^1)+Ad_gr`$
$`=`$ $`{\displaystyle \underset{i}{}}h_i{\displaystyle \frac{}{\lambda ^i}}+(Ad_gr{\displaystyle \underset{i}{}}h_i{\displaystyle \frac{g}{\lambda ^i}}g^1).`$
Here in the second from the last equality, we used $`Ad_gh_i=h_i`$ since $`gG^H`$. Let
$$r_g=Ad_gr\underset{i}{}h_i\frac{g}{\lambda ^i}g^1.$$
(8)
Combining with Proposition 2.1, we thus have proved the following:
###### Proposition 2.10
Assume that $`g:h^{}G^H`$ is a smooth map. Then
1. $`\mathrm{\Lambda }_{r_g}=Ad_g\mathrm{\Lambda }_r`$;
2. $`r`$ is a triangular dynamical $`r`$-matrix iff $`r_g`$ is a triangular dynamical $`r`$-matrix.
3. rank$`r_g`$=rank$`r`$; in particular, if $`r`$ is non-degenerate, so is $`r_g`$.
This proposition naturally leads us to the notion of gauge transformations on dynamical $`r`$-matrices, which was first introduced by Etingof and Varchenko . Recall that triangular dynamical $`r`$-matrices $`r_1`$ and $`r_2`$ are said to be gauge equivalent if there exists a smooth function $`g:h^{}G^H`$ such that $`r_2=(r_1)_g`$.
Remark Although non-degenerate triangular dynamical $`r`$-matrices are preserved by gauge transformations, splittable dynamical $`r`$-matrices in general are not. For example, the trivial triangular dynamical $`r`$-matrix $`r=0`$ is always splittable. However $`r_g=h_i\frac{g}{\lambda ^i}g^1`$ is never splittable unless $`G^H=H`$.
By $`(g,h)`$, we denote the quotient space of the space of all triangular dynamical $`r`$-matrices $`r:h^{}^2g`$ by gauge transformations, which is called the moduli space of triangular dynamical $`r`$-matrices.
Next we will introduce the dynamical $`r`$-matrix cohomology $`H_r^{}(g,h)`$, whose second cohomology group describes the tangent space of the moduli space $`(g,h)`$. As we will see in Section 6, the second cohomology group $`H_r^2(g,h)`$ is connected with the classification of quantizations of $`r`$ when it is non-degenerate.
Consider $`C^k=C^{\mathrm{}}(h^{},(^kg)^H)`$ (or equivalently denoted as $`C^{\mathrm{}}(h^{},(^kg)^h)`$), and define a differential $`\delta _r:C^kC^{k+1}`$ by
$$\delta _r\tau =\underset{i}{}h_i\frac{\tau }{\lambda ^i}+[r,\tau ],\tau C^k.$$
(9)
###### Proposition 2.11
$`\delta _r:C^kC^{k+1}`$ is well-defined and $`\delta _r^2=0`$.
Proof. It is clear that $`\delta _r\tau `$ is in $`C^{\mathrm{}}(h^{},(^{k+1}g)^H)`$ provided that $`\tau C^{\mathrm{}}(h^{},(^kg)^H)`$. For any $`\tau C^k=C^{\mathrm{}}(h^{},(^kg)^H)`$, $`\tau `$ can be naturally considered as a section of $`^kA`$, and
$`[\mathrm{\Lambda },\tau ]`$
$`=`$ $`[{\displaystyle \underset{i}{}}h_i{\displaystyle \frac{}{\lambda ^i}}+r,\tau ]`$
$`=`$ $`{\displaystyle \underset{i}{}}h_i{\displaystyle \frac{\tau }{\lambda ^i}}+[r,\tau ]`$
$`=`$ $`\delta _r\tau .`$
Since $`[\mathrm{\Lambda },\mathrm{\Lambda }]=0`$, it thus follows that $`\delta _r^2=0`$.
$`\mathrm{}`$
Hence the cochain complex $`\delta _r:C^kC^{k+1}`$ defines a cohomology, called the dynamical $`r`$-matrix cohomology, and denoted by $`H_r^{}(g,h)`$. Two remarks are in order.
Remark (1). The cochain complex $`\delta _r:C^kC^{k+1}`$ is in fact a subcomplex of the Lie algebroid cohomology cochain complex $`d_{}:\mathrm{\Gamma }(^kA)\mathrm{\Gamma }(^{k+1}A),d_{}X=[\mathrm{\Lambda },X]`$. Therefore it is easy to see that such a cochain complex is always defined for an arbitrary dynamical $`r`$-matrix, which is not necessary triangular.
(2). When $`r`$ is triangular, $`H_r^{}(g,h)`$ can be naturally identified with a “special” $`G\times H`$-invariant Poisson cohomology of the Poisson manifold $`(M,\pi )`$, i.e., the cohomology obtained by restricting the Poisson cochain complex to $`G\times H`$-invariant multi-vector fields tangent to the fibers of the fibration: $`h^{}\times Gh^{}`$.
###### Proposition 2.12
If $`g:h^{}G^H`$ is a smooth map, then
1. $`\delta _{r_g}Ad_g=Ad_g\delta _r`$;
2. $`Ad_g:(C^{},\delta _r)(C^{},\delta _{r_g})`$ induces an isomorphism $`H_r^{}(g,h)H_{r_g}^{}(g,h)`$.
Proof. For any $`\tau C^{\mathrm{}}(h^{},(^kg)^H)`$,
$`(Ad_g\delta _r)\tau `$ $`=`$ $`Ad_g[\mathrm{\Lambda },\tau ]`$
$`=`$ $`[Ad_g\mathrm{\Lambda },Ad_g\tau ]`$
$`=`$ $`[\mathrm{\Lambda }_{r_g},Ad_g\tau ]`$
$`=`$ $`(\delta _{r_g}Ad_g)\tau .`$
The conclusion thus follows immediately.
$`\mathrm{}`$
As a consequence, we conclude that $`H_r^{}(g,h)`$ only depends on the gauge equivalence class of the dynamical $`r`$-matrix. For this reason, we also denote this group by $`H_{[r]}^{}(g,h)`$.
###### Proposition 2.13
For any triangular dynamical $`r`$-matrix $`r:h^{}^2g`$, $`T_{[r]}(g,h)H_{[r]}^2(g,h)`$.
Proof. In Equation (2), replace $`r`$ by $`r+t\tau `$ and take the derivative at $`t=0`$, one obtains the linearization equation: $`_ih_i\frac{\tau }{\lambda ^i}+[r,\tau ]=0`$, i.e., $`\delta _r\tau =0`$. It is clear that $`\tau `$ is of zero weight since $`r+t\tau `$ is of zero weight.
To compute the tangent space to the gauge orbit at $`r`$, one needs to compute $`\frac{d}{dt}|_{t=0}(r_{\mathrm{exp}tf})`$, for $`fC^{\mathrm{}}(h^{},g^H)`$. Now $`r_{\mathrm{exp}tf}=Ad_{\mathrm{exp}tf}r_ih_i\frac{\mathrm{exp}tf}{\lambda ^i}(\mathrm{exp}tf)^1`$. It is thus simple to see that $`\frac{d}{dt}|_{t=0}(r_{\mathrm{exp}tf})=[f,r]_ih_i\frac{f}{\lambda ^i}=\delta _rf`$. The conclusion thus follows immediately.
$`\mathrm{}`$
Given a Lie algebra $`g`$, one may also consider classical triangular dynamical $`r`$-matrices $`r_{\mathrm{}}:h^{}(^2g)[[\mathrm{}]]`$ valued in $`g[[\mathrm{}]]`$ such that $`r_{\mathrm{}}(\lambda )=r(\lambda )+\mathrm{}r_1(\lambda )+\mathrm{}`$. The gauge transformation can be defined formally in an obvious way. Thus one can form the moduli space $`(g[[\mathrm{}]],h)`$. Assume that $`r:h^{}^2g`$ is a triangular classical dynamical $`r`$-matrix. From Proposition 2.13, it follows that $`T_{[r]}(g[[\mathrm{}]],h)H_{[r]}^2(g,h)[[\mathrm{}]]`$. By a formal neighbourhood of $`r`$ in $`(g[[\mathrm{}]],h)`$, denoted by $`_r(g[[\mathrm{}]],h)`$, we mean the subset in $`(g[[\mathrm{}]],h)`$ consisting of the classes of those elements $`r+O(\mathrm{})`$. Then $`H_{[r]}^2(g,h)[[\mathrm{}]]`$ can be considered as a linearization of $`_r(g[[\mathrm{}]],h)`$. In general, these two spaces are different. However, when $`r`$ is non-degenerate, they expect to be isomorphic, which should follow from Moser lemma.
In fact, as we will see in the next theorem, when $`r`$ is non-degenerate, $`H_{[r]}^{}(g,h)`$ is isomorphic to the relative Lie algebra cohomology.
###### Theorem 2.14
If $`r:h^{}^2g`$ is a non-degenerate dynamical $`r`$-matrix, then $`H_{[r]}^{}(g,h)`$ is isomorphic to $`H^{}(g,h)`$, the relative Lie algebra cohomology of the pair $`(g,h)`$.
Proof. Since $`r`$ is non-degenerate, $`(M,\pi )`$ is a symplectic manifold. As it is well known, $`\pi ^\mathrm{\#}:\mathrm{\Omega }^{}(M)X^{}(M)`$ induces an isomorphism between the de Rham cohomology cochain complex and the Poisson cohomology cochain complex. Now a $`k`$-mutivector field $`PX^k(M)`$ is in $`C^k`$ iff (i) $`P`$ is left $`G`$-invariant and right $`H`$-invariant; and (ii) $`d\lambda ^i\text{ }\text{ }P=0,i=1,\mathrm{},l`$. This, however, is equivalent to that (i) $`(\pi ^\mathrm{\#})^1P`$ is both left $`G`$-invariant and right $`H`$-invariant; and (ii) $`\stackrel{}{h_i}\text{ }\text{ }(\pi ^\mathrm{\#})^1P=0`$, because $`\pi ^\mathrm{\#}(d\lambda ^i)=\stackrel{}{h_i}`$, $`i=1,\mathrm{},l`$, and $`\pi `$ is $`G\times H`$-invariant. Note that a $`k`$-form $`\omega \mathrm{\Omega }^k(M)`$ is $`H`$-invariant and satisfies $`\stackrel{}{h_i}\text{ }\text{ }\omega =0,i=1,\mathrm{},l`$, iff $`\omega `$ is the pull back of a $`k`$-form on the quotient space $`M/H`$, i.e, $`\omega =p^{}\omega ^{}`$, where $`p:MM/H`$ is the projection and $`\omega ^{}\mathrm{\Omega }^k(M/H)`$. Moreover, $`\omega `$ is left $`G`$-invariant iff $`\omega ^{}`$ is left $`G`$-invariant since the left $`G`$-action on $`M`$ commutes with the right $`H`$-action. In summary, we have proved that the space $`(\pi ^\mathrm{\#})^1(C^k)`$ can be naturally identified with the space of left $`G`$-invariant $`k`$-forms on $`M/Hh^{}\times G/H`$. Under such an identification, the differential $`\delta _r`$ goes to the de-Rham differential. Hence $`H_{[r]}^k(g,h)`$ is isomorphic to the invariant de-Rham cohomology $`H^k(h^{}\times G/H)^G`$. Since $`G`$ does not act on the first factor $`h^{}`$, the latter is isomorphic to $`H^k(G/H)^G`$, which is in turn isomorphic to the relative Lie algebra cohomology $`H^k(g,h)`$ .
$`\mathrm{}`$
## 3 Quantization and star products
In this section, we investigate the relation between quantizations of a triangular dynamical $`r`$-matrix and star products on its associated Poisson manifold $`(M,\pi )`$. The main theme is to show that quantizing $`r`$ is equivalent to finding a certain special type of star products on $`M`$. Let us first introduce the precise definition of a quantization.
###### Definition 3.1
Let $`r:h^{}^2g`$ be a triangular dynamical $`r`$-matrix. A quantization of $`r`$ is an element $`F(\lambda )=1+\mathrm{}F_1(\lambda )+O(\mathrm{}^2)C^{\mathrm{}}(h^{},UgUg)[[\mathrm{}]]`$ satisfying
1. the zero weight condition: $`[1h+h1,F(\lambda )]=0,hh`$;
2. the shifted cocycle condition:
$$(\mathrm{\Delta }id)F(\lambda )F^{12}(\lambda \frac{1}{2}\mathrm{}h^{(3)})=(id\mathrm{\Delta })F(\lambda )F^{23}(\lambda +\frac{1}{2}\mathrm{}h^{(1)});$$
(10)
3. the normal condition:
$$(ϵid)F(\lambda )=1;(idϵ)F(\lambda )=1;\text{and}$$
(11)
4. the quantization condition: $`F_1^{12}(\lambda )F_1^{21}(\lambda )=r(\lambda )`$,
where $`\mathrm{\Delta }:UgUgUg`$ is the standard comultiplication, $`ϵ:UgC`$ is the counit map, and $`F^{12}(\lambda \frac{1}{2}\mathrm{}h^{(3)}),F^{23}(\lambda +\frac{1}{2}\mathrm{}h^{(1)})`$ are $`UgUgUg`$-valued functions on $`h^{}`$ defined by
$`F^{12}(\lambda {\displaystyle \frac{1}{2}}\mathrm{}h^{(3)})`$ $`=`$ $`F(\lambda )1{\displaystyle \frac{\mathrm{}}{2}}{\displaystyle \underset{i}{}}{\displaystyle \frac{F}{\lambda ^i}}h_i+{\displaystyle \frac{1}{2!}}({\displaystyle \frac{\mathrm{}}{2}})^2{\displaystyle \underset{i_1i_2}{}}{\displaystyle \frac{^2F}{\lambda ^{i_1}\lambda ^{i_2}}}h_{i_1}h_{i_2}`$ (12)
$`+\mathrm{}+{\displaystyle \frac{1}{k!}}({\displaystyle \frac{\mathrm{}}{2}})^k{\displaystyle \frac{^kF}{\lambda ^{i_1}\mathrm{}\lambda ^{i_k}}h_{i_1}\mathrm{}h_{i_k}}+\mathrm{},`$
and similarly for $`F^{23}(\lambda +\frac{1}{2}\mathrm{}h^{(1)})`$.
The relation between this definition of quantizations and the well known quantum dynamical Yang-Baxter equation (QDYBE) is explained by the following proposition, which can be proved by a straightforward verification.
###### Proposition 3.2
If $`F(\lambda )`$ is a quantization of a triangular dynamical $`r`$-matrix $`r(\lambda ):h^{}^2g`$, then $`R(\lambda )=F^{21}(\lambda )^1F^{12}(\lambda )`$ can be written as $`R(\lambda )=1+\mathrm{}r(\lambda )+O(\mathrm{}^2)`$ and satisfies the quantum dynamical Yang-Baxter equation (QDYBE):
$$R^{12}(\lambda \frac{1}{2}\mathrm{}h^{(3)})R^{13}(\lambda +\frac{1}{2}\mathrm{}h^{(2)})R^{23}(\lambda \frac{1}{2}\mathrm{}h^{(1)})=R^{23}(\lambda +\frac{1}{2}\mathrm{}h^{(1)})R^{13}(\lambda \frac{1}{2}\mathrm{}h^{(2)})R^{12}(\lambda +\frac{1}{2}\mathrm{}h^{(3)}).$$
(13)
Remark This is a symmetrized version of QDYBE, which is known to be equivalent to the non-symmetrized QDYBE:
$$R^{12}(\lambda +\mathrm{}h^{(3)})R^{13}(\lambda )R^{23}(\lambda +\mathrm{}h^{(1)})=R^{23}(\lambda )R^{13}(\lambda +\mathrm{}h^{(2)})R^{12}(\lambda ).$$
The reason for us to choose the symmetrized QDYBE in this paper is because it is related to the Weyl quantization, while the non-symmetrized QDYBE is related to the normal ordering quantization, as indicated in . Since we will use Fedosov method later on, the Weyl quantization is obviously of some advantage.
To proceed, we need some preparation on notations. Let $`𝒜=𝒟Ug[[\mathrm{}]]`$, where $`𝒟`$ is the algebra of smooth differential operators on $`h^{}`$. Then $`𝒟Ug`$ can be naturally identified with the algebra of left $`G`$-invariant differential operators on $`M`$. Hence $`𝒜`$ becomes a Hopf algebroid with base algebra $`R=C^{\mathrm{}}(h^{})[[\mathrm{}]]`$. The comultiplication
$$\mathrm{\Delta }:𝒜𝒜_R𝒜𝒟_{C^{\mathrm{}}(h^{})}𝒟UgUg[[\mathrm{}]]$$
is a natural extension of the comultiplications on $`𝒟`$ and on $`Ug`$:
$$\mathrm{\Delta }(Du)=\mathrm{\Delta }D\mathrm{\Delta }u,D𝒟,\text{ and }uUg,$$
where $`\mathrm{\Delta }D`$ is the bidifferential operator on $`h^{}`$ given by $`(\mathrm{\Delta }D)(f,g)=D(fg),f,gC^{\mathrm{}}(h^{})`$ and $`\mathrm{\Delta }uUgUg`$ is the usual comultiplication on $`Ug`$. Let us fix a basis in $`h`$, say $`\{h_1,\mathrm{},h_l\}`$, and let $`\{\xi _1,\mathrm{},\xi _l\}`$ be its dual basis, which in turn defines a coordinate system $`(\lambda ^1,\mathrm{},\lambda ^l)`$ on $`h^{}`$.
Set
$$\theta =\frac{1}{2}\underset{i=1}{\overset{l}{}}(h_i\frac{}{\lambda ^i}\frac{}{\lambda ^i}h_i)𝒜𝒜,\text{ and }\mathrm{\Theta }=\mathrm{exp}\mathrm{}\theta 𝒜𝒜.$$
(14)
Note that $`\theta `$, and hence $`\mathrm{\Theta }`$, is independent of the choice of a basis in $`h`$.
For each $`D𝒟Ug`$, we denote by $`\stackrel{}{D}`$ its corresponding left $`G`$-invariant differential operator on $`M=h^{}\times G`$. We also use a similar notation to denote multi-differential operators on $`M`$ as well. Now let $`r(\lambda ):h^{}^2g`$ be a triangular dynamical $`r`$-matrix, and $`M=h^{}\times G`$ its associated (regular) Poisson manifold with Poisson tensor $`\pi =_i\stackrel{}{h_i}\frac{}{\lambda ^i}+\stackrel{}{r(\lambda )}`$. It is simple to see that the Poisson brackets on $`C^{\mathrm{}}(M)`$ can be described as follows:
1. for any $`f,gC^{\mathrm{}}(h^{})`$, $`\{f,g\}=0`$;
2. for any $`fC^{\mathrm{}}(h^{})`$ and $`gC^{\mathrm{}}(G)`$, $`\{f,g\}=_i\frac{f}{\lambda ^i}\stackrel{}{h_i}g`$;
3. for any $`f,gC^{\mathrm{}}(G)`$, $`\{f,g\}=\stackrel{}{r(\lambda )}(f,g)`$.
This Poisson bracket relation naturally motivates the following theorem, which is indeed the main theorem of this section.
###### Theorem 3.3
Let $`(M,\pi )`$ be the Poisson manifold associated to a triangular dynamical $`r`$-matrix as in Proposition 2.2. Assume that $`_{\mathrm{}}`$ is a $`G\times H`$-invariant star product on $`(M,\pi )`$ satisfying the properties:
1. for any $`f,gC^{\mathrm{}}(h^{})`$,
$$f(\lambda )_{\mathrm{}}g(\lambda )=f(\lambda )g(\lambda );$$
2. for any $`f(\lambda )C^{\mathrm{}}(h^{})`$ and $`g(x)C^{\mathrm{}}(G)`$,
$`f(\lambda )_{\mathrm{}}g(x)`$ $`=`$ $`\stackrel{}{\mathrm{\Theta }}(f,g)={\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}({\displaystyle \frac{\mathrm{}}{2}})^k{\displaystyle \frac{1}{k!}}{\displaystyle \frac{^kf}{\lambda ^{i_1}\mathrm{}\lambda ^{i_k}}}\stackrel{}{h_{i_1}}\mathrm{}\stackrel{}{h_{i_k}}g,`$
$`g(x)_{\mathrm{}}f(\lambda )`$ $`=`$ $`\stackrel{}{\mathrm{\Theta }}(g,f)={\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}({\displaystyle \frac{\mathrm{}}{2}})^k{\displaystyle \frac{1}{k!}}\stackrel{}{h_{i_1}}\mathrm{}\stackrel{}{h_{i_k}}g{\displaystyle \frac{^kf}{\lambda ^{i_1}\mathrm{}\lambda ^{i_k}}};`$
3. there is a smooth map $`F:h^{}UgUg[[\mathrm{}]]`$ such that for any $`f(x),g(x)C^{\mathrm{}}(G)`$,
$$f_{\mathrm{}}g=\stackrel{}{F(\lambda )}(f,g).$$
(15)
Then $`F(\lambda )`$ is a quantization of the dynamical $`r`$-matrix $`r(\lambda )`$. Conversely, any quantization of $`r(\lambda )`$ corresponds to a $`G\times H`$-invariant star product on $`M`$ satisfying the properties (i)-(iii).
A $`G\times H`$-invariant star product on $`M`$ with properties (i)-(iii) is called a compatible star product. In other words, Theorem 3.3 can be stated that a quantization of $`r(\lambda )`$ is equivalent to a compatible star-product on $`M`$.
To prove Theorem 3.3, we need several lemmas.
###### Lemma 3.4
$`\mathrm{\Theta }`$ satisfies the equation:
$$[(\mathrm{\Delta }id)\mathrm{\Theta }]\mathrm{\Theta }^{12}=[(id\mathrm{\Delta })\mathrm{\Theta }]\mathrm{\Theta }^{23}\text{ in }𝒜𝒜𝒜.$$
(16)
Proof. Note that both sides of Equation (16) normally are elements in $`𝒜_R𝒜_R𝒜`$. In our situation, however, they indeed can be considered as elements in $`𝒜𝒜𝒜`$.
Now
$`[(\mathrm{\Delta }id)\mathrm{\Theta }]\mathrm{\Theta }^{12}`$
$`=`$ $`[(\mathrm{\Delta }id)\mathrm{exp}\mathrm{}\theta ]\mathrm{exp}\mathrm{}\theta ^{12}`$
$`=`$ $`\mathrm{exp}\mathrm{}[(\mathrm{\Delta }id)\theta +\theta ^{12}]`$
$`=`$ $`\mathrm{exp}{\displaystyle \frac{1}{2}}\mathrm{}{\displaystyle \underset{i=1}{\overset{k}{}}}(h_i1{\displaystyle \frac{}{\lambda ^i}}+1h_i{\displaystyle \frac{}{\lambda ^i}}+h_i{\displaystyle \frac{}{\lambda ^i}}1{\displaystyle \frac{}{\lambda ^i}}1h_i1{\displaystyle \frac{}{\lambda ^i}}h_i{\displaystyle \frac{}{\lambda ^i}}h_i1).`$
Here in the second equality we used the fact that $`(\mathrm{\Delta }id)\theta `$ and $`\theta ^{12}`$ commute in $`𝒜𝒜𝒜`$.
A similar computation leads to the same expression for $`[(id\mathrm{\Delta })\mathrm{\Theta }]\mathrm{\Theta }^{23}`$. This proves Equation (16).
$`\mathrm{}`$
###### Lemma 3.5
$`D_1,D_2,D_3𝒜`$, and $`f_1(\lambda )C^{\mathrm{}}(h^{})`$, $`f_2(x)C^{\mathrm{}}(G)`$, and $`g(\lambda ,x)C^{\mathrm{}}(h^{}\times G)`$,
$$\stackrel{}{[(\mathrm{\Delta }id)F(\lambda )(D_1D_2D_3)]}(f_1(\lambda ),f_2(x),g(\lambda ,x))=\stackrel{}{[F^{23}(\lambda )(D_1D_2D_3)]}(f_1(\lambda ),f_2(x),g(\lambda ,x)).$$
Proof. Write $`F(\lambda )=a_{\alpha \beta }(\lambda )u_\alpha u_\beta `$, with $`u_\alpha ,u_\beta Ug`$ and $`a_{\alpha \beta }(\lambda )C^{\mathrm{}}(h^{})[[\mathrm{}]]`$. Then
$$((\mathrm{\Delta }id)F(\lambda ))(D_1D_2D_3)=a_{\alpha \beta }(\lambda )\mathrm{\Delta }u_\alpha (D_1D_2)u_\beta D_3.$$
Hence
$`\stackrel{}{[(\mathrm{\Delta }id)F(\lambda )(D_1D_2D_3)]}(f_1(\lambda ),f_2(x),g(\lambda ,x))`$
$`=`$ $`{\displaystyle a_{\alpha \beta }(\lambda )\stackrel{}{\mathrm{\Delta }u_\alpha (D_1D_2)}(f_1(\lambda ),f_2(x))(\stackrel{}{u_\beta D_3}g)(\lambda ,x)}`$
$`=`$ $`{\displaystyle a_{\alpha \beta }(\lambda )\stackrel{}{u_\alpha }[(\stackrel{}{D_1}f_1)(\lambda )(\stackrel{}{D_2}f_2)(x)](\stackrel{}{u_\beta D_3}g)(\lambda ,x)}`$
$`=`$ $`{\displaystyle a_{\alpha \beta }(\lambda )(\stackrel{}{D_1}f_1)(\lambda )((\stackrel{}{u_\alpha D_2})f_2)(x)(\stackrel{}{u_\beta D_3}g)(\lambda ,x)}`$
$`=`$ $`\stackrel{}{D_1F(\lambda )(D_2D_3)}(f_1(\lambda ),f_2(x),g(\lambda ,x))`$
$`=`$ $`\stackrel{}{F^{23}(\lambda )(D_1D_2D_3)}(f_1(\lambda ),f_2(x),g(\lambda ,x)).`$
$`\mathrm{}`$
###### Corollary 3.6
$`f_1(\lambda )C^{\mathrm{}}(h^{}),f_2(x)C^{\mathrm{}}(G)`$ and $`g(\lambda ,x)C^{\mathrm{}}(h^{}\times G)`$,
$$\stackrel{}{F(\lambda )\mathrm{\Theta }}(f_1(\lambda )_{\mathrm{}}f_2(x),g(\lambda ,x))=\stackrel{}{\mathrm{\Theta }}(f_1(\lambda ),\stackrel{}{F(\lambda )\mathrm{\Theta }}(f_2(x),g(\lambda ,x))).$$
Proof.
$`\stackrel{}{F(\lambda )\mathrm{\Theta }}(f_1(\lambda )_{\mathrm{}}f_2(x),g(\lambda ,x))`$
$`=`$ $`\stackrel{}{F(\lambda )\mathrm{\Theta }}(\stackrel{}{\mathrm{\Theta }}(f_1(\lambda ),f_2(x)),g(\lambda ,x))`$
$`=`$ $`\stackrel{}{(\mathrm{\Delta }id)(F(\lambda )\mathrm{\Theta })\mathrm{\Theta }^{12}}(f_1(\lambda ),f_2(x),g(\lambda ,x))`$
$`=`$ $`\stackrel{}{(\mathrm{\Delta }id)F(\lambda )(\mathrm{\Delta }id)\mathrm{\Theta }\mathrm{\Theta }^{12}}(f_1(\lambda ),f_2(x),g(\lambda ,x))\text{ (by Lemma }\text{3.4}\text{)}`$
$`=`$ $`\stackrel{}{(\mathrm{\Delta }id)F(\lambda )(id\mathrm{\Delta })\mathrm{\Theta }\mathrm{\Theta }^{23}}(f_1(\lambda ),f_2(x),g(\lambda ,x))\text{ (by Lemma }\text{3.5}\text{)}`$
$`=`$ $`\stackrel{}{F^{23}(\lambda )(id\mathrm{\Delta })\mathrm{\Theta }\mathrm{\Theta }^{23}}(f_1(\lambda ),f_2(x),g(\lambda ,x)).`$
Let us write $`\mathrm{\Theta }=D_\alpha D_\beta `$. Then $`(id\mathrm{\Delta })\mathrm{\Theta }=D_\alpha \mathrm{\Delta }D_\beta `$, and
$`\stackrel{}{F^{23}(\lambda )(id\mathrm{\Delta })\mathrm{\Theta }\mathrm{\Theta }^{23}}(f_1(\lambda ),f_2(x),g(\lambda ,x))`$
$`=`$ $`{\displaystyle \stackrel{}{[D_\alpha F(\lambda )\mathrm{\Delta }D_\beta \mathrm{\Theta }]}(f_1(\lambda ),f_2(x),g(\lambda ,x))}`$
$`=`$ $`{\displaystyle (\stackrel{}{D_\alpha }f_1)(\lambda )\stackrel{}{F(\lambda )\mathrm{\Delta }D_\beta \mathrm{\Theta }}(f_2(x),g(\lambda ,x))}.`$
Using the expansion $`\mathrm{\Theta }=_{k=0}^{\mathrm{}}(\frac{\mathrm{}}{2})^k\frac{1}{k!}(_{i=1}^l(h_i\frac{}{\lambda _i}\frac{}{\lambda _i}h_i))^k`$, one obtains that
$`\stackrel{}{F(\lambda )\mathrm{\Theta }}(f_1(\lambda )_{\mathrm{}}f_2(x),g(\lambda ,x))`$
$`=`$ $`{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}({\displaystyle \frac{\mathrm{}}{2}})^k{\displaystyle \frac{1}{k!}}{\displaystyle \frac{^kf_1(\lambda )}{\lambda ^{i_1}\mathrm{}\lambda ^{i_k}}}\stackrel{}{F(\lambda )\mathrm{\Delta }(h_{i_1}\mathrm{}h_{i_k})\mathrm{\Theta }}(f_2(x),g(\lambda ,x))`$
$`=`$ $`{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}({\displaystyle \frac{\mathrm{}}{2}})^k{\displaystyle \frac{1}{k!}}{\displaystyle \frac{^kf_1(\lambda )}{\lambda ^{i_1}\mathrm{}\lambda ^{i_k}}}\stackrel{}{\mathrm{\Delta }(h_{i_1}\mathrm{}h_{i_k})F(\lambda )\mathrm{\Theta }}(f_2(x),g(\lambda ,x))`$
$`=`$ $`{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}({\displaystyle \frac{\mathrm{}}{2}})^k{\displaystyle \frac{1}{k!}}{\displaystyle \frac{^kf_1(\lambda )}{\lambda ^{i_1}\mathrm{}\lambda ^{i_k}}}\stackrel{}{h_{i_1}}\mathrm{}\stackrel{}{h_{i_k}}[\stackrel{}{F(\lambda )\mathrm{\Theta }}(f_2(x),g(\lambda ,x))]`$
$`=`$ $`\stackrel{}{\mathrm{\Theta }}(f_1(\lambda ),\stackrel{}{F(\lambda )\mathrm{\Theta }}(f_2(x),g(\lambda ,x))).`$
Here the second equality follows from the fact that $`F(\lambda )`$ is of zero weight, i.e., $`F(\lambda )(\mathrm{\Delta }h)=(\mathrm{\Delta }h)F(\lambda )`$, $`hh`$. This concludes the proof.
$`\mathrm{}`$
###### Proposition 3.7
Under the same hypothesis as in Theorem 3.3, we have
(1). for any $`f(\lambda )C^{\mathrm{}}(h^{})`$ and $`g(\lambda ,x)C^{\mathrm{}}(h^{}\times G)`$,
$`f(\lambda )_{\mathrm{}}g(\lambda ,x)`$ $`=`$ $`\stackrel{}{\mathrm{\Theta }}(f,g)={\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}({\displaystyle \frac{\mathrm{}}{2}})^k{\displaystyle \frac{1}{k!}}{\displaystyle \frac{^kf}{\lambda ^{i_1}\mathrm{}\lambda ^{i_k}}}\stackrel{}{h_{i_1}}\mathrm{}\stackrel{}{h_{i_k}}g,`$ (17)
$`g(\lambda ,x)_{\mathrm{}}f(\lambda )`$ $`=`$ $`\stackrel{}{\mathrm{\Theta }}(g,f)={\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}({\displaystyle \frac{\mathrm{}}{2}})^k{\displaystyle \frac{1}{k!}}\stackrel{}{h_{i_1}}\mathrm{}\stackrel{}{h_{i_k}}g{\displaystyle \frac{^kf}{\lambda ^{i_1}\mathrm{}\lambda ^{i_k}}};`$ (18)
(2). for any $`f(\lambda ,x)C^{\mathrm{}}(h^{}\times G)`$ and $`g(x)C^{\mathrm{}}(G)`$,
$`f(\lambda ,x)_{\mathrm{}}g(x)`$ $`=`$ $`(\stackrel{}{F(\lambda )\mathrm{\Theta }})(f,g)={\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}({\displaystyle \frac{\mathrm{}}{2}})^k{\displaystyle \frac{1}{k!}}\stackrel{}{F(\lambda )}({\displaystyle \frac{^kf}{\lambda ^{i_1}\mathrm{}\lambda ^{i_k}}},\stackrel{}{h_{i_1}}\mathrm{}\stackrel{}{h_{i_k}}g),`$ (19)
$`g(x)_{\mathrm{}}f(\lambda ,x)`$ $`=`$ $`(\stackrel{}{F(\lambda )\mathrm{\Theta }})(g,f)={\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}({\displaystyle \frac{\mathrm{}}{2}})^k{\displaystyle \frac{1}{k!}}\stackrel{}{F(\lambda )}(\stackrel{}{h_{i_1}}\mathrm{}\stackrel{}{h_{i_k}}g,{\displaystyle \frac{^kf}{\lambda ^{i_1}\mathrm{}\lambda ^{i_k}}}).`$ (20)
Proof. We will prove Equation (17) first. For that, it suffices to show this for $`g(\lambda ,x)=g_1(\lambda )_{\mathrm{}}g_2(x)`$, $`g_1(\lambda )C^{\mathrm{}}(h^{})`$ and $`g_2(x)C^{\mathrm{}}(G)`$, since, at each point, the $`C^{\mathrm{}}`$-jet space of $`C^{\mathrm{}}(h^{}\times G)[[\mathrm{}]]`$ is spanned by the $`C^{\mathrm{}}`$-jets of this type of functions. Now
$`f(\lambda )_{\mathrm{}}g(\lambda ,x)`$ $`=`$ $`f(\lambda )_{\mathrm{}}(g_1(\lambda )_{\mathrm{}}g_2(x))`$
$`=`$ $`(f(\lambda )_{\mathrm{}}g_1(\lambda ))_{\mathrm{}}g_2(x)`$
$`=`$ $`(f(\lambda )g_1(\lambda ))_{\mathrm{}}g_2(x)`$
$`=`$ $`\stackrel{}{\mathrm{\Theta }}(f(\lambda )g_1(\lambda ),g_2(x))`$
$`=`$ $`\stackrel{}{\mathrm{\Theta }}(\stackrel{}{\mathrm{\Theta }}(f(\lambda ),g_1(\lambda )),g_2(x))`$
$`=`$ $`\stackrel{}{[(\mathrm{\Delta }id)\mathrm{\Theta }]\mathrm{\Theta }^{12}}(f(\lambda ),g_1(\lambda ),g_2(x))\text{ (by Lemma }\text{3.4}\text{)}`$
$`=`$ $`\stackrel{}{[(id\mathrm{\Delta })\mathrm{\Theta }]\mathrm{\Theta }^{23}}(f(\lambda ),g_1(\lambda ),g_2(x))`$
$`=`$ $`\stackrel{}{\mathrm{\Theta }}(f(\lambda ),\stackrel{}{\mathrm{\Theta }}(g_1(\lambda ),g_2(x)))`$
$`=`$ $`\stackrel{}{\mathrm{\Theta }}(f(\lambda ),g_1(\lambda )_{\mathrm{}}g_2(x))`$
$`=`$ $`\stackrel{}{\mathrm{\Theta }}(f(\lambda ),g(\lambda ,x)).`$
Equation (18) can be proved similarly.
To prove Equation (19), similarly we may assume that $`f(\lambda ,x)=f_1(\lambda )_{\mathrm{}}f_2(x)`$, for $`f_1(\lambda )C^{\mathrm{}}(h^{})`$ and $`f_2(x)C^{\mathrm{}}(G)`$. Then
$`f(\lambda ,x)_{\mathrm{}}g(x)`$ $`=`$ $`(f_1(\lambda )_{\mathrm{}}f_2(x))_{\mathrm{}}g(x)`$
$`=`$ $`f_1(\lambda )_{\mathrm{}}(f_2(x)_{\mathrm{}}g(x))\text{ (using Equation (}\text{17}\text{))}`$
$`=`$ $`\stackrel{}{\mathrm{\Theta }}(f_1(\lambda ),f_2(x)_{\mathrm{}}g(x))`$
$`=`$ $`\stackrel{}{\mathrm{\Theta }}(f_1(\lambda ),\stackrel{}{F(\lambda )}(f_2(x),g(x)))`$
$`=`$ $`\stackrel{}{\mathrm{\Theta }}(f_1(\lambda ),\stackrel{}{F(\lambda )\mathrm{\Theta }}(f_2(x),g(x)))\text{ (by Corollary }\text{3.6}\text{)}`$
$`=`$ $`\stackrel{}{F(\lambda )\mathrm{\Theta }}(f_1(\lambda )_{\mathrm{}}f_2(x),g(x))`$
$`=`$ $`\stackrel{}{F(\lambda )\mathrm{\Theta }}(f(\lambda ,x),g(x)).`$
Equation (20) can also be proved similarly.
$`\mathrm{}`$
We are now ready to prove the main theorem of the section.
###### Theorem 3.8
Under the same hypothesis as in Theorem 3.3, $`\stackrel{}{F(\lambda )\mathrm{\Theta }}`$ is the formal bidifferential operator defining the star product $`_{\mathrm{}}`$, i.e., for any $`f(\lambda ,x),g(\lambda ,x)C^{\mathrm{}}(h^{}\times G)`$,
$$f(\lambda ,x)_{\mathrm{}}g(\lambda ,x)=\stackrel{}{F(\lambda )\mathrm{\Theta }}(f,g).$$
Proof. We may assume that $`f(\lambda ,x)=f_1(\lambda )_{\mathrm{}}f_2(x)`$, for $`f_1(\lambda )C^{\mathrm{}}(h^{})`$ and $`f_2(x)C^{\mathrm{}}(G)`$. Then
$`f(\lambda ,x)_{\mathrm{}}g(\lambda ,x)`$
$`=`$ $`(f_1(\lambda )_{\mathrm{}}f_2(x))_{\mathrm{}}g(\lambda ,x)`$
$`=`$ $`f_1(\lambda )_{\mathrm{}}(f_2(x)_{\mathrm{}}g(\lambda ,x))\text{ (by Proposition }\text{3.7}\text{)}`$
$`=`$ $`\stackrel{}{\mathrm{\Theta }}(f_1(\lambda ),\stackrel{}{F(\lambda )\mathrm{\Theta }}(f_2(x),g(\lambda ,x)))\text{ (by Corollary }\text{3.6}\text{)}`$
$`=`$ $`\stackrel{}{F(\lambda )\mathrm{\Theta }}(f_1(\lambda )_{\mathrm{}}f_2(x),g(\lambda ,x))`$
$`=`$ $`\stackrel{}{F(\lambda )\mathrm{\Theta }}(f(\lambda ,x),g(\lambda ,x)).`$
This concludes the proof.
$`\mathrm{}`$
Finally, before proving Theorem 3.3, we need the following result, which connects the shifted cocycle condition with the associativity of a star-product.
###### Proposition 3.9
Under the same hypothesis as in Theorem 3.3, $`f_1(x),f_2(x),f_3(x)C^{\mathrm{}}(G)`$,
1. $`\stackrel{}{(\mathrm{\Delta }id)F(\lambda )F^{12}(\lambda \frac{1}{2}\mathrm{}h^{(3)})}(f_1(x),f_2(x),f_3(x))=(f_1(x)_{\mathrm{}}f_2(x))_{\mathrm{}}f_3(x)`$;
2. $`\stackrel{}{(id\mathrm{\Delta })F(\lambda )F^{23}(\lambda +\frac{1}{2}\mathrm{}h^{(1)})}(f_1(x),f_2(x),f_3(x))=f_1(x)_{\mathrm{}}(f_2(x)_{\mathrm{}}f_3(x))`$.
Proof. From Equation (12), it follows that
$$(\mathrm{\Delta }id)F(\lambda )F^{12}(\lambda \frac{1}{2}\mathrm{}h^{(3)})=\frac{1}{k!}(\frac{\mathrm{}}{2})^k[(\mathrm{\Delta }id)F(\lambda )](\frac{^kF}{\lambda ^{i_1}\mathrm{}\lambda ^{i_k}}h_{i_1}\mathrm{}h_{i_k}).$$
Hence
$`\stackrel{}{(\mathrm{\Delta }id)F(\lambda )F^{12}(\lambda {\displaystyle \frac{1}{2}}\mathrm{}h^{(3)})}(f_1(x),f_2(x),f_3(x))`$
$`=`$ $`{\displaystyle \frac{1}{k!}(\frac{\mathrm{}}{2})^k\stackrel{}{F(\lambda )}[\stackrel{}{\frac{^kF(\lambda )}{\lambda ^{i_1}\mathrm{}\lambda ^{i_k}}}(f_1(x),f_2(x)),(\stackrel{}{h_{i_1}}\mathrm{}\stackrel{}{h_{i_k}}f_3)(x)]}`$
$`=`$ $`{\displaystyle \frac{1}{k!}(\frac{\mathrm{}}{2})^k\stackrel{}{F(\lambda )}[\frac{^k(f_1_{\mathrm{}}f_2)}{\lambda ^{i_1}\mathrm{}\lambda ^{i_k}},(\stackrel{}{h_{i_1}}\mathrm{}\stackrel{}{h_{i_k}}f_3)(x)](\text{using Equation (}\text{19}\text{)})}`$
$`=`$ $`(f_1(x)_{\mathrm{}}f_2(x))_{\mathrm{}}f_3(x).`$
The second identity can be proved similarly.
$`\mathrm{}`$
Proof of Theorem 3.3 Since $`_{\mathrm{}}`$ is invariant under the right $`H`$-action, $`\stackrel{}{F(\lambda )}`$ is right $`H`$-invariant. This implies that $`F(\lambda )`$ is $`Ad_H`$-invariant, and therefore is of zero weight. The normal condition follows from the fact that $`1`$ is the unit of the star algebra, i.e., $`1_{\mathrm{}}f=f_{\mathrm{}}1=f`$. And the shifted cocycle condition follows from the associativity of the star product together with Proposition 3.9. Finally, let us write $`F(\lambda )=1+\mathrm{}F_1(\lambda )+O(\mathrm{}^2)`$. Since $`_{\mathrm{}}`$ is a star product quantizing $`\pi `$, it follows that $`\stackrel{}{(F_1(\lambda )F_1^{21}(\lambda ))}(f,g)=\{f,g\}=\stackrel{}{r(\lambda )}(f,g),f,gC^{\mathrm{}}(G)`$. Hence it follows that $`F_1(\lambda )F_1^{21}(\lambda )=r(\lambda )`$.
Conversely, if $`F(\lambda )`$ is a quantization of $`r(\lambda )`$, according to Theorem 7.5 in , $`\stackrel{}{F(\lambda )\mathrm{\Theta }}`$ is indeed an associator and therefore defines a star product on $`M=h^{}\times G`$. It is simple to see that this star product is a quantization of $`\pi `$ and satisfies Properties (i)-(iii) in Theorem 3.3.
$`\mathrm{}`$
We end this section by the following
Remark Bordemann et. al. found an explicit formula for a star-product on $`R\times SU(2)`$ using a quantum analogue of Marsden-Weinstein reduction. It would be interesting to see if this is a compatible star-product.
## 4 Symplectic connections
From now on, we will confine ourselves mostly to non-degenerate triangular dynamical $`r`$-matrices. In this case, the corresponding Poisson manifolds are in fact symplectic, and therefore can be quantized by Fedosov method . As is well known, Fedosov quantization relies on the choice of a symplectic connection. Serving as a preliminary, this section is devoted to the discussion on symplectic connections. We will start with some general notations and constructions.
Let $``$ be a torsion-free symplectic connection on a symplectic manifold $`(M,\omega )`$. Define the symplectic curvature by
$$R(X,Y,Z,W)=\omega (X,R(Z,W)Y),X,Y,Z,WX(M),$$
(21)
where $`R(Z,W)Y=_Z_WY_W_ZY_{[Z,W]}Y`$ is the usual curvature tensor of $``$.
###### Proposition 4.1
1. $`R(X,Y,Z,W)`$ is skew symmetric with respect to $`Z`$ and $`W`$, and symmetric with respect to $`X`$ and $`Y`$, i.e.,
$$R(X,Y,Z,W)=R(X,Y,W,Z),R(X,Y,Z,W)=R(Y,X,Z,W).$$
(22)
2. The following Bianchi’s identity holds:
$$R(X,Y,Z,W)+R(X,Z,W,Y)+R(X,W,Y,Z)=0.$$
(23)
Proof. It is clear by definition that $`R(X,Y,Z,W)`$ is skew symmetric with respect to $`Z`$ and $`W`$. Now since $``$ is a symplectic connection, then
$`\omega (X,_Z_WY)`$
$`=`$ $`Z(\omega (X,_WY))\omega (_ZX,_WY)`$
$`=`$ $`Z(W\omega (X,Y))Z\omega (_WX,Y)W\omega (_ZX,Y)+\omega (_W_ZX,Y).`$
Similarly,
$$\omega (X,_W_ZY)=W(Z\omega (X,Y))W\omega (_ZX,Y)Z\omega (_WX,Y)+\omega (_Z_WX,Y).$$
Hence
$$\omega (X,_{[Z,W]}Y)=[Z,W](\omega (X,Y))\omega (_{[Z,W]}X,Y).$$
Thus
$`R(X,Y,Z,W)`$ $`=`$ $`\omega (X,R(Z,W)Y)`$
$`=`$ $`\omega (X,_Z_WY_W_ZY_{[Z,W]}Y)`$
$`=`$ $`\omega (_Z_WX_W_ZX_{[Z,W]}X,Y)`$
$`=`$ $`\omega (Y,R(Z,W)X)`$
$`=`$ $`R(Y,X,Z,W).`$
This concludes the proof of (i). Finally, (ii) follows from the usual Bianchi’s identity for a torsion-free connection.
$`\mathrm{}`$
Symplectic connections always exist on any symplectic manifold. In fact, there is a standard procedure to construct a torsion-free symplectic connection from an arbitrary torsion-free linear connection . Since such a construction is essential to our discussion here, let us recall it briefly below.
Assume that $`^0`$ is a torsion-free linear connection on a symplectic manifold $`M`$. Then any linear connection on $`M`$ can be written as
$$_XY=_X^0Y+S(X,Y),X,YX(M),$$
(24)
where $`S`$ is a $`(2,1)`$-tensor on $`M`$. Clearly, $``$ is torsion-free iff $`S`$ is symmetric, i.e., $`S(X,Y)=S(Y,X)`$, $`X,YX(M)`$. And $``$ is symplectic iff $`_X\omega =0`$. The latter is equivalent to
$$\omega (S(X,Y),Z)\omega (S(X,Z),Y)=(_X^0\omega )(Y,Z),X,Y,ZX(M).$$
(25)
###### Lemma 4.2
If $`^0`$ is a torsion-free linear connection, and $`S`$ is a $`(2,1)`$-tensor defined by the equation:
$$\omega (S(X,Y),Z)=\frac{1}{3}[(_X^0\omega )(Y,Z)+(_Y^0\omega )(X,Z)],$$
(26)
then $`_XY=_X^0Y+S(X,Y)`$ is a torsion-free symplectic connection. Moreover, if $`M`$ is a symplectic $`G`$-space and $`^0`$ is a $`G`$-invariant connection, then $``$ is also $`G`$-invariant.
Proof. Clearly, $`S(X,Y)`$, defined in this way, is symmetric with respect to $`X`$ and $`Y`$. Now
$`\omega (S(X,Y),Z)\omega (S(X,Z),Y)`$
$`=`$ $`{\displaystyle \frac{1}{3}}[(_X^0\omega )(Y,Z)+(_Y^0\omega )(X,Z)]{\displaystyle \frac{1}{3}}[(_X^0\omega )(Z,Y)+(_Z^0\omega )(X,Y)]`$
$`=`$ $`{\displaystyle \frac{1}{3}}[(_X^0\omega )(Y,Z)+(_Y^0\omega )(X,Z)+(_X^0\omega )(Y,Z)+(_Z^0\omega )(Y,X)]`$
$`=`$ $`(_X^0\omega )(Y,Z),`$
where the last step follows from the identity:
$$(_X^0\omega )(Y,Z)+(_Y^0\omega )(Z,X)+(_Z^0\omega )(X,Y)=0.$$
This means that $``$ is a torsion-free symplectic connection. The second statement is obvious according to Equation (26).
$`\mathrm{}`$
Now we retain to the case that $`M=h^{}\times G`$, the symplectic manifold associated with a non-degenerate triangular dynamical $`r`$-matrix $`r`$, which is our main subject of interest in the present paper. The main result is the following
###### Theorem 4.3
Assume that $`r:h^{}^2g`$ is a non-degenerate triangular dynamical $`r`$-matrix. Let $`M=h^{}\times G`$ be equipped with the symplectic structure as in Corollary 2.5. Then $`M`$ admits a $`G\times H`$-invariant torsion-free symplectic connection $``$ satisfying the property that $`_X\stackrel{}{h}=0,XX(M),hh`$.
We need a couple of lemmas first.
###### Lemma 4.4
Assume that $`g`$ admits a reductive decomposition $`g=hm`$, i.e., $`[h,m]m`$. Then, the following equations define a biinvariant torsion-free linear connection $`^0`$ on $`M`$:
$`\begin{array}{ccc}_X^0\frac{}{\lambda ^i}=0,\hfill & _X^0\stackrel{}{h}=0,\hfill & _X^0\stackrel{}{e}=0;\hfill \\ _\stackrel{}{h}^0\frac{}{\lambda ^i}=0,\hfill & _\stackrel{}{h}^0\stackrel{}{h_1}=0,\hfill & _\stackrel{}{h}^0\stackrel{}{e}=\stackrel{}{[h,e]};\hfill \\ _\stackrel{}{e}^0\frac{}{\lambda ^i}=0,\hfill & _\stackrel{}{e}^0\stackrel{}{h}=0,\hfill & _{\stackrel{}{e_1}}^0\stackrel{}{e_2}=\frac{1}{2}\stackrel{}{[e_1,e_2]},\hfill \end{array}`$ (30)
where $`XX(h^{})`$, $`h,h_1h`$, and $`e,e_1,e_2m`$.
Proof. This follows from a straightforward verification.
$`\mathrm{}`$
###### Lemma 4.5
Given a Lie algebra $`g`$, if there exists a non-degenerate triangular dynamical $`r`$-matrix $`r:h^{}^2g`$, then $`g`$ admits a reductive decomposition $`g=hm`$ so that $`[h,m]m`$.
Proof. Fixing any $`\lambda h^{}`$, we take $`m=r(\lambda )^\mathrm{\#}h^{}`$. Since $`r(\lambda )`$ is non-degenerate, by definition, we have $`g=h+m`$. On the other hand, it is clear that $`\text{dim}m\text{dim}h^{}=\text{dim}g\text{dim}h`$. Hence, $`\text{dim}h+\text{dim}m\text{dim}g`$. Therefore $`g=h+m`$ must be a direct sum. For any $`hh`$ and $`\xi g^{}`$, since $`r(\lambda )`$ is of zero weight, we have $`[h,r(\lambda )^\mathrm{\#}\xi ]=r(\lambda )^\mathrm{\#}(ad_h^{}\xi )`$. Since $`ad_h^{}\xi h^{}`$ for any $`\xi g^{}`$, it follows that $`m=r(\lambda )^\mathrm{\#}h^{}`$ is stable under the adjoint action of $`h`$.
$`\mathrm{}`$
Remark Note that, in our proof above, the decomposition $`g=hm`$ depends on the choice of a particular point $`\lambda h^{}`$. It is not clear if $`m=r(\lambda )^\mathrm{\#}h^{}`$ is independent of $`\lambda `$.
Proof of Theorem 4.3 According to Lemma 4.5, we may find a reductive decomposition $`g=hm`$ such that $`[h,m]m`$. Let $`^0`$ be the $`G`$-biinvariant torsion-free connection on $`M`$ as in Lemma 4.4. According to Lemma 4.2, one can construct a torsion-free symplectic connection $``$ on $`M`$. Since the symplectic structure is $`G\times H`$-invariant, the resulting symplectic connection $``$ is $`G\times H`$-invariant. It remains to show that $``$ still satisfies the condition that $`_X\stackrel{}{h}=0,XX(M)`$ and $`hh`$. The latter is equivalent to that $`S(\stackrel{}{h},X)=0`$. To show this identity, first note that $`XX(M)`$, $`_\stackrel{}{h}^0X=L_\stackrel{}{h}X`$, since $`^0`$ is torsion-free and $`_X^0\stackrel{}{h}=0`$. Hence $`_\stackrel{}{h}^0\omega =L_\stackrel{}{h}\omega `$. However, $`L_\stackrel{}{h}\omega =0`$ since $`\omega `$ is invariant under the right $`H`$-action. Thus, we have $`_\stackrel{}{h}^0\omega =0`$. According to Equation (26), $`YX(M)`$, $`\omega (S(\stackrel{}{h},X),Y)=\frac{1}{3}[(_\stackrel{}{h}^0\omega )(X,Y)+(_X^0\omega )(\stackrel{}{h},Y)]=\frac{1}{3}(_X^0\omega )(\stackrel{}{h},Y)`$. This implies that $`\omega ^b(S(\stackrel{}{h},X))=\frac{1}{3}(\stackrel{}{h}\text{ }\text{ }_X^0\omega )=\frac{1}{3}_X^0(\stackrel{}{h}\text{ }\text{ }\omega )`$ since $`_X^0\stackrel{}{h}=0`$. Finally, for any $`i`$, $`\stackrel{}{h_i}\text{ }\text{ }\omega =d\lambda ^i`$ and from the table in Lemma 4.4, it is easy to check that $`_X^0(d\lambda ^i)=0,i=1,\mathrm{},l`$. It thus follows that $`S(\stackrel{}{h_i},X)=0,i=1,\mathrm{},l`$. This concludes the proof.
$`\mathrm{}`$
In the case that $`r(\lambda )^2m`$, the symplectic connection can be described more explicitly.
###### Proposition 4.6
Suppose that $`g=hm`$ is a reductive decomposition, $`\{h_1,\mathrm{},h_l\}`$ is a basis of $`h`$, and $`\{e_1,\mathrm{},e_m\}`$ is a basis of $`m`$. Suppose that $`r(\lambda )=_{ij}r^{ij}(\lambda )e_ie_j`$ is a non-degenerate triangular dynamical $`r`$-matrix. Then the symplectic connection on $`M`$ obtained from $`^0`$, using the standard construction as in Lemma 4.4, has the following form:
$`\begin{array}{ccc}_{\frac{}{\lambda ^i}}\frac{}{\lambda ^j}=0,\hfill & _{\frac{}{\lambda ^i}}\stackrel{}{h_j}=0,\hfill & _{\frac{}{\lambda ^i}}\stackrel{}{e_j}=_kd_{ij}^k(\lambda )\stackrel{}{e_k};\hfill \\ _{\stackrel{}{h_i}}\frac{}{\lambda ^j}=0,\hfill & _{\stackrel{}{h_i}}\stackrel{}{h_j}=0,\hfill & _{\stackrel{}{h_i}}\stackrel{}{e_j}=\stackrel{}{[h_i,e_j]};\hfill \\ _{\stackrel{}{e_i}}\frac{}{\lambda ^j}=_kd_{ij}^k(\lambda )\stackrel{}{e_k},\hfill & _{\stackrel{}{e_i}}\stackrel{}{h_j}=0,\hfill & _{\stackrel{}{e_i}}\stackrel{}{e_j}=\frac{1}{2}\stackrel{}{[e_i,e_j]}+_kf_{ij}^k(\lambda )\stackrel{}{e_k},\hfill \end{array}`$ (34)
where $`d_{ij}^k(\lambda )`$ and $`f_{ij}^k(\lambda )`$ are smooth functions on $`h^{}`$.
Proof. The proof is essentially a straightforward computation. We omit it here.
$`\mathrm{}`$
###### Corollary 4.7
Under the same hypothesis as in Proposition 4.6, if $`\{h_{}^1,\mathrm{},h_{}^l,e_{}^1,\mathrm{},e_{}^m\}`$ denotes the dual basis of $`\{h_1,\mathrm{},h_l,e_1,\mathrm{},e_m\}`$, then
$`\begin{array}{ccc}_{\frac{}{\lambda ^i}}d\lambda ^j=0,\hfill & _{\frac{}{\lambda ^i}}\stackrel{}{h_{}^j}=0,\hfill & _{\frac{}{\lambda ^i}}\stackrel{}{e_{}^j}=_kd_{ik}^j(\lambda )\stackrel{}{e_{}^k};\hfill \\ _{\stackrel{}{h_i}}d\lambda ^j=0,\hfill & _{\stackrel{}{h_i}}\stackrel{}{h_{}^j}=0,\hfill & _{\stackrel{}{h_i}}\stackrel{}{e_{}^j}=\stackrel{}{ad_{h_i}^{}e_{}^j};\hfill \\ _{\stackrel{}{e_i}}d\lambda ^j=0,\hfill & _{\stackrel{}{e_i}}\stackrel{}{h_{}^j}=\frac{1}{2}_ka_{ik}^j\stackrel{}{e_{}^k},\hfill & _{\stackrel{}{e_i}}\stackrel{}{e_{}^j}=_kd_{ik}^j(\lambda )d\lambda ^k(\frac{1}{2}a_{ik}^j+f_{ik}^j(\lambda ))\stackrel{}{e_{}^k},\hfill \end{array}`$ (38)
where the coadjoint action is defined by $`<ad_u^{}\xi ,v>=<\xi ,[u,v]>`$, $`u,vg`$ and $`\xi g^{}`$, and the constants $`a_{ij}^k`$ are defined by the equation $`[e_i,e_j]=_ka_{ij}^kh_k(\text{mod }m)`$.
We end this section by generalizing Theorem 4.3 to the splittable triangular dynamical $`r`$-matrix case. According to Corollary 2.9, one may reduce a splittable triangular dynamical $`r`$-matrix to a non-degenerate one by considering the Lie subalgebra $`g_1g`$. Thus immediately we obtain the following
###### Corollary 4.8
Assume that $`r:h^{}^2g`$ is a splittable triangular dynamical $`r`$-matrix. Let $`M=h^{}\times G`$ be its associated Poisson manifold as in Proposition 2.2, which admits a (regular) symplectic foliation. Then there exists a $`G\times H`$-invariant torsion-free leafwise Poisson connection $``$ satisfying $`_X\stackrel{}{h}=0`$, for any $`hh`$ and any vector field $`XX(M)`$ tangent to the symplectic foliation.
However, when a triangular dynamical $`r`$-matrix $`r`$ is not splittable, such a Poisson connection may not exist. We give a counterexample below.
Example 4.9. Consider a two dimensional Lie algebra $`g`$ with basis $`\{h,e\}`$ satisfying the bracket relation $`[h,e]=ah`$, where $`a`$ is a fixed constant. Let $`h=Rh`$ and $`r(\lambda )=f(\lambda )he`$, where $`f(\lambda )`$ is a smooth function. It is simple to see that $`r(\lambda )`$ is a triangular dynamical $`r`$-matrix of rank zero, and it is not splittable unless $`a=0`$. Nevertheless, $`r(\lambda )`$ defines a regular rank 2 Poisson structure on the three dimensional space $`M=R\times G`$ with the Poisson tensor $`\pi =\stackrel{}{h}\frac{d}{d\lambda }+f(\lambda )\stackrel{}{h}\stackrel{}{e}`$, where $`G`$ is a 2-dimensional Lie group integrating the Lie algebra $`g`$. It is simple to see that the symplectic foliation of $`M`$ is spanned by the vector fields $`\stackrel{}{h}`$ and $`\frac{d}{d\lambda }+f(\lambda )\stackrel{}{e}`$. Let us denote $`X=\frac{d}{d\lambda }+f(\lambda )\stackrel{}{e}`$. Then, we have $`[\stackrel{}{h},X]=af(\lambda )\stackrel{}{h}`$. Now suppose that $``$ is a $`G\times H`$-invariant torsion-free leafwise Poisson connection on $`M`$ satisfying the condition that $`_\stackrel{}{h}\stackrel{}{h}=0`$ and $`_X\stackrel{}{h}=0`$. Since $``$ is torsion-free, it follows that $`_\stackrel{}{h}X=[\stackrel{}{h},X]=af(\lambda )\stackrel{}{h}`$. Assume that $`_XX=b(\lambda ,x)\stackrel{}{h}+c(\lambda ,x)X`$, where $`b(\lambda ,x)`$ and $`c(\lambda ,x)`$ are smooth functions on $`M`$. Then, $`_X\pi =_X(\stackrel{}{h}X)=\stackrel{}{h}_XX=c(\lambda ,x)\stackrel{}{h}X`$. Since $``$ is a Poisson connection, it follows that $`c(\lambda ,x)=0`$. Finally, we still need to check that $``$ is $`G\times H`$-invariant. It is clear that $``$ is $`G`$-invariant iff the function $`b(\lambda ,x)`$ is independent of $`xG`$ (which will be denoted by $`b(\lambda )`$). For it to be invariant under the right $`H`$-action, one needs the following condition:
$$_{[\stackrel{}{h},X]}X+_X[\stackrel{}{h},X]=[\stackrel{}{h},_XX]=[\stackrel{}{h},b(\lambda )\stackrel{}{h}]=0.$$
It thus follows that $`_{(af(\lambda )\stackrel{}{h})}X+_X(af(\lambda )\stackrel{}{h})=0`$, which implies that $`f^2(\lambda )a^2\stackrel{}{h}+a(\frac{df}{d\lambda })\stackrel{}{h}=0`$. Therefore, we arrive at the following equation (under the assumption that $`a0`$):
$$\frac{df}{d\lambda }=af^2(\lambda ).$$
(39)
In conclusion, we have proved that such a connection does not exist unless $`f(\lambda )`$ is a solution of the above equation. It would be interesting to find out what is the geometric meaning of this equation.
Remark Our quantization method does not work for this particular example. It is thus very natural to ask whether this dynamical $`r`$-matrix is still quantizable. Etingof and Nikshych recently has given an affimative answer to this question using the so called vertex-IRF transformation method . Their method indeed works for a large class of dynamical $`r`$-matrices called “completely degenerate”, which somehow is opposite to the non-degenerate ones considered in this paper. It would be very interesting to see whether one could combine these two methods together to completely solve the quantization problem for arbitary triangular dynamical $`r`$-matrices.
## 5 Compatible Fedosov star products
In this section, we consider Fedosov star products on a symplectic Hamiltonian $`H`$-space $`M`$, where $`H`$ is an Abelian group. For the reader’s convenience, we will give a brief account of the general construction of Fedosov star products in Appendix. Readers may refer to that section for various notations and formulas that are used here. What is eventually relevant to our situation is the case when $`M`$ is the symplectic manifold $`h^{}\times G`$ corresponding to a nondegenerate dynamical $`r`$-matrix. However, we believe that our general presentation would be of its own interest. We can now state the main result of this section.
###### Theorem 5.1
Let $`H`$ be an Abelian group and $`M`$ a symplectic Hamiltonian $`H`$-space with an equivariant momentum map $`J:Mh^{}`$. Assume that $`J`$ is a submersion, and there exists a $`H`$-invariant symplectic connection $``$ such that $`\stackrel{}{h}`$ is parallel for any $`hh`$, i.e., $`_X\stackrel{}{h}=0`$, $`XX(M)`$. Let $`_{\mathrm{}}`$ be the corresponding Fedosov star product on $`M`$ with Weyl curvature $`\mathrm{\Omega }=\omega +\mathrm{}\omega _1+\mathrm{}+\mathrm{}^i\omega _i+\mathrm{}Z^2(M)[[\mathrm{}]]`$, which satisfies the condition that $`i_\stackrel{}{h}\omega _i=0,i1,hh`$. Then for any $`f(\lambda )C^{\mathrm{}}(h^{})`$ and $`g(x)C^{\mathrm{}}(M)`$, we have
$`(J^{}f)_{\mathrm{}}g(x)`$ $`=`$ $`{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}({\displaystyle \frac{\mathrm{}}{2}})^k{\displaystyle \frac{1}{k!}}J^{}({\displaystyle \frac{^kf}{\lambda ^{i_1}\mathrm{}\lambda ^{i_k}}})\stackrel{}{h_{i_1}}\mathrm{}\stackrel{}{h_{i_k}}g;`$
$`g(x)_{\mathrm{}}(J^{}f)`$ $`=`$ $`{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}({\displaystyle \frac{\mathrm{}}{2}})^k{\displaystyle \frac{1}{k!}}(\stackrel{}{h_{i_1}}\mathrm{}\stackrel{}{h_{i_k}}g)J^{}({\displaystyle \frac{^kf}{\lambda ^{i_1}\mathrm{}\lambda ^{i_k}}}).`$
Here $`\stackrel{}{h}`$ denotes the corresponding Hamiltonian vector field on $`M`$ generated by $`hh`$.
Remark From Theorem 5.1, it follows that $`J^{}:C^{\mathrm{}}(h^{})[[\mathrm{}]]C^{\mathrm{}}(M)[[\mathrm{}]]`$ is an algebra homomorphism, where $`C^{\mathrm{}}(h^{})[[\mathrm{}]]`$ is equipped with pointwise multiplication. In other words, $`J^{}`$ is a quantum momentum map . It would be interesting to see how to generalize this result to the case when $`H`$ is not Abelian .
Applying Theorem 5.1 to the symplectic manifold $`M=h^{}\times G`$ associated to a nondegenerate triangular dynamical $`r`$-matrix, and using Theorem 4.3, we obtain the following
###### Corollary 5.2
Let $`r:h^{}^2g`$ be a nondegenerate triangular dynamical $`r`$-matrix, and $`M=h^{}\times G`$ its associated symplectic manifold. Let $``$ be the symplectic connection on $`M`$ as in Theorem 4.3. Suppose that $`\mathrm{\Omega }=\omega +\mathrm{}\omega _1+\mathrm{}+\mathrm{}^i\omega _i+\mathrm{}Z^2(M)^G[[\mathrm{}]]`$ satisfies the condition that $`i_\stackrel{}{h}\omega _i=0,i1,hh`$. Then the Fedosov star product on $`M`$ corresponding to $`(,\mathrm{\Omega })`$ is a compatible star product.
Combining with Theorem 3.3, we are lead to the following main result of the paper.
###### Theorem 5.3
Any nondegenerate triangular dynamical $`r`$-matrix is quantizable.
More generally, if $`r`$ is a splittable triangular dynamical $`r`$-matrix, according to Corollary 4.8, the corresponding Poisson manifold $`M=h^{}\times G`$ admits a $`G\times H`$-invariant leafwise (w.r.t. the symplectic foliation) Poisson connection such that $`_X\stackrel{}{h}=0,hh`$. Applying Theorem 5.1 leafwisely, we thus have the following
###### Theorem 5.4
Any splittable triangular dynamical $`r`$-matrix is quantizable.
The rest of the section is devoted to the proof of Theorem 5.1. We will start with the following
###### Proposition 5.5
Under the same hypothesis as in Theorem 5.1, we have
1. For any $`(r,s)`$-type tensor $`S𝒯^{(r,s)}M`$ and $`hh`$, we have $`_\stackrel{}{h}S=L_\stackrel{}{h}S`$.
2. $`XX(M)`$ and $`i1`$, $`_X(J^{}d\lambda ^i)=0`$.
3. Given any $`\theta \mathrm{\Omega }^1(M)`$, if $`\stackrel{}{h}\text{ }\text{ }\theta =0`$, then $`\stackrel{}{h}\text{ }\text{ }_X\theta =0`$, $`XX(M)`$.
4. $`R(X,Y,Z,W)=0`$, if any of the vectors $`X,Y,Z,W`$ is tangent to the $`H`$-orbits.
5. $`_\stackrel{}{h}R=0`$, $`hh`$.
Proof. (i). Since $``$ is torsion-free, for any vector field $`XX(M)`$, we have
$$_\stackrel{}{h}X=_X\stackrel{}{h}+[\stackrel{}{h},X]=[\stackrel{}{h},X]=L_\stackrel{}{h}X.$$
This implies that $`_\stackrel{}{h}\theta =L_\stackrel{}{h}\theta `$ for any one form $`\theta \mathrm{\Omega }^1(M)`$. Therefore $`_\stackrel{}{h}S=L_\stackrel{}{h}S`$ for any $`(r,s)`$-type tensor $`S𝒯^{(r,s)}M`$.
(ii). Since $`J:Mh^{}`$ is a momentum map, it follows that $`J^{}d\lambda ^i=\omega ^b\stackrel{}{h_i}`$, where $`\omega ^b:X(M)\mathrm{\Omega }^1(M)`$ is the isomorphism induced by the symplectic structure $`\omega `$. Hence $`_X(J^{}d\lambda ^i)=_X(\omega ^b\stackrel{}{h_i})=\omega ^b(_X\stackrel{}{h_i})=0`$, since $``$ is a symplectic connection.
(iii). We have $`_X(\stackrel{}{h}\text{ }\text{ }\theta )=(_X\stackrel{}{h})\text{ }\text{ }\theta +\stackrel{}{h}\text{ }\text{ }_X\theta =\stackrel{}{h}\text{ }\text{ }_X\theta `$. The claim thus follows.
(iv). Let $`\mathrm{\Phi }`$ denote the $`H`$-action on $`M`$. For any $`hh`$, since $``$ is $`H`$-invariant, it follows that $`W,YX(M)`$, $`_{(\mathrm{\Phi }_{\mathrm{exp}th}W)}(\mathrm{\Phi }_{\mathrm{exp}th}Y)=\mathrm{\Phi }_{\mathrm{exp}th}(_WY)`$. Taking the derivative at $`t=0`$, one obtains that
$$_{[\stackrel{}{h},W]}Y+_W[\stackrel{}{h},Y]=[\stackrel{}{h},_WY].$$
Hence,
$`R(\stackrel{}{h},W)Y`$ $`=`$ $`_\stackrel{}{h}_WY_W_\stackrel{}{h}Y_{[\stackrel{}{h},W]}Y`$
$`=`$ $`[\stackrel{}{h},_WY]_W[\stackrel{}{h},Y]_{[\stackrel{}{h},W]}Y`$
$`=`$ $`0.`$
On the other hand, we know that $`R(Z,W)\stackrel{}{h}=0`$, since $`\stackrel{}{h}`$ is parallel by assumption. This means that $`R(X,Y,Z,W)=0`$ if $`Y=\stackrel{}{h}`$ or $`Z=\stackrel{}{h}`$. Since $`R(X,Y,Z,W)`$ is antisymmetric with respect to $`W,Z`$, and symmetric with respect to $`X,Y`$ according to Proposition 4.1, the conclusion thus follows.
(v). Since both the connection $``$ and the symplectic structure $`\omega `$ are $`H`$-invariant, the symplectic curvature $`R`$, as defined by Equation (21), is also $`H`$-invariant. Hence, for any $`hh`$, according to (i), $`_\stackrel{}{h}R=L_\stackrel{}{h}R=0`$.
This completes the proof of the proposition.
$`\mathrm{}`$
By $`KTM`$, we denote the integrable distribution on $`M`$ corresponding to the $`H`$-orbits, and $`K^{}`$ its conormal subbundle. That is, a covector $`\theta `$ is in $`K^{}`$ iff $`<\theta ,\stackrel{}{h}>=0,hh`$. For any $`xM`$, by $`\text{pol}(K_x^{})`$, we denote the polynomials on $`T_xM`$ generated by those linear functions corresponding to covectors in $`K_x^{}`$. By $`W_x^{}`$, we denote the formal power series in $`\mathrm{}`$ with coefficients in $`\text{pol}(K_x^{})`$. Clearly $`W_x^{}`$ is a subalgebra of the Weyl algebra $`W_x`$. Let $`W^{}=_{xM}W_x^{}`$ be the subbundle of $`W`$. We also consider $`W^{}^qK^{}`$, a subbundle of $`W^qT^{}M`$, whose space of sections is denoted by $`\mathrm{\Gamma }W^{}(\mathrm{\Lambda }^{})^q`$. As before, let us fix a basis $`\{h_1,\mathrm{},h_l\}`$ of $`h`$, and denote by $`(\lambda ^1,\mathrm{},\lambda ^l)`$ its induced coordinate system on $`h^{}`$. Since $`J:Mh^{}`$ is a momentum map, we have $`X_{J^{}\lambda ^i}=\stackrel{}{h_i},i=1,\mathrm{},l`$. It thus follows that $`J_{}\stackrel{}{h_i}=J_{}X_{J^{}\lambda ^i}=0,i=1,\mathrm{},l`$, since $`h`$ is Abelian. Next we need to extend $`\{\stackrel{}{h_1},\mathrm{},\stackrel{}{h_l}\}`$ to a set of (local) vector fields which constitutes a basis of tangent fibers of $`M`$. For this purpose, let $`\{u_1,\mathrm{},u_m\}`$ be (local) vector fields on $`M`$ tangent to the J-fibers such that $`\{\stackrel{}{h_1},\mathrm{},\stackrel{}{h_l},u_1,\mathrm{},u_m\}`$ constitutes a basis of the tangent spaces of the J-fibers. Choose (local) vector fields $`\{v_1,\mathrm{},v_l\}`$ on $`M`$ such that $`J_{}v_i=\frac{}{\lambda ^i},i=1,\mathrm{},l`$, which is always possible since $`J`$ is a submersion. It is easy to see that locally $`\{\stackrel{}{h_1},\mathrm{},\stackrel{}{h_l},v_1,\mathrm{},v_l,u_1,\mathrm{},u_m\}`$ constitutes a basis of the tangent fibers of $`M`$. Let $`\{\stackrel{}{h_{}^1},\mathrm{},\stackrel{}{h_{}^l},v_{}^1,\mathrm{}v_{}^l,u_{}^1,\mathrm{},u_{}^m\}`$ be its dual basis. Then any section of $`W\mathrm{\Lambda }`$ can be written as
$$a=\mathrm{}^ka_{k,i_1\mathrm{}i_p,j_1\mathrm{}j_q}y_{}^{i_1}\mathrm{}y_{}^{i_p}x_{}^{j_1}\mathrm{}x_{}^{j_q},$$
(40)
where all $`y_{}^i`$’s and $`x_{}^i`$’s are either $`\stackrel{}{h_{}^i},v_{}^i`$ or $`u_{}^i`$, and the coefficients $`a_{k,i_1\mathrm{}i_p,j_1\mathrm{}j_q}`$ are covariant tensors symmetric with respect to $`i_1\mathrm{}i_p`$ and antisymmetric in $`j_1\mathrm{}j_q`$. It is simple to see that a section $`a`$ belongs to $`\mathrm{\Gamma }W^{}(\mathrm{\Lambda }^{})^q`$ iff there are no terms involving explicit $`h_{}^i`$’s in the above expression.
###### Lemma 5.6
1. For any $`i=1,\mathrm{},l`$, $`J^{}d\lambda ^i=v_{}^i`$;
2. for any $`i,j`$, $`_{\stackrel{}{h_i}}v_{}^j=0`$, and $`_{\stackrel{}{h_i}}h_{}^j`$ and $`_{\stackrel{}{h_i}}u_{}^j`$ belong to $`\mathrm{\Gamma }K^{}`$;
3. for any $`i,j`$, $`\pi (v_{}^i,h_{}^j)=\delta _{ij},\pi (v_{}^i,v_{}^j)=0,\pi (v_{}^i,u_{}^j)=0`$;
4. the commutatant of $`\{v_{}^1,\mathrm{},v_{}^l\}`$ in $`\mathrm{\Gamma }W`$ is $`\mathrm{\Gamma }W^{}`$.
Proof. (i) $`<J^{}d\lambda ^i,v_j>=<d\lambda ^i,J_{}v_j>=<d\lambda ^j,\frac{}{\lambda ^i}>=\delta _{ij}`$. Similarly, we have $`<J^{}d\lambda ^i,u_j>=0`$ and $`<J^{}d\lambda ^i,h_j>=0`$. Therefore, $`J^{}d\lambda ^i=v_{}^i`$.
(ii) According to Proposition 5.5, $`_{\stackrel{}{h_i}}v_{}^j=_{\stackrel{}{h_i}}(J^{}d\lambda ^j)=0`$. Also, $`k`$, $`<_{\stackrel{}{h_i}}h_{}^j,\stackrel{}{h_k}>`$
$`=_{\stackrel{}{h_i}}<h_{}^j,\stackrel{}{h_k}><h_{}^j,_{\stackrel{}{h_i}}\stackrel{}{h_k}>=0`$. Hence it follows that $`_{\stackrel{}{h_i}}h_{}^j\mathrm{\Gamma }K^{}`$. Similarly, we can prove that $`_{\stackrel{}{h_i}}u_{}^j\mathrm{\Gamma }K^{}`$.
(iii) We have $`\pi (v_{}^i,h_{}^j)=<\pi ^\mathrm{\#}(J^{}d\lambda ^i),h_{}^j>=<\stackrel{}{h_i},h_{}^j>=\delta _{ij}`$. Similarly, we can show that $`\pi (v_{}^i,v_{}^j)=0`$ and $`\pi (v_{}^i,u_{}^j)=0`$.
(iv) Assume that $`a\mathrm{\Gamma }W`$ such that $`[a,v_{}^i]=0,i=1,\mathrm{},l`$. It thus follows that $`\{a,v_{}^i\}=0`$, where the Poisson bracket refers to the one corresponding to the fiberwise symplectic structure on $`TM`$. Thus $`a\mathrm{\Gamma }W^{}`$ according to (iii).
$`\mathrm{}`$
###### Lemma 5.7
1. $`\mathrm{\Gamma }W^{}\mathrm{\Lambda }^{}`$ is closed under the multiplication $``$ as defined by Equation (55).
2. $`\mathrm{\Gamma }W^{}\mathrm{\Lambda }^{}`$ is closed under both the operators $`\delta `$ and $`\delta ^1`$.
3. $`\mathrm{\Gamma }W^{}\mathrm{\Lambda }^{}`$ is invariant under the covariant derivative $`_X,XX(M)`$.
Proof. (i) and (ii) are obvious. For (iii), note that $`\mathrm{\Gamma }(K^{})`$ is invariant under the covariant derivative $`_X`$ according to Proposition 5.5 (iii). Hence $`\mathrm{\Gamma }W^{}\mathrm{\Lambda }^{}`$ is also invariant.
$`\mathrm{}`$
As an immediate consequense, we have the following
###### Corollary 5.8
If $`a\mathrm{\Gamma }W^{}\mathrm{\Lambda }^{}`$ and $`_\stackrel{}{h}a=0,hh`$, then $`a\mathrm{\Gamma }W^{}\mathrm{\Lambda }^{}`$.
To prove Theorem 5.1, we start with the following
###### Lemma 5.9
Under the same hypthesis as in Theorem 5.1, we have $`\gamma _0=\delta ^1\stackrel{~}{\mathrm{\Omega }}\mathrm{\Gamma }W^{}\mathrm{\Lambda }^{}`$ and $`_\stackrel{}{h}\gamma _0=0,hh`$.
Proof. According to Equation (66), we know that $`\stackrel{~}{\mathrm{\Omega }}=\mathrm{\Omega }\omega +R=R+\mathrm{}\omega _1+\mathrm{}^2\omega _2+\mathrm{}`$. By assumption, we have $`\omega _i\mathrm{\Gamma }W^{}\mathrm{\Lambda }^{},i1`$. On the other hand, according to Proposition 5.5 (iv), we know that $`R\mathrm{\Gamma }W^{}\mathrm{\Lambda }^{}`$. Therefore, $`\stackrel{~}{\mathrm{\Omega }}\mathrm{\Gamma }W^{}\mathrm{\Lambda }^{}`$. Hence $`\gamma _0\mathrm{\Gamma }W^{}\mathrm{\Lambda }^{}`$ by Lemma 5.7.
Finally, note that for any $`hh`$, $`L_\stackrel{}{h}\omega _i=i_\stackrel{}{h}(d\omega _i)+d(i_\stackrel{}{h}\omega _i)=d(i_\stackrel{}{h}\omega _i)=0`$. According to Proposition 5.5, we have $`L_\stackrel{}{h}R=_\stackrel{}{h}R=0`$. Hence $`L_\stackrel{}{h}\stackrel{~}{\mathrm{\Omega }}=0`$. It thus follows that $`_\stackrel{}{h}\gamma _0=L_\stackrel{}{h}\gamma _0=L_\stackrel{}{h}\delta ^1\stackrel{~}{\mathrm{\Omega }}=\delta ^1L_\stackrel{}{h}\stackrel{~}{\mathrm{\Omega }}=0`$.
$`\mathrm{}`$
###### Proposition 5.10
Under the same hypthesis as in Theorem 5.1, the element $`\gamma `$, defined as in Theorem A.2, belongs to $`\mathrm{\Gamma }W^{}\mathrm{\Lambda }^{}`$ and satisfies $`_\stackrel{}{h}\gamma =0,hh`$.
Proof. We prove this proposition by induction. Assume that $`\gamma _n\mathrm{\Gamma }W^{}\mathrm{\Lambda }^{}`$ and $`_\stackrel{}{h}\gamma _n=0,hh`$. It suffices to show that $`\gamma _{n+1}`$ satisfies the same conditions. By Equation (68), $`\gamma _{n+1}`$ and $`\gamma _n`$ are related by the following equation:
$$\gamma _{n+1}=\gamma _0+\delta ^1(\gamma _n+\frac{i}{\mathrm{}}\gamma _n^2),n0.$$
(41)
According to Corollary 5.8, we have $`\gamma _n\mathrm{\Gamma }W^{}\mathrm{\Lambda }^{}`$. On the other hand, by Lemma 5.7, $`\gamma _n^2\mathrm{\Gamma }W^{}\mathrm{\Lambda }^{}`$. Hence $`\gamma _{n+1}\mathrm{\Gamma }W^{}\mathrm{\Lambda }^{}`$ according to Lemma 5.7 and Lemma 5.9.
Now
$`_\stackrel{}{h}(\gamma _{n+1})`$
$`=`$ $`_\stackrel{}{h}\gamma _0+_\stackrel{}{h}\delta ^1(\gamma _n+{\displaystyle \frac{i}{\mathrm{}}}\gamma _n^2)`$
$`=`$ $`L_\stackrel{}{h}\delta ^1(\gamma _n+{\displaystyle \frac{i}{\mathrm{}}}\gamma _n^2)`$
$`=`$ $`\delta ^1(L_\stackrel{}{h}\gamma _n+{\displaystyle \frac{i}{\mathrm{}}}L_\stackrel{}{h}\gamma _n^2)`$
$`=`$ $`0.`$
Here, in the last step, we used the relation $`L_\stackrel{}{h}=L_\stackrel{}{h}`$, which follows from the fact that the symplectic connection is $`H`$-invariant. This concludes the proof.
$`\mathrm{}`$
As in Appendix, for any $`aC^{\mathrm{}}(M)`$, we denote by $`\stackrel{~}{a}W_D`$ its parallel lift, i.e., $`D\stackrel{~}{a}=0`$ and $`\stackrel{~}{a}|_{y=0}=a`$. Theorem 5.1 is in fact an immediate consequence of the following
###### Proposition 5.11
Under the same hypothesis as in Theorem 5.1,
1. if $`a=J^{}f`$ for $`fC^{\mathrm{}}(h^{})`$, then
$$\stackrel{~}{a}=\underset{k=0}{\overset{\mathrm{}}{}}\frac{1}{k!}J^{}(\frac{^kf}{\lambda ^{i_1}\mathrm{}\lambda ^{i_k}})v_{}^{i_1}\mathrm{}v_{}^{i_k};$$
(42)
2. for any $`aC^{\mathrm{}}(M)`$,
$$\stackrel{~}{a}=\underset{k=0}{\overset{\mathrm{}}{}}\frac{1}{k!}(\stackrel{}{h_{i_1}}\mathrm{}\stackrel{}{h_{i_k}}a)h_{}^{i_1}\mathrm{}h_{}^{i_k}+T,$$
where the reminder $`T`$ does not contain any terms which are pure polynomials of $`h_{}^i`$’s.
Proof. For (i), it suffices to prove that $`\stackrel{~}{a}`$ given by Equation (42) is a parallel section. According to Proposition 5.10 and Lemma 5.6, we have $`[\gamma ,\stackrel{~}{a}]=0`$. Thus it follows that $`D\stackrel{~}{a}=\delta \stackrel{~}{a}+\stackrel{~}{a}`$, which clearly vanishes since $`v_{}^i=0`$ by Proposition 5.5 (ii) and Lemma 5.6 (ii).
For (ii), recall that $`\stackrel{~}{a}`$ is determined by the iteration formula
$$a_{n+1}=a+\delta ^1(a_n+[\frac{i}{\mathrm{}}\gamma ,a_n]).$$
(43)
So it suffices to prove that
$$a_n=\underset{k=0}{\overset{n}{}}\frac{1}{k!}(\stackrel{}{h_{i_1}}\mathrm{}\stackrel{}{h_{i_k}}a_0)h_{}^{i_1}\mathrm{}h_{}^{i_k}+T_n,$$
where each term in the reminder $`T_n`$ is not a pure polynomial of $`h_{}^i`$’s. This can be proved by induction again.
Assume that this assertion holds for $`a_n`$. To show that it still holds for $`a_{n+1}`$, we need to analyze which terms in $`a_n`$ would produce pure polynomials of $`h_{}^i`$’s out of Equation (43). Since $`\gamma \mathrm{\Gamma }W^{}\mathrm{\Lambda }^{}`$, we may ignore $`\delta ^1[\frac{i}{\mathrm{}}\gamma ,a_n]`$ and only consider $`\delta ^1a_n=\delta ^1(_i_{\stackrel{}{h_i}}a_nh_{}^i+_i_{v_i}a_nv_{}^i+_i_{u_i}a_nu_{}^i)`$. From this, it is clear that those terms containing pure polynomials of $`h_{}^i`$’s arise only from $`\delta ^1(_i_{\stackrel{}{h_i}}a_nh_{}^i)`$. Now a general term in $`a_n`$ has the form $`\mathrm{}^ka_{\alpha \beta \gamma }(x)v_{}^\alpha h_{}^\beta u_{}^\gamma `$, where $`\alpha ,\beta `$ and $`\gamma `$ are multi-indexes. However,
$`_{\stackrel{}{h_i}}(a_{\alpha \beta \gamma }(x)v_{}^\alpha h_{}^\beta u_{}^\gamma )`$
$`=`$ $`(\stackrel{}{h_i}a_{\alpha \beta \gamma }(x))v_{}^\alpha h_{}^\beta u_{}^\gamma +a_{\alpha \beta \gamma }(x)(_{\stackrel{}{h_i}}v_{}^\alpha )h_{}^\beta u_{}^\gamma +a_{\alpha \beta \gamma }(x)v_{}^\alpha (_{\stackrel{}{h_i}}h_{}^\beta )u_{}^\gamma +a_{\alpha \beta \gamma }(x)v_{}^\alpha h_{}^\beta (_{\stackrel{}{h_i}}u_{}^\gamma )`$
$`=`$ $`(\stackrel{}{h_i}a_{\alpha \beta \gamma }(x))v_{}^\alpha h_{}^\beta u_{}^\gamma +a_{\alpha \beta \gamma }(x)v_{}^\alpha (_{\stackrel{}{h_i}}h_{}^\beta )u_{}^\gamma +a_{\alpha \beta \gamma }(x)v_{}^\alpha h_{}^\beta (_{\stackrel{}{h_i}}u_{}^\gamma ).`$
According to Lemma 5.6, neither $`_{\stackrel{}{h_i}}h_{}^\beta `$ nor $`_{\stackrel{}{h_i}}u_{}^\gamma `$ will be a pure polynomial of $`h_{}^i`$’s. Hence to produce a pure $`h_{}^i`$-polynomial term, one needs that $`\alpha =\gamma =0`$. And in this case, the resulting pure $`h_{}^i`$-polynomial term is $`\mathrm{}^k(\stackrel{}{h_i}a_{0\beta 0}(x))h_{}^\beta `$. In conclusion, only pure $`h_{}^i`$-polynomial terms in $`a_n`$ can give rise to pure $`h_{}^i`$-polynomial terms in $`\delta ^1a_n`$. Hence the pure $`h_{}^i`$-polynomial terms in $`a_{n+1}`$ is $`a_0+_{k=0}^n\frac{1}{k!}\frac{1}{k+1}\stackrel{}{h_i}(\stackrel{}{h_{i_1}}\mathrm{}\stackrel{}{h_{i_k}}a_0)h_{}^ih_{}^{i_1}\mathrm{}h_{}^{i_k}`$, which clearly equals to $`_{k=0}^{n+1}\frac{1}{k!}(\stackrel{}{h_{i_1}}\mathrm{}\stackrel{}{h_{i_k}}a_0)h_{}^{i_1}\mathrm{}h_{}^{i_k}`$. This concludes the proof.
$`\mathrm{}`$
## 6 Classification
This section is devoted to the classification of quantization of a non-degenerate triangular dynamical $`r`$-matrix. Our method relies heavily on the classification result of star products on a symplectic manifold. First, let us introduce the following:
###### Definition 6.1
Two quantizations $`F(\lambda )`$ and $`E(\lambda )`$ of a triangular dynamical $`r`$-matrix are said to be equivalent if there exists a $`T(\lambda ):h^{}(Ug)^h[[\mathrm{}]]`$ satisfying the condition that $`T(\lambda )=1(\text{mod }\mathrm{})`$ and $`ϵ(T(\lambda ))=1`$ such that
$$E(\lambda )=\mathrm{\Delta }T(\lambda )^1F(\lambda )T_1(\lambda \frac{1}{2}\mathrm{}h^{(2)})T_2(\lambda +\frac{1}{2}\mathrm{}h^{(1)}).$$
(44)
To justify this definition, we need the following result, which interprets this equivalence in terms of star products.
###### Theorem 6.2
Given a compatible star product $`_{\mathrm{}}`$ on the Poisson manifold $`(M,\pi )`$ associated to a triangular dynamical $`r`$-matrix $`r(\lambda )`$, assume that $`T(\lambda ):h^{}(Ug)^h[[\mathrm{}]]`$ satisfies the condition that $`T(\lambda )=1(\text{mod }\mathrm{})`$ and $`ϵ(T(\lambda ))=1`$. Then the $``$-product:
$$f\stackrel{~}{_{\mathrm{}}}g=\stackrel{}{T}^1(\stackrel{}{T}f_{\mathrm{}}\stackrel{}{T}g),f,gC^{\mathrm{}}(M)$$
(45)
is still a compatible star-product. Moreover, if $`f,gC^{\mathrm{}}(G)`$,
$$f\stackrel{~}{_{\mathrm{}}}g=\stackrel{}{E(\lambda )}(f,g),$$
(46)
where $`E(\lambda )`$ is given by Equation (44).
Thus we are lead to the following
###### Definition 6.3
Compatible star-products $`_{\mathrm{}}`$ and $`\stackrel{~}{_{\mathrm{}}}`$ are said to be strongly equivalent iff they are related by Equation (45) for some $`T(\lambda ):h^{}(Ug)^h[[\mathrm{}]]`$ satisfying the property that $`T(\lambda )=1(\text{mod }\mathrm{})`$ and $`ϵ(T(\lambda ))=1`$.
An immediate consequence of Theorem 6.2 is the following:
###### Corollary 6.4
If $`F(\lambda )`$ is a quantization of a triangular dynamical $`r`$-matrix $`r:h^{}^2g`$ and $`T(\lambda ):h^{}(Ug)^h[[\mathrm{}]]`$ satisfies the condition that $`T(\lambda )=1(\text{mod }\mathrm{})`$ and $`ϵ(T(\lambda ))=1`$, then
$$E(\lambda )=\mathrm{\Delta }T(\lambda )^1F(\lambda )T_1(\lambda \frac{1}{2}\mathrm{}h^{(2)})T_2(\lambda +\frac{1}{2}\mathrm{}h^{(1)})$$
is also a quantization of $`r(\lambda )`$.
Due to this fact, Definition (6.1) is well justified. Indeed, Theorem 6.2 allows us to reduce the classification problem of quantizations of a triangular dynamical $`r`$-matrix to that of strongly equivalent star products on $`M`$.
Remark Let $`R_E(\lambda )=E^{21}(\lambda )^1E^{12}(\lambda )`$ and $`R_F(\lambda )=F^{21}(\lambda )^1F^{12}(\lambda )`$. It is easy to see that they are related by
$$R_E(\lambda )=T_2(\lambda \frac{1}{2}\mathrm{}h^{(1)})^1T_1(\lambda +\frac{1}{2}\mathrm{}h^{(2)})^1R_F(\lambda )T_1(\lambda \frac{1}{2}\mathrm{}h^{(2)})T_2(\lambda +\frac{1}{2}\mathrm{}h^{(1)}).$$
(47)
Alternatively, we may define a quantization of a triangular dynamical $`r`$-matrix $`r(\lambda )`$ to be an element $`R(\lambda )=1+\mathrm{}r(\lambda )+\mathrm{}UgUg[[\mathrm{}]]`$ satisfying the QDYBE, and define an equivalence of quantizations by Equation (47). This definition sounds weaker than our original one. We do not know, however, at this moment whether these two definitions are equivalent. It would be interesting to have this clarified.
To prove Theorem 6.2, we need a lemma.
###### Lemma 6.5
Assume that $`T(\lambda ):h^{}(Ug)^h[[\mathrm{}]]`$ is as in Theorem 6.2, then
1. $`(ad_\theta )^n(T1)=(\frac{1}{2})^n_{i_1\mathrm{}i_n}\frac{^nT}{\lambda ^{i_1}\mathrm{}\lambda ^{i_n}}h_{i_1}\mathrm{}h_{i_n}`$;
2. $`\mathrm{\Theta }(T1)\mathrm{\Theta }^1=T_1(\lambda \frac{1}{2}\mathrm{}h^{(2)})`$;
3. $`\mathrm{\Theta }(1T)\mathrm{\Theta }^1=T_2(\lambda +\frac{1}{2}\mathrm{}h^{(1)})`$;
4. $`\mathrm{\Theta }(TT)\mathrm{\Theta }^1=T_1(\lambda \frac{1}{2}\mathrm{}h^{(2)})T_2(\lambda +\frac{1}{2}\mathrm{}h^{(1)})`$.
Proof. (i) We prove this equation by induction. Obviously, it holds for $`n=0`$. Assume that it holds for $`n=k`$. Now
$`(ad_\theta )^{k+1}(T1)`$
$`=`$ $`ad_\theta [({\displaystyle \frac{1}{2}})^k{\displaystyle \underset{i_1\mathrm{}i_k}{}}{\displaystyle \frac{^kT}{\lambda ^{i_1}\mathrm{}\lambda ^{i_k}}}h_{i_1}\mathrm{}h_{i_k}]`$
$`=`$ $`({\displaystyle \frac{1}{2}})^k{\displaystyle \frac{1}{2}}{\displaystyle \underset{i_1\mathrm{}i_k}{}}{\displaystyle \underset{i}{}}([h_i{\displaystyle \frac{}{\lambda ^i}},{\displaystyle \frac{^kT}{\lambda ^{i_1}\mathrm{}\lambda ^{i_k}}}h_{i_1}\mathrm{}h_{i_k}]`$
$`[{\displaystyle \frac{}{\lambda ^i}}h_i,{\displaystyle \frac{^kT}{\lambda ^{i_1}\mathrm{}\lambda ^{i_k}}}h_{i_1}\mathrm{}h_{i_k}])`$
$`=`$ $`({\displaystyle \frac{1}{2}})^{k+1}{\displaystyle \underset{i_1\mathrm{}i_{k+1}}{}}{\displaystyle \frac{^{k+1}T}{\lambda ^{i_1}\mathrm{}\lambda ^{i_{k+1}}}}h_{i_1}\mathrm{}h_{i_{k+1}}.`$
(ii) We have
$`\mathrm{\Theta }(T1)\mathrm{\Theta }^1`$
$`=`$ $`\mathrm{exp}(\mathrm{}ad_\theta )(T1)`$
$`=`$ $`{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{k!}}(\mathrm{}ad_\theta )^k(T1)`$
$`=`$ $`{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{k!}}({\displaystyle \frac{\mathrm{}}{2}})^k{\displaystyle \frac{^kT}{\lambda ^{i_1}\mathrm{}\lambda ^{i_k}}}h_{i_1}\mathrm{}h_{i_k}`$
$`=`$ $`T_1(\lambda {\displaystyle \frac{1}{2}}\mathrm{}h^{(2)}).`$
(iii) is proved similarly, and (iv) follows from (ii) and (iii).
$`\mathrm{}`$
Proof of Theorem 6.2 If $`f,gC^{\mathrm{}}(h^{})`$, then $`\stackrel{}{T}f=f`$ and $`\stackrel{}{T}g=g`$ since $`ϵ(T)=1`$. Hence
$$\stackrel{}{T}^1(\stackrel{}{T}f_{\mathrm{}}\stackrel{}{T}g)=fg.$$
Now if $`fC^{\mathrm{}}(h^{})`$ and $`gC^{\mathrm{}}(G)`$,
$`\stackrel{}{T}f_{\mathrm{}}\stackrel{}{T}g`$
$`=`$ $`f_{\mathrm{}}\stackrel{}{T}g`$
$`=`$ $`\stackrel{}{\mathrm{\Theta }}(f,\stackrel{}{T}g)`$
$`=`$ $`\stackrel{}{\mathrm{\Theta }(1T)}(f,g)\text{ (by Lemma }\text{6.5}\text{)}`$
$`=`$ $`\stackrel{}{T_2(\lambda +{\displaystyle \frac{1}{2}}\mathrm{}h^{(1)})\mathrm{\Theta }}(f,g)`$
$`=`$ $`{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{k!}}({\displaystyle \frac{\mathrm{}}{2}})^k\stackrel{}{(h_{i_1}\mathrm{}h_{i_k}{\displaystyle \frac{^kT}{\lambda ^{i_1}\mathrm{}\lambda ^{i_k}}})\mathrm{\Theta }}(f,g)`$
$`=`$ $`\stackrel{}{(1T)\mathrm{\Theta }}(f,g)`$
$`=`$ $`\stackrel{}{T}(f_{\mathrm{}}g).`$
Here in the last equality, we used the fact that $`\stackrel{}{T}`$ does not involve any derivative $`\frac{}{\lambda ^i}`$. So we have proved that $`\stackrel{}{T}^1(\stackrel{}{T}f_{\mathrm{}}\stackrel{}{T}g)=\stackrel{}{\mathrm{\Theta }}(f,g)`$.
Finally, assume that $`f,gC^{\mathrm{}}(G)`$. According to Theorem 3.8,
$`\stackrel{}{T}f_{\mathrm{}}\stackrel{}{T}g`$
$`=`$ $`\stackrel{}{(F(\lambda )\mathrm{\Theta }})(\stackrel{}{T}f,\stackrel{}{T}g)`$
$`=`$ $`\stackrel{}{F(\lambda )\mathrm{\Theta }(TT)}(f,g)\text{ (by Lemma }\text{6.5}\text{)}`$
$`=`$ $`\stackrel{}{F(\lambda )T_1(\lambda {\displaystyle \frac{1}{2}}\mathrm{}h^{(2)})T_2(\lambda +{\displaystyle \frac{1}{2}}\mathrm{}h^{(1)})\mathrm{\Theta }}(f,g).`$
It thus follows that
$`\stackrel{}{T}^1(\stackrel{}{T}f_{\mathrm{}}\stackrel{}{T}g)`$
$`=`$ $`\stackrel{}{\mathrm{\Delta }T(\lambda )^1F(\lambda )T_1(\lambda {\displaystyle \frac{1}{2}}\mathrm{}h^{(2)})T_2(\lambda +{\displaystyle \frac{1}{2}}\mathrm{}h^{(1)})\mathrm{\Theta }}(f,g)`$
$`=`$ $`\stackrel{}{E(\lambda )\mathrm{\Theta }}(f,g).`$
This concludes the proof.
$`\mathrm{}`$
.
The rest of the section is devoted to the classification of strongly equivalent classes of compatible star products on $`M=h^{}\times G`$. The classification of star products on a general symplectic manifold was studied by many authors, for example, see . Here we follow the elementary approach due to Bertelson, Bieliavsky and Gutt concerning invariant star products.
First we prove
###### Theorem 6.6
Let $`M=h^{}\times G`$ be the symplectic manifold corresponding to a non-degenerate dynamical $`r`$-matrix $`r(\lambda )`$. Two compatible Fedosov $``$-products are strongly equivalent iff their Weyl curvatures $`\mathrm{\Omega }_{}`$ and $`\mathrm{\Omega }`$ are strongly cohomologous, i.e., $`\mathrm{\Omega }_{}\mathrm{\Omega }=d\theta `$, where $`\theta \mathrm{\Omega }^1(M)[[\mathrm{}]]`$ is $`G\times H`$-invariant and satisfies $`i_\stackrel{}{h}\theta =0`$, $`hh`$.
From now on, in this section, by $`M`$ we always mean the symplectic manifold $`h^{}\times G`$ associated with a non-degenerate dynamical $`r`$-matrix. Let $`g=hm`$ be a reductive decomposition as in Lemma 4.5, and $`\{h_1,\mathrm{},h_l\}`$ a basis in $`h`$, and $`\{e_1,\mathrm{},e_m\}`$ a basis of $`m`$. If we choose $`v_i=\frac{}{\lambda ^i}`$ and $`u_i=\stackrel{}{e_i}`$, then $`\{\stackrel{}{h_1},\mathrm{},\stackrel{}{h_l},v_1,\mathrm{},v_l,u_1\mathrm{},u_m\}`$ constitutes a local (in fact global in this case) basis of tangent fibers of $`M`$, which satisfies all the required properties as in the construction preceding Lemma 5.6. In what follows, we will fix such a choice, and denote by $`\{\stackrel{}{h_{}^1},\mathrm{},\stackrel{}{h_{}^l},v_{}^1,\mathrm{}v_{}^l,u_{}^1,\mathrm{},u_{}^m\}`$ its dual basis.
###### Lemma 6.7
Assume that $`D`$ is an Abelian connection defining a compatible $``$-product on $`M`$ as in Corollary 5.2. For any $`aC^{\mathrm{}}(M)`$, let
$$\stackrel{~}{a}=\mathrm{}^kD_{k,\alpha \beta \gamma }(a)v_{}^\alpha h_{}^\beta u_{}^\gamma \mathrm{\Gamma }(W)$$
(48)
be its parallel lift, where $`\alpha ,\beta `$ and $`\gamma `$ are multi-indexes, and $`D_{k,\alpha \beta \gamma }`$ are certain differential operators on $`M`$. If an operator $`D_{k,\alpha \beta \gamma }`$ involves a derivative of $`\lambda h^{}`$, then the corresponding term satisfies $`|\alpha |>0`$.
Proof. As it is known, $`\stackrel{~}{a}`$ is given by the iteration formula
$$a_{n+1}=a_0+\delta ^1(a_n+[\frac{i}{\mathrm{}}\gamma ,a_n]),$$
so it suffices to show that $`a_n`$ possesses such a property for any $`n`$, which we shall prove by induction.
Assume that $`a_n`$ possesses this property, and we need to show that so does $`a_{n+1}`$. Let $`\mathrm{}^kD_{k,\alpha \beta \gamma }(a)v_{}^\alpha h_{}^\beta u_{}^\gamma `$ be a term in $`a_{n+1}`$, where $`D_{k,\alpha \beta \gamma }`$ involves a derivative of $`\lambda `$. There are two possible sources that this term may come from. One is from $`\delta ^1[\frac{i}{\mathrm{}}\gamma ,a_n]`$. Since this operation does not affect the part involving derivatives on $`a`$, so it must come from a term having the form:
$$\delta ^1[\frac{i}{\mathrm{}}\gamma ,\mathrm{}^k^{}D_{k,\alpha \beta \gamma }(a)v_{}^\alpha ^{}h_{}^\beta ^{}u_{}^\gamma ^{}],$$
(49)
where $`\mathrm{}^k^{}D_{k,\alpha \beta \gamma }(a)v_{}^\alpha ^{}h_{}^\beta ^{}u_{}^\gamma ^{}`$ is one of the terms in $`a_n`$. By assumption, we know that $`|\alpha ^{}|>0`$. Since $`\gamma \mathrm{\Gamma }W^{}\mathrm{\Lambda }^{}`$, it follows from Lemma 5.6 that any resulting term in Equation (49) has at least a factor $`v_{}^\alpha ^{}`$.
Another possible source is from $`\delta ^1(a_n)`$. Now
$$\delta ^1a_n=\underset{i}{}(_{\frac{}{\lambda ^i}}a_n)v_{}^i+\underset{i}{}(_{\stackrel{}{h_i}}a_n)h_{}^i+\underset{i}{}(_{\stackrel{}{e_i}}a_n)u_{}^i.$$
If it arises from the first term, we are done. Assume that it comes from the second term: $`(_{\stackrel{}{h_i}}a_n)h_{}^i`$. Let $`\mathrm{}^kD_{k,\alpha \beta \gamma }(a)v_{}^\alpha h_{}^\beta u_{}^\gamma `$ be a general term in $`a_n`$, then
$$_{\stackrel{}{h_i}}(\mathrm{}^kD_{k,\alpha \beta \gamma }(a)v_{}^\alpha h_{}^\beta u_{}^\gamma )=\mathrm{}^k(\stackrel{}{h_i}D_{k,\alpha \beta \gamma })(a)v_{}^\alpha h_{}^\beta u_{}^\gamma +\mathrm{}^kD_{k,\alpha \beta \gamma }(a)v_{}^\alpha (_{\stackrel{}{h_i}}h_{}^\beta )u_{}^\gamma +\mathrm{}^kD_{\alpha ,\beta ,\gamma }(a)v_{}^\alpha h_{}^\beta (_{\stackrel{}{h_i}}u_{}^\gamma ).$$
From this equation, it is clear that $`D_{k,\alpha \beta \gamma }`$ must already contain some derivative of $`\lambda h^{}`$. The conclusion thus follows from the inductive assumption. A similar argument applies when it arises from the last term $`(_{\stackrel{}{e_i}}a_n)u_{}^i`$. This concludes the proof.
$`\mathrm{}`$
Proof of Theorem 6.6 Our proof here is essentially a modification of the proof of Corollary 5.5.4 in .
“Necessity.” Let
$`D`$ $`=`$ $`\delta ++{\displaystyle \frac{i}{\mathrm{}}}[\gamma ,],\text{and}`$
$`D_{}`$ $`=`$ $`\delta ++{\displaystyle \frac{i}{\mathrm{}}}[\gamma _{},]`$
be the Abelian connections with Weyl curvatures $`\mathrm{\Omega }`$ and $`\mathrm{\Omega }_{}`$, respectively, and $`A:W_DW_D_{}`$ an isomorphism of algebras. It is standard that $`A`$ lifts to an automorphism of the Weyl bundle $`W`$, which will be denoted by the same symbol $`A:WW`$. Then $`A`$ is $`G\times H`$-equivariant. As in , we may assume that $`A(a)=UaU^1`$ for some $`U\mathrm{\Gamma }W_+`$, $`aW`$. We may also assume that $`U`$ is $`G\times H`$-invariant since $`A`$ is $`G\times H`$-equivariant. By assumption, we also know that $`Aa=a`$ if $`a=_0^{\mathrm{}}\frac{1}{k!}\frac{^ka_0}{\lambda ^{i_1}\mathrm{}\lambda ^{i_k}}v_{}^{i_1}\mathrm{}v_{}^{i_k}a_0C^{\mathrm{}}(h^{})`$, which is the parallel lift of $`a_0`$ according to Proposition 5.11. This implies that $`U`$ commutes with $`v_{}^i,i=1,\mathrm{},l`$, and therefore $`U\mathrm{\Gamma }W_+^{}`$ according to Lemma 5.6. Consider another Abelian connection: $`\stackrel{~}{D}a=(ADA^1)(a)=UD(U^1aU)U^1=Da[DUU^1,a]`$. Then $`\stackrel{~}{D}`$ has the same Weyl curvature as $`D`$ (see Theorem 5.5.3 and Corollary 5.5.4 in ), which is assumed to be $`\mathrm{\Omega }`$. On the other hand,
$`D_{}a\stackrel{~}{D}a`$ $`=`$ $`{\displaystyle \frac{i}{\mathrm{}}}[\gamma _{}\gamma i\mathrm{}(DUU^1),a]`$ (50)
$`=`$ $`{\displaystyle \frac{i}{\mathrm{}}}[\mathrm{\Delta }\gamma ,a].`$
Since $`U\mathrm{\Gamma }W_+^{}`$ and $`\gamma _{},\gamma \mathrm{\Gamma }W^{}\mathrm{\Lambda }^{}`$, it follows that $`\mathrm{\Delta }\gamma \mathrm{\Gamma }W^{}\mathrm{\Lambda }^{}`$. It is also clear that $`\mathrm{\Delta }\gamma `$ is $`G\times H`$-invariant. Moreover, from Equation (50), it follows that $`[\mathrm{\Delta }\gamma ,a]=0`$, if $`aW_D_{}`$. Hence $`\mathrm{\Delta }\gamma `$ is a scalar form. Thus $`\mathrm{\Omega }_{}\mathrm{\Omega }=d\mathrm{\Delta }\gamma `$. Clearly, $`\mathrm{\Delta }\gamma `$ is $`G\times H`$-invariant and $`i_\stackrel{}{h}\mathrm{\Delta }\gamma =0,hh`$.
“Sufficientity”. Assume that $`\mathrm{\Omega }_{}\mathrm{\Omega }=d\theta `$, $`\theta \mathrm{\Omega }^1(M)[[\mathrm{}]]`$ is $`G\times H`$-invariant and $`i_\stackrel{}{h}\theta =0`$, $`hh`$. Let $`\mathrm{\Omega }(t)=\mathrm{\Omega }+td\theta `$ and $`D_t=\delta ++\frac{i}{\mathrm{}}[\gamma (t),]`$ be the Abelian connection with Weyl curvature $`\mathrm{\Omega }(t)`$, where $`\gamma (t)`$ is as in Theorem A.2 satisfying $`\delta ^1\gamma (t)=0`$.
Let $`H(t)\mathrm{\Gamma }W`$ be the solution of the equation $`D_tH(t)=\theta +\dot{\gamma }(t)`$ satisfying $`H(t)|_{y=0}=0`$. Then $`H(t)`$ is $`G\times H`$-invariant since $`D_t,\theta ,\gamma (t)`$ are all $`G\times H`$-invariant. On the other hand, since $`\gamma (t)\mathrm{\Gamma }W^{}\mathrm{\Lambda }^{}`$ according to Proposition 5.10, and $`\theta \mathrm{\Gamma }\mathrm{\Lambda }^{}`$ by assumption, it follows that $`H(t)\mathrm{\Gamma }W^{}`$.
According to Theorem 5.5.3 , the solution of the Heisenberg equation:
$$\frac{d\stackrel{~}{a}}{dt}+[H(t),\stackrel{~}{a}]=0$$
(51)
establishes an isomorphism $`W_DW_D_{}`$, which is given by $`\stackrel{~}{a}(0)\stackrel{~}{a}(1)`$. In fact, $`D_t\stackrel{~}{a}(t)=0`$ if $`D\stackrel{~}{a}(0)=0`$.
Clearly, this correspondence is $`G\times H`$-equivariant since $`H(t)`$ is $`G\times H`$-invariant. So its corresponding formal differential operator $`T:C^{\mathrm{}}(M)[[\mathrm{}]]C^{\mathrm{}}(M)[[\mathrm{}]]`$ is $`G\times H`$-invariant. Finally it remains to show that $`T`$, as a formal differential operator, does not involve any derivative of $`\lambda h^{}`$.
To show this, for any $`aC^{\mathrm{}}(M)`$, let $`\stackrel{~}{a}W_D`$ be its parallel lift, and $`\stackrel{~}{a}(t)`$ the solution of Equation (51) satisfying the initial condition $`\stackrel{~}{a}(0)=\stackrel{~}{a}`$. Then $`D_t\stackrel{~}{a}(t)=0`$. Also, let $`a(t)=\stackrel{~}{a}(t)|_{y=0}`$. Write
$$\stackrel{~}{a}(t)=\mathrm{}^kD_{t,k,\alpha \beta \gamma }(a(t))v_{}^\alpha h_{}^\beta u_{}^\gamma .$$
If an operator $`D_{t,k,\alpha \beta \gamma }`$ involves a derivative to $`\lambda h^{}`$, we know that $`\alpha 0`$ according to Lemma 6.7. Since $`H(t)\mathrm{\Gamma }W^{}`$, it thus follows that $`[H(t),D_{t,k,\alpha \beta \gamma }(a(t))v_{}^\alpha h_{}^\beta u_{}^\gamma ]|_{y=0}=0`$. This implies that $`[H(t),\stackrel{~}{a}(t)]|_{y=0}=𝒟_ta(t)`$, where $`𝒟_t`$ is a formal differential operator on $`M`$ involving no derivatives of $`\lambda h^{}`$. Now Equation (51) implies that
$$\frac{da(t)}{dt}+𝒟_t(a(t))=0.$$
Therefore the equivalence operator $`T:C^{\mathrm{}}(M)[[\mathrm{}]]C^{\mathrm{}}(M)[[\mathrm{}]]`$, which sends $`a(0)`$ to $`a(1)`$, does not involve any derivative of $`\lambda h^{}`$. This concludes the proof.
$`\mathrm{}`$
As in , by $`C_{diff,0}^k(M)`$, we denote the space of differential Hochschild $`k`$-cochains on $`C^{\mathrm{}}(M)`$ (i.e. k-multidifferential operators on $`M`$) vanishing on constants, and denote by $`b:C_{diff,0}^k(M)C_{diff,0}^{k+1}(M)`$ the Hochschild coboundary operator.
###### Proposition 6.8
Suppose that $`_{\mathrm{}}`$ and $`_{\mathrm{}}^{}`$ are two compatible star-products on $`M`$:
$$u_{\mathrm{}}v=\underset{k=0}{\overset{\mathrm{}}{}}\mathrm{}^kC_k(u,v),u_{\mathrm{}}^{}v=\underset{k=0}{\overset{\mathrm{}}{}}\mathrm{}^kC_k^{}(u,v),u,vC^{\mathrm{}}(M).$$
Assume that $`_{\mathrm{}}`$ and $`_{\mathrm{}}^{}`$ coincide with each other up to order $`n`$, i.e., $`C_k=C_k^{},0kn`$. Then
1. $`(C_{n+1}C_{n+1}^{})(u,v)=\stackrel{}{B}(u,v)+(b\stackrel{}{E})(u,v)`$, where $`BC^2(h^{},(^2g)^h)`$ is a $`\delta _r`$ 2-cocycle (i.e., $`\delta _rB=0`$), and $`E:h^{}(Ug)^h`$. Here $`\delta _r`$ denotes the coboundary operator defined by Equation (9).
2. $`C_1=\frac{1}{2}\{,\}+b\stackrel{}{c_1}`$ for some $`c_1C^{\mathrm{}}(h^{},(Ug)^h)`$.
3. If $`B=\delta _rX,XC^{\mathrm{}}(h^{},g^h)`$, then the formal operator $`T=1+\mathrm{}^n\stackrel{}{X}+\mathrm{}^{n+1}\stackrel{}{E_1}`$ transforms $`_{\mathrm{}}`$ to another star-product, which coincides with $`_{\mathrm{}}^{}`$ up to order $`n+1`$. Here $`E_1=E(u)[X,c_1]`$.
Proof. We use a similar argument as in .
(i). By definition, if either $`u`$ or $`v`$ is in $`C^{\mathrm{}}(h^{})`$, we have $`u_{\mathrm{}}v=u_{\mathrm{}}^{}v=\stackrel{}{\mathrm{\Theta }}(u,v)`$, which implies that $`(C_{n+1}C_{n+1}^{})(u,v)=0`$.
On the other hand, as it is well known, $`C_{n+1}C_{n+1}^{}`$ is a Hochschild 2-cocycle . Hence we may write
$$C_{n+1}C_{n+1}^{}=S+bT,$$
where $`S\mathrm{\Gamma }(^2TM)`$ and $`T`$ is a Hochschild 1-cochain. Since $`S`$ and $`bT`$ are, respectively, the skew-symmetric and symmetric parts of $`C_{n+1}C_{n+1}^{}`$, they share many common properties as $`C_{n+1}C_{n+1}^{}`$. In particular, both of them are $`G\times H`$-invariant and vanish when one of the argument $`u`$ or $`v`$ belongs to $`C^{\mathrm{}}(h^{})`$. This implies that $`S=\stackrel{}{B}`$, for some $`BC^{\mathrm{}}(h^{},(^2g)^h)`$. It is also standard that $`S`$ satisfies the equation: $`[\pi ,S]=0`$, which is equivalent to $`\delta _rB=0`$ according to the remark following Proposition 2.11.
Now $`M=h^{}\times G`$ clearly admits a $`G\times H`$-invariant (in fact G-biinvariant) connection. Since $`bT`$ is $`G\times H`$-invariant, according to Proposition 2.1 in , we can assume that $`T`$ is a $`G\times H`$-invariant 1-cochain. Since $`(bT)(u,v)=0`$, $`u,vC^{\mathrm{}}(h^{})`$, we have $`u(Tv)T(uv)+(Tu)v=0`$. On the other hand, since $`Tu`$ is $`G`$-invariant, it must be a function of $`\lambda h^{}`$ only, i.e., $`TuC^{\mathrm{}}(h^{})`$. Hence the restriction of the operator $`T`$ to $`C^{\mathrm{}}(h^{})`$ defines a vector field $`Y`$ on $`h^{}`$. Now since $`(bT)(u,v)=0,uC^{\mathrm{}}(h^{})`$, it follows that
$$(TY)(uv)=u(TY)(v),uC^{\mathrm{}}(h^{}),vC^{\mathrm{}}(M).$$
Hence $`TY`$ does not involve any derivative with respect to $`\lambda h^{}`$. Since $`TY`$ is $`G\times H`$-invariant, it follows that $`TY=\stackrel{}{E}`$, for some $`E:h^{}(Ug)^h`$. Therefore, $`bT=b\stackrel{}{E}`$.
(ii) It is standard that $`C_1=\frac{1}{2}\{,\}+bc_1^{}`$, where $`c_1^{}`$ is a Hochschild 1-cochain. By repeating a similar argument as in (i), we can prove that $`c_1^{}`$ can be chosen so that $`c_1^{}=\stackrel{}{c_1}`$ for some $`c_1C^{\mathrm{}}(h^{},(Ug)^h)`$.
(iii) If $`B=\delta _rX`$, then $`\stackrel{}{B}=[\pi ,\stackrel{}{X}]`$ according to the remark following Proposition 2.11. It is easy to check that the operator $`T=1+\mathrm{}^n\stackrel{}{X}+\mathrm{}^{n+1}\stackrel{}{E_1}`$ transforms $`_{\mathrm{}}`$ to another star-product, which coincides with $`_{\mathrm{}}^{}`$ up to order $`n+1`$.
$`\mathrm{}`$
As a consequence, we have
###### Corollary 6.9
If $`r`$ is a non-degenerate triangular dynamical $`r`$-matrix and $`M=h^{}\times G`$ its associated symplectic manifold, then every compatible $``$-product on $`M`$ is strongly equivalent to a Fedosov $``$-product as constructed in Corollary 5.2.
Proof. This follows essentially from the same argument as in the proof of Proposition 4.1 in . We will omit it here.
$`\mathrm{}`$
Combing with Theorem 6.6, we thus have proved:
###### Theorem 6.10
Let $`M=h^{}\times G`$ be the symplectic manifold associated with a non-degenerate triangular dynamical $`r`$-matrix $`r:h^{}^2g`$. Then the equivalent classes of compatible $``$-products on $`M`$ are classified by the relative Lie algebra cohomology (with coefficients being formal power series of $`\mathrm{}`$) $`H^2(g,h)[[\mathrm{}]]`$.
Using Theorem 6.2, we are thus lead to the following
###### Theorem 6.11
The equivalence classes of quantization of a non-degenerate triangular dynamical $`r`$-matrix $`r:h^{}^2g`$ are classified by the relative Lie algebra cohomology (with coefficients being formal power series of $`\mathrm{}`$) $`H^2(g,h)[[\mathrm{}]]`$.
Remark It would be interesting to see if this theorem can be proved by directly applying the usual classification theorem of star products on a symplectic manifold. One of the difficulties is that the characteristic class of a star product is usually difficult to computer. Recently, Tsygan comes up a nice way of redefining the characteristic class using the jet bundle. This may shed some new light on our problem.
Inspired by Kontesvich’s formality theorem, we end this section with the following:
Conjecture For an arbitrary classical triangular dynamical $`r`$-matrix $`r:h^{}^2g`$, the quantization is classified by $`_r(g[[\mathrm{}]],h)`$, the formal neighbourhood of $`r`$ in the moduli space $`(g[[\mathrm{}]],h)`$.
## Appendix A Appendix
In this section, we recall some basic ingredients of the Fedosov construction of $``$-products on a symplectic manifold, as well as some useful notations, which are used throughout the paper. For details, readers should consult .
Let $`(M,\omega )`$ be a symplectic manifold of dimension $`2n`$. Then, each tangent space $`T_xM`$ is equipped with a linear symplectic structure, which can be quantized using the standard Moyal-Weyl product. The resulting space is denoted by $`W_x`$. More precisely,
###### Definition A.1
A formal Weyl algebra $`W_x`$ associated to $`T_xM`$ is an associative algebra with a unit over $`C`$, whose elements consist of formal power series in $`\mathrm{}`$ with coefficients being formal polynomials in $`T_xM`$. In other words, each element has the form:
$$a(y,\mathrm{})=\mathrm{}^ka_{k,\alpha }y^\alpha $$
(52)
where $`y=(y^1,\mathrm{},y^{2n})`$ is a linear coordinate system on $`T_xM`$, $`\alpha =(\alpha _1,\mathrm{},\alpha _{2n})`$ is a multi-index, $`y^\alpha =(y^1)^{\alpha _1}\mathrm{}(y^{2n})^{\alpha _{2n}}`$, and $`a_{k,\alpha }`$ are constants. The product is defined according to the Moyal-Weyl rule:
$$ab=\underset{k=0}{\overset{\mathrm{}}{}}(\frac{\mathrm{}}{2})^k\frac{1}{k!}\pi ^{i_1j_1}\mathrm{}\pi ^{i_kj_k}\frac{^ka}{y^{i_1}\mathrm{}y^{i_k}}\frac{^kb}{y^{j_1}\mathrm{}y^{j_k}}.$$
(53)
Let $`W=_{xM}W_x`$. Then $`W`$ is a bundle of algebras over $`M`$, called the Weyl bundle. Its space of sections $`\mathrm{\Gamma }W`$ forms an associative algebra with unit under the fiberwise multiplications. One may think of $`W`$ as a “quantum tangent bundle” of $`M`$, whose space of sections $`\mathrm{\Gamma }W`$ gives rise to a deformation quantization for the tangent bundle $`TM`$, considered as a Poisson manifold with fiberwise linear symplectic structures. As in , by $`W^+`$ we denote the extension of the algebra $`W`$ consisting of those elements described as follows:
1. elements $`aW^+`$ are given by series (52), but the powers of $`\mathrm{}`$ can be both positive and negative;
2. the total degree $`2k+|\alpha |`$ of any term of the series is nonnegative;
3. there exists a finite number of terms with a given nonnegative total degree.
The center $`Z(W)`$ of $`\mathrm{\Gamma }W`$ consists of sections not containing $`y^{}s`$, thus can be naturally identified with $`C^{\mathrm{}}(M)[[\mathrm{}]]`$. By assigning degrees to $`y^{}s`$ and $`\mathrm{}`$ with $`\text{deg}y^i=1`$ and $`\text{deg}\mathrm{}=2`$, there is a natural filtration
$$C^{\mathrm{}}(M)\mathrm{\Gamma }(W_1)\mathrm{}\mathrm{\Gamma }(W_i)\mathrm{\Gamma }(W_{i+1})\mathrm{}\mathrm{\Gamma }(W)$$
with respect to the total degree (e.g., any individual term in the summation of the RHS of Equation (52) has degree $`2k+|\alpha |`$.)
A differential $`q`$-form with values in $`W`$ is a section of the bundle $`W^qT^{}M`$, which can be expressed locally as
$$a(x,y,\mathrm{},dx)=\mathrm{}^ka_{k,i_1\mathrm{}i_p,j_1\mathrm{}j_q}y^{i_1}\mathrm{}y^{i_p}dx^{j_1}\mathrm{}dx^{j_q}.$$
(54)
Here the coefficient $`a_{k,i_1\mathrm{}i_p,j_1\mathrm{}j_q}`$ is a covariant tensor symmetric with respect to $`i_1\mathrm{}i_p`$ and antisymmetric in $`j_1\mathrm{}j_q`$. For short, we denote the space of these sections by $`\mathrm{\Gamma }W\mathrm{\Lambda }^q`$. There is an associative product $``$ on $`\mathrm{\Gamma }W\mathrm{\Lambda }^{}`$, which naturally extends the multiplication $``$ on $`\mathrm{\Gamma }W`$ and the wedge product on $`\mathrm{\Lambda }^{}`$:
$$(a\theta )(b\omega )=(ab)(\theta \omega ),a,b\mathrm{\Gamma }W,\text{ and }\theta ,\omega \mathrm{\Lambda }^{}.$$
(55)
The usual exterior derivative on differential forms extends, in a straightforward way, to an operator $`\delta `$ on $`W`$-valued differential forms:
$$\delta a=\underset{i}{}dx^i\frac{a}{y^i},a\mathrm{\Gamma }W\mathrm{\Lambda }^{}.$$
(56)
By $`\delta ^1`$, we denote its “inverse” operator defined by:
$$\delta ^1a=\underset{i}{}\frac{1}{p+q}y^i(\frac{}{x^i}\text{ }\text{ }a)$$
(57)
when $`p+q>0`$, and $`\delta ^1a=0`$ when $`p+q=0`$, where $`a\mathrm{\Gamma }W\mathrm{\Lambda }^q`$ is homogeneous of degree $`p`$ in $`y`$.
There is a “Hodge”-decomposition:
$$a=\delta \delta ^1a+\delta ^1\delta a+a_{00},a\mathrm{\Gamma }W\mathrm{\Lambda }^{},$$
(58)
where $`a_{00}(x)`$ is the constant term of $`a`$, i.e, the $`0`$-form term of $`a|_{y=0}`$ or $`a_{00}(x)=a(x,0,0,0)`$. The operator $`\delta `$ possesses most of the basic properties of the usual exterior derivatives. For example,
$$\delta ^2=0\text{and }(\delta ^1)^2=0.$$
It is also clear that both $`\delta `$ and $`\delta ^1`$ commute with the Lie derivative, i.e., $`XX(M)`$,
$`L_X\delta =\delta L_X,\text{ and }L_X\delta ^1=\delta ^1L_X.`$ (59)
Let $``$ be a torsion-free symplectic connection on $`M`$ and
$$:\mathrm{\Gamma }W\mathrm{\Gamma }W\mathrm{\Lambda }^1$$
be its induced covariant derivative.
Consider a connection on $`W`$ of the form:
$$D=\delta ++\frac{i}{\mathrm{}}[\gamma ,],$$
(60)
with $`\gamma \mathrm{\Gamma }W\mathrm{\Lambda }^1`$.
Clearly, $`D`$ is a derivation with respect to the Moyal-Weyl product, i.e.,
$$D(ab)=aDb+Dab,a,b\mathrm{\Gamma }W.$$
(61)
A simple calculation yields that
$$D^2a=[\frac{i}{\mathrm{}}\mathrm{\Omega },a],a\mathrm{\Gamma }W,$$
(62)
where
$$\mathrm{\Omega }=\omega R+\delta \gamma \gamma \frac{i}{\mathrm{}}\gamma ^2.$$
(63)
Here $`R=\frac{1}{4}R_{ijkl}y^iy^jdx^kdx^l`$, and $`R_{ijkl}=\omega _{im}R_{jkl}^m`$ is the curvature tensor of the symplectic connection as defined by Equation (21).
A connection of the form (60) is called Abelian if $`\mathrm{\Omega }`$ is a scalar 2-form, i.e., $`\mathrm{\Omega }\mathrm{\Omega }^2(M)[[\mathrm{}]]`$. It is called a Fedosov connection if it is Abelian and in addition $`\gamma \mathrm{\Gamma }W_3\mathrm{\Lambda }^1`$. For an Abelian connection, the Bianchi identity implies that $`d\mathrm{\Omega }=D\mathrm{\Omega }=0`$, i.e., $`\mathrm{\Omega }Z^2(M)[[\mathrm{}]]`$. In this case, $`\mathrm{\Omega }`$ is called the Weyl curvature.
###### Theorem A.2
(Fedosov ) Let $``$ be a torsion-free symplectic connection, and $`\mathrm{\Omega }=\omega +\mathrm{}\omega _1+\mathrm{}Z^2(M)[[\mathrm{}]]`$ a perturbation of the symplectic form in the space $`Z^2(M)[[\mathrm{}]]`$. There exists a unique $`\gamma \mathrm{\Gamma }W_3\mathrm{\Lambda }^1`$ such that $`D`$, given by Equation (60), is a Fedosov connection, which has $`\mathrm{\Omega }`$ as the Weyl curvature and satisfies
$$\delta ^1\gamma =0.$$
Proof. It suffices to solve the equation:
$$\omega R+\delta \gamma \gamma \frac{i}{\mathrm{}}\gamma ^2=\mathrm{\Omega }.$$
(64)
This is equivalent to
$$\delta \gamma =\stackrel{~}{\mathrm{\Omega }}+\gamma +\frac{i}{\mathrm{}}\gamma ^2,$$
(65)
where
$$\stackrel{~}{\mathrm{\Omega }}=\mathrm{\Omega }\omega +R.$$
(66)
Applying the operator $`\delta ^1`$ to Equation (65) and using the Hodge decomposition (Equation (58)), we obtain
$$\gamma =\delta ^1\stackrel{~}{\mathrm{\Omega }}+\delta ^1(\gamma +\frac{i}{\mathrm{}}\gamma ^2).$$
(67)
Note that $`\gamma _{00}=0`$ since $`\gamma `$ is a $`1`$-form.
Take $`\gamma _0=\delta ^1\stackrel{~}{\mathrm{\Omega }}`$, and consider the following iteration equation:
$$\gamma _{n+1}=\gamma _0+\delta ^1(\gamma _n+\frac{i}{\mathrm{}}\gamma _n^2),n0.$$
(68)
Since the operator $``$ preserves the filtration and $`\delta ^1`$ raises it by $`1`$, $`\gamma _n`$ defined by Equation (68) converges to a unique $`\gamma \mathrm{\Gamma }W\mathrm{\Lambda }^1`$, which is clearly a solution to Equation (67). Moreover since $`\gamma _0`$ is at least of degree 3, $`\gamma `$ is indeed an element in $`\mathrm{\Gamma }W_3\mathrm{\Lambda }^1`$.
$`\mathrm{}`$
Theorem A.2 indicates that a Fedosov connection $`D`$ is uniquely determined by a torsion-free symplectic connection $``$ and a Weyl curvature $`\mathrm{\Omega }=_{i=0}^{\mathrm{}}\mathrm{}^i\omega _iZ^2(M)[[\mathrm{}]]`$. For this reason, we will say that $`D`$ is a Fedosov connection corresponding to the pair $`(,\mathrm{\Omega })`$.
If $`D`$ is a Fedosov connection, the space of all parallel sections $`W_D`$ automatically becomes an associative algebra. Let $`\sigma `$ denote the projection from $`W_D`$ to its center $`C^{\mathrm{}}(M)[[\mathrm{}]]`$ defined by $`\sigma (a)=a|_{y=0}`$.
###### Theorem A.3
(Fedosov ) For any $`a_0(x,\mathrm{})C^{\mathrm{}}(M)[[\mathrm{}]]`$ there is a unique section $`aW_D`$ such that $`\sigma (a)=a_0`$. Therefore, $`\sigma `$ establishes an isomorphism between $`W_D`$ and $`C^{\mathrm{}}(M)[[\mathrm{}]]`$ as vector spaces.
Proof. The equation $`Da=0`$ can be written as
$$\delta a=a+[\frac{i}{\mathrm{}}\gamma ,a].$$
Applying the operator $`\delta ^1`$, it follows from the Hodge decomposition (Equation (58)) that
$$a=a_0+\delta ^1(a+[\frac{i}{\mathrm{}}\gamma ,a]).$$
(69)
In analogue to the proof of Theorem A.2, we can solve this equation by the iteration formula:
$$a_{n+1}=a+\delta ^1(a_n+[\frac{i}{\mathrm{}}\gamma ,a_n]).$$
(70)
$`\mathrm{}`$
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# Optical conductivity of one-dimensional Mott insulators
## Abstract
We calculate the optical conductivity of one-dimensional Mott insulators at low energies using a field theory description. The square root singularity at the optical gap, characteristic of band insulators, is generally absent and appears only at the Luther-Emery point. We also show that only few particle processes contribute significantly to the optical conductivity over a wide range of frequencies and that the bare perturbative regime is recovered only at very large energies. We discuss possible applications of our results to quasi one-dimensional organic conductors.
Measurements of dynamical properties and in particular the optical conductivity $`\sigma (\omega )`$ are supposed to provide a stringent test of the existing theories of quasi one-dimensional (1D) systems. The behaviour of $`\sigma (\omega )`$ in the metallic regime is easily understood in terms of the Tomonaga-Luttinger theory . The situation in the Mott insulating phase is much more complicated as a spectral gap is dynamically generated by interactions. Here $`\sigma (\omega )`$ has until now only been studied by perturbative methods , which are expected to work well at high and intermediate frequencies but are not applicable to the most interesting regime of frequencies close to the optical gap. The purpose of the present work is to determine $`\sigma (\omega )`$ in 1D Mott insulators for all frequencies much smaller than the bandwidth, which is the large scale in the field theory approach to the problem. In particular we obtain for the first time the true behaviour of $`\sigma (\omega )`$ just above the optical gap.
An important property of one-dimensional systems that significantly simplifies our analysis is spin-charge separation, which occurs at energies much smaller than the bandwidth. In this regime $`\sigma (\omega )`$ is determined solely by the charge degrees of freedom. The standard description of the charge sector of 1D Mott insulator is given by the sine-Gordon model (SGM)
$$H_{sG}=\text{d}x\left[4\pi (\mathrm{\Pi })^2+\frac{1}{16\pi }(_x\varphi )^2+2\mu \mathrm{cos}(\beta \varphi )\right].$$
(1)
Here the momentum and coordinate densities obey the standard commutation relation $`[\mathrm{\Pi }(x),\varphi (y)]=i\delta (xy)`$. Throughout this letter we set the charge velocity and $`\mathrm{}`$ equal to one.
The cosine term in the Hamiltonian is related to Umklapp processes and the value of the sine-Gordon coupling constant $`\beta `$ is determined by the interactions. The Umklapp processes are relevant for $`\beta ^2<1`$ and dynamically generate a spectral gap $`M`$, which is related to $`\mu `$ by (30). For $`1/2<\beta ^2<1`$ the spectral gap is related to optical gap $`\mathrm{\Delta }`$ (i.e. the gap seen in the optical absorption) by $`\mathrm{\Delta }=2M`$ whereas for $`\beta ^2<1/2`$ solitonic bound states are formed below $`2M`$.
Our calculations of $`\sigma (\omega )`$ are based on the exact solution of the SGM and in particular on the work of Smirnov . We confine our analysis to the repulsive regime $`1/2<\beta ^2<1`$, where the excitation spectrum consists of charged particles and holes (solitons and anti-solitons), which do not form bound states. At the “Luther-Emery” point $`\beta ^2=1/2`$ the SGM is equivalent to the theory of free spinless massive Dirac fermions. In this limit the solitons become non-interacting particles and the Mott insulator turns into a conventional band insulator. In the limit $`\beta ^21`$ the SGM acquires an SU(2) symmetry and describes the Hubbard model at half-filling in the regime of weak interactions and $`\sigma (\omega )`$ was recently determined in .
The optical conductivity is related to the imaginary part of the current-current correlation function, $`\chi (\omega ,q)=j_qj_q`$ by
$$\sigma (\omega >0)=\mathrm{Im}\left\{\chi (\omega ,q=0)\right\}/\omega .$$
(2)
The current density operator is proportional to the momentum density
$$j_q=A^{1/2}\mathrm{\Pi }_q,\mathrm{\Pi }_q=𝑑x\mathrm{\Pi }(t,x)e^{iqx}.$$
(3)
The non-universal coefficient $`A^{1/2}`$ depends on the detailed structure of the underlying microscopic lattice model.
Using the spectral representation one can express the optical conductivity at $`T=0`$ as a sum over matrix elements of the zero wave vector Fourier component of the momentum operator:
$$\sigma (\omega >0)=\frac{A}{\omega }\underset{n}{}|0|\mathrm{\Pi }_0|n|^2\delta [\omega (E_nE_0)].$$
(4)
Here $`|0`$ and $`|n`$ represent the ground state and excited states with energies $`E_0`$ and $`E_n`$ respectively. The difficulties in computation of the optical response are related to the fact that one requires not only the knowledge of the spectrum $`E_n`$, but also of the matrix elements of the momentum operator. The exact expressions for the matrix elements $`n|\mathrm{\Pi }_0|0`$ are extracted from the exact solution by means of the so-called form factor bootstrap procedure . This approach is particularly efficient for strongly interacting integrable models with spectral gaps, because for a given energy $`\omega `$ the spectral representation for the imaginary part contains only a finite number of terms (in the absence of bound states at most $`[\omega /\mathrm{\Delta }]`$ terms). In practice the spectral sum is found to converge extremely rapidly, so that a very good approximate description can be obtained by taking into account intermediate states with at most four particles . The multi-particle matrix elements become essential only at very high energies where the field theory can no longer be used to describe the underlying lattice model anyway.
In order to compute (4) we need to introduce a suitable spectral representation. In the parameter regime we study, the spectrum contains only solitons and anti-solitons with relativistic dispersion $`e(p)=\sqrt{p^2+M^2}`$. It is useful to parametrize the spectrum in terms of a rapidity variable $`\theta `$
$$p=M\mathrm{sinh}\theta ,e=M\mathrm{cosh}\theta ,$$
(5)
Solitons and anti-solitons are distinguished by the internal index $`\epsilon =\pm `$. A state of $`n`$ solitons/anti-solitons with rapidities $`\{\theta _k\}`$ and internal indices $`\{\epsilon _k\}`$ is denoted by: $`|\theta _n\mathrm{}\theta _1_{\epsilon _n\mathrm{}\epsilon _1}`$. Its total energy $`E`$, momentum $`P`$ and electric charge $`Q`$ are:
$$P=M\underset{k=1}{\overset{n}{}}\mathrm{sinh}\theta _k,E=M\underset{k=1,}{\overset{n}{}}\mathrm{cosh}\theta _k,Q\underset{k=1}{\overset{n}{}}\epsilon _k.$$
(6)
In terms of this basis $`\sigma (\omega )`$ is expressed as
$`\sigma (\omega )`$ $`=`$ $`{\displaystyle \frac{2\pi ^2A}{\omega }}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{\epsilon _i}{}}{\displaystyle \frac{\text{d}\theta _1\mathrm{}\text{d}\theta _n}{(2\pi )^nn!}\left|f^j(\theta _1\mathrm{}\theta _n)_{\epsilon _1\mathrm{}\epsilon _n}\right|^2}`$ (7)
$`\times `$ $`\delta (M{\displaystyle \underset{k}{}}\mathrm{sinh}\theta _k)\delta (\omega M{\displaystyle \underset{k}{}}\mathrm{cosh}\theta _k)`$ (8)
$`=`$ $`\sigma _2(\omega )+\sigma _4(\omega )+\mathrm{}`$ (9)
Here
$$f^j(\theta _1\mathrm{}\theta _n)_{\epsilon _1\mathrm{}\epsilon _n}0|j(0,0)|\theta _n\mathrm{}\theta _1_{\epsilon _n\mathrm{}\epsilon _1}$$
(10)
are the form factors of the current operator, $`\sigma _2(\omega )`$ and $`\sigma _4(\omega )`$ represent the contributions from 2 and 4-particle processes and the dots indicate processes involving higher number of (anti)solitons. We note that as a consequence of symmetry properties only intermediate states with an even number of particles contribute to this correlation function. From (9) it is easy to see that only 2-particle processes contribute up to energies $`\omega =4M`$, only 2 and 4-particle processes up to $`\omega =6M`$ and so on.
The form factors (10) have been determined in and can be used to calculate the first few terms in the expansion (9). We find
$$\sigma _2(\omega )=\frac{2A\mathrm{\Theta }(\omega 2M)}{\omega ^2\sqrt{\omega ^24M^2}}|f(\theta )|^2,$$
(11)
where $`\mathrm{\Theta }(x)`$ is the Heaviside function,
$`f(\theta )=f^j(\theta )_+=f^j(\theta )_+={\displaystyle \frac{2\pi M}{i\beta }}{\displaystyle \frac{\mathrm{sinh}\theta /2}{\mathrm{cosh}\left(\frac{\theta +i\pi }{2\xi }\right)}}`$ (12)
$`\times \mathrm{exp}\left\{{\displaystyle _0^{\mathrm{}}}𝑑t{\displaystyle \frac{\mathrm{sinh}^2t(1i\theta /\pi )\mathrm{sinh}t(\xi 1)}{t\mathrm{sinh}2t\mathrm{cosh}t\mathrm{sinh}t\xi }}\right\}`$ (13)
and
$$\theta =2\mathrm{a}\mathrm{r}\mathrm{c}\mathrm{c}\mathrm{o}\mathrm{s}\mathrm{h}(\stackrel{~}{\omega }),\xi =\beta ^2/(1\beta ^2),\stackrel{~}{\omega }=\omega /2M.$$
(14)
The four particle contribution is of the form
$`\sigma _4(\omega )={\displaystyle \frac{\mathrm{\Theta }(\omega 4M)}{192\omega \pi ^2M^2}}{\displaystyle \underset{\epsilon _i}{}}{\displaystyle \underset{\sigma =\pm }{}}{\displaystyle _a^a}𝑑\theta {\displaystyle _{b(\theta )}^{b(\theta )}}𝑑\gamma `$ (15)
$`\times |f^j(g{\displaystyle \frac{\sigma \alpha }{2}},g+{\displaystyle \frac{\sigma \alpha }{2}}g+\theta +\gamma ,g\theta +\gamma )_{\epsilon _1\mathrm{}\epsilon _4}|^2`$ (16)
$`\times \left\{\left(\sqrt{\mathrm{cosh}^2\theta \mathrm{sinh}^2\gamma +\stackrel{~}{\omega }^2}\mathrm{cosh}\theta \mathrm{cosh}\gamma \right)^21\right\}^{\frac{1}{2}}`$ (17)
$`\times \left[\mathrm{cosh}^2\theta \mathrm{sinh}^2\gamma +\stackrel{~}{\omega }^2\right]^{\frac{1}{2}},`$ (18)
where
$`a`$ $`=`$ $`\mathrm{arccosh}(\stackrel{~}{\omega }1),b(\theta )=\mathrm{arccosh}\left[{\displaystyle \frac{\stackrel{~}{\omega }^21\mathrm{cosh}^2\theta }{2\mathrm{cosh}\theta }}\right],`$ (19)
$`g`$ $`=`$ $`\mathrm{ln}\left[{\displaystyle \frac{\mathrm{cosh}(\alpha /2)+\mathrm{exp}(\gamma )\mathrm{cosh}\theta }{\stackrel{~}{\omega }}}\right],`$ (20)
$`\alpha `$ $`=`$ $`2\mathrm{a}\mathrm{r}\mathrm{c}\mathrm{c}\mathrm{o}\mathrm{s}\mathrm{h}\left[\sqrt{\mathrm{cosh}^2\theta \mathrm{sinh}^2\gamma +\stackrel{~}{\omega }^2}\mathrm{cosh}\theta \mathrm{cosh}\gamma \right].`$ (21)
The four particle form factor is given by
$`f^j(\beta _1,\mathrm{},\beta _4)_{++}={\displaystyle \frac{4\pi ^3}{\beta }}\xi Md^2{\displaystyle \underset{k<l}{}}\zeta (\beta _k\beta _l)`$ (22)
$`\times `$ $`{\displaystyle \underset{m,n=1,2}{}}\left(\mathrm{sinh}[(\beta _{2+m}\beta _ni\pi )/\xi ]\right)^1`$ (23)
$`\times `$ $`2\mathrm{sinh}[(\beta _4+\beta _3\beta _1\beta _22\pi i)/2\xi ]`$ (24)
$`\times `$ $`\mathrm{exp}({\displaystyle \frac{1}{\xi }}{\displaystyle \underset{k}{}}\beta _k){\displaystyle \frac{\text{d}\alpha }{i\pi }\underset{k}{}\phi (\alpha \beta _k)\mathrm{cosh}(\alpha \frac{1}{2}\underset{k}{}\beta _k)}`$ (25)
$`\times `$ $`\mathrm{\Delta }(e^{2\alpha /\xi }|e^{2\beta _1/\xi },e^{2\beta _2/\xi }|e^{2\beta _3/\xi },e^{2\beta _4/\xi }).`$ (26)
The different orderings for (22) that appear in (15) can be obtained using the following property of the form factors
$`f^j(\theta _1\mathrm{}\theta _i\theta _{i+1}\mathrm{}\theta _n)_{\epsilon _1\mathrm{}\epsilon _i\epsilon _{i+1}\mathrm{}\epsilon _n}S_{\epsilon _i^{}\epsilon _{i+1}^{}}^{\epsilon _i\epsilon _{i+1}}(\theta _i\theta _{i+1})`$ (27)
$`=f^j(\theta _1\mathrm{}\theta _{i+1}\theta _i\mathrm{}\theta _n)_{\epsilon _1\mathrm{}\epsilon _{i+1}^{}\epsilon _i^{}\mathrm{}\epsilon _n},`$ (28)
where $`S`$ is the two body scattering matrix. The various functions appearing in Eqs (22) and (28) can be found in (note that our definition of $`\xi `$ differs by a factor of $`\pi `$).
The two and (one hundred times the) four-particle contributions to $`\sigma (\omega )`$ for $`\beta ^2=0.9`$ are presented in Fig.1. Most importantly, the square root singularity, being a characteristic feature of band insulators, is suppressed by the momentum dependence of the soliton-antisoliton form factor and reappears only for the Luther-Emery point $`\beta ^2=1/2`$. This effect was noted previously for the Hubbard model at half-filling which corresponds to the special SU(2)-symmetric point $`\beta ^2=1`$. We find that for any $`\beta ^21/2`$ there is a square root “shoulder” $`\sigma (\omega )\sqrt{\omega \mathrm{\Delta }}`$ for $`\omega /\mathrm{\Delta }11`$ as is shown in the inset of Fig.1. In the vicinity of the Luther-Emery point $`\beta ^2=1/2`$ we obtain the following analytical expression valid for $`\stackrel{~}{\omega }11`$:
$`\sigma (\omega ){\displaystyle \frac{\sqrt{\stackrel{~}{\omega }^21}}{[\stackrel{~}{\omega }^21]+\xi ^2\mathrm{sin}^2\gamma }},\gamma =\pi \left({\displaystyle \frac{1}{2\beta ^2}}1\right).`$ (29)
The square root singularity above $`\omega =\mathrm{\Delta }`$ for $`\beta ^2=1/2`$ is replaced by a maximum occurring at $`\omega /\mathrm{\Delta }1\gamma ^2`$.
The four particle contribution to $`\sigma `$ is seen to be insignificant at low energies and becomes larger than the two particle contribution only at $`\omega 180M`$ for $`\beta ^2=0.9`$. This suggests that the optical conductivity is well described by the combination of 2 and 4-particle contributions up to several hundred times the mass gap. Computation of higher order terms in Eq.(9) becomes cumbersome and probably of no physical interest, since the previous analysis suggests that they become important outside the region of applicability of the field theory approach to physical systems.
At frequencies much larger than the gap it is possible to determine $`\sigma (\omega )`$ by perturbative methods. The leading asymptotics can be calculated by “conformal perturbation theory” . Here the cosine interaction in (1) is considered as a (relevant) perturbation of the Gaussian model and correlation functions are calculated in a perturbative expansion in powers of the scale $`\mu `$, which then can be expressed in terms of the physical gap $`M`$ as
$$\mu =\frac{\mathrm{\Gamma }(\beta ^2)}{\pi \mathrm{\Gamma }(1\beta ^2)}\left[M\frac{\sqrt{\pi }\mathrm{\Gamma }(1/2+\xi /2)}{2\mathrm{\Gamma }(\xi /2)}\right]^{22\beta ^2}.$$
(30)
We find to leading order
$`\sigma (\omega )=2^{94\beta ^2}\left({\displaystyle \frac{\pi ^2\beta }{\mathrm{\Gamma }(2\beta ^2)}}\right)^2\mu ^2\omega ^{(4\beta ^25)}`$ (31)
$`={\displaystyle \frac{8\pi ^3\beta ^2}{\omega \mathrm{\Gamma }^2(1\beta ^2)\mathrm{\Gamma }^2(\frac{1}{2}+\beta ^2)}}\left[{\displaystyle \frac{\mathrm{\Gamma }(\frac{\xi }{2})}{2\sqrt{\pi }\mathrm{\Gamma }(\frac{1+\xi }{2})}}{\displaystyle \frac{\omega }{M}}\right]^{4\beta ^24}.`$ (32)
We emphasize that the ratio of the coefficients of the high- and low-energy asymptotics (32), (11) is fixed ,. In other words, the amplitude of the power law in (32) is tied to the overall factor in (11) and the form factor expansion must approach the perturbative result in the large-$`\omega `$ limit. A comparison between the form factor results and (32) is shown in Fig.2. We see that the asymptotic regime is not yet reached at energies as high as $`\omega 1000M`$ (in practical terms this implies that perturbation theory cannot be used to make contact with experiment). We note that the contributions due to intermediate states with 6,8,10 … particles are all positive and will make the agreement of the form factor sum with perturbation theory in the region $`\omega 1000M`$ only worse. A good way to overcome these deficiencies of bare perturbation theory is to carry out a renormalisation-group (RG) improvement as performed in . In Zamolodchikov’s scheme the RG equations for the Sine-Gordon model are given by
$$\frac{dg_{}}{dt}=\frac{g_{}g_{}}{1+\frac{g_{}}{2}},\frac{dg_{}}{dt}=\frac{g_{}^2}{1+\frac{g_{}}{2}}.$$
(33)
The solution of (33) is
$$g_{}=4\frac{1\beta ^2}{\beta ^2}\frac{\sqrt{q}}{1q},g_{}=2\frac{1\beta ^2}{\beta ^2}\frac{1+q}{1q},$$
(34)
where
$$q\left(\frac{(1q)\beta ^2}{4(1\beta ^2)}\right)^{2\beta ^22}=e^{(44\beta ^2)(tt_0)}.$$
(35)
Using $`tt_0=\mathrm{ln}\left(\frac{\sqrt{\pi }e^{3/4}M}{2^{3/2}\omega }\right)`$ we can reexpress (32) up to higher order terms as
$$\sigma (\omega )=\frac{\pi ^3\beta ^6g_{}^2}{2\omega \mathrm{\Gamma }^2(2\beta ^2)\mathrm{\Gamma }^2(\frac{1}{2}+\beta ^2)}\left[\frac{\mathrm{\Gamma }(\frac{\xi }{2})e^{3/4}\sqrt{\xi }}{2^{7/2}\mathrm{\Gamma }(\frac{1+\xi }{2})}\right]^{4\beta ^24}.$$
(36)
The RG improved result (36) for $`\sigma (\omega )`$ is compared to the form factor result (sum of the two and four-particle contributions) in the inset of Fig.2. The agreement is rather good down to energies of the order of $`5M`$.
One possible realisation of a 1D Mott insulator are the $`(\mathrm{TMTSF})_2\mathrm{X}`$ Bechgaard salts . These materials are highly anisotropic and can be modelled as weakly coupled, quarter-filled chains. At energies or temperatures above the 1D-3D crossover scale $`E_{\mathrm{cr}}`$ the interchain coupling becomes ineffective and a description in terms of a purely 1D model with charge sector (1) should be possible . At present there is some uncertainty regarding the value of $`E_{\mathrm{cr}}`$ because interactions can renormalize its bare value, set by the interchain coupling, downwards . There is a lot of ambiguity in fitting our results to the data. The value of the optical gap $`2M`$ is not known and, as discussed above, we cannot calculate the overall normalisation of $`\sigma (\omega )`$. We therefore use these as parameters in order to obtain a good fit at large $`\omega `$ (where the theory is expected to work best as 3D effects are unimportant) to the data for any given value of $`\beta `$. We obtain reasonable agreement with the data for $`\beta ^20.9`$, which corresponds to a Luttinger liquid parameter of $`K_\rho 0.23`$. This value is consistent with previous estimates (see the discussion in ).
As is clear from Fig. 3, the model (1) seems to apply well at high energies, but becomes inadequate at energies of the order of about $`10`$ times the Mott gap ($`1600/\mathrm{cm}`$ in $`(\mathrm{TMTSF})_2\mathrm{PF}_6`$). Spectral weight is transferred to lower energies and physics beyond that of a pure 1D Mott insulator emerges. There are at least two mechanisms that should be taken into account in this range of energies. Firstly, a small dimerization occurs in the 1D chains and will almost certainly affect the structure of $`\sigma (\omega )`$ around its maximum. Secondly, the interchain hopping is no longer negligible and ought to be taken into account.
In summary, we have exactly calculated $`\sigma (\omega )`$ for a pure 1D Mott insulator in a low-energy effective field theory approach. We have determined the threshold behaviour for the first time and found it to exhibit a universal square root increase for any $`\beta ^2>1/2`$. This is in contrast to the well-known suqare-root singularity that appears at the Luther-Emery point $`\beta ^2=1/2`$. In the “low” energy region ($`\omega /\mathrm{\Delta }<50`$) the optical conductivity is dominated by the two-particle form factor contribution with a small correction from four-particle processes. This means that the entire optical transport is dominated by two-particle processes! We furthermore have shown that the leading asymptotic behaviour obtained in perturbation theory is a good approximation only at extremely large frequencies, whereas RG-improved perturbation theory works well over a large region of energies.
We are grateful to A. Schwartz for generously providing us with the experimantal data and to F. Gebhard, T. Giamarchi, E. Jeckelmann and S. Lukyanov for important comments and discussions. We thank the Isaac Newton Institute for Mathematical Sciences, where this work was completed, for hospitality.
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# Group representations and the Euler characteristic of elliptically fibered Calabi–Yau threefolds
## 1. The group $`G`$
###### Definition 1.1.
An elliptic Calabi–Yau threefold with section is a proper, flat map $`\pi :XB`$ from a nonsingular projective complex threefold $`X`$ with trivial canonical bundle to a nonsingular surface $`B`$, whose general fiber is an elliptic curve, and which admits a section $`\sigma :BX`$. (During certain parts of our discussion, we shall also assume that the rank of the Mordell–Weil group $`MW(X/B)`$ of the elliptic fibration is zero.)
Any such $`X`$ is a resolution of a possibly singular, Weierstrass model $`\overline{\pi }:WB`$ . $`W`$ can be described (locally) by a “Weierstrass equation”
(1.1)
$$y^2=x^3+fx+g,$$
where $`f`$ and $`g`$ are sections of line bundles on the base $`B`$.
###### Lemma 1.2.
In this set up we can naturally associate a reductive Lie group $`G`$ to the fibration as follows. Let $`[E]^{}`$ be the orthogonal complement within $`H_4(X)`$ of the elliptic fiber $`E`$, and let $`\mathrm{\Lambda }`$ be the cokernel of the natural map
$$\pi ^1:H_2(B)[E]^{}.$$
Then $`\mathrm{\Lambda }`$ serves as the coroot lattice of $`G`$, and $`\mathrm{\Lambda }U(1)`$ serves as the Cartan subgroup. Moreover, to each component of the discriminant locus is associated a local factor of the group, determined by the generic Kodaira fiber along that component and by the monodromy, as follows:
| generic Kodaira fiber | | $`I_n`$ | $`I_n`$ | $`II`$ | $`III`$ | $`IV`$ | $`IV`$ |
| --- | --- | --- | --- | --- | --- | --- | --- |
| monodromy | | $`_2`$ | $`\{e\}`$ | $`\{e\}`$ | $`\{e\}`$ | $`_2`$ | $`\{e\}`$ |
| local group factor | | $`\mathrm{Sp}([\frac{n}{2}])`$ | $`\mathrm{SU}(n)`$ | $`\{e\}`$ | $`\mathrm{SU}(2)`$ | $`\mathrm{Sp}(1)`$ | $`\mathrm{SU}(3)`$ |
| | | | | | | | |
| | $`I_0^{}`$ | $`I_n^{}`$ | $`I_n^{}`$ | $`IV^{}`$ | $`IV^{}`$ | $`III^{}`$ | $`II^{}`$ |
| | $`_3`$ or $`𝔖_3`$ | $`_2`$ | $`\{e\}`$ | $`_2`$ | $`\{e\}`$ | $`\{e\}`$ | $`\{e\}`$ |
| | $`G_2`$ | $`\mathrm{SO}(2n+7)`$ | $`\mathrm{SO}(2n+8)`$ | $`F_4`$ | $`E_6`$ | $`E_7`$ | $`E_8`$ |
###### Proof.
To specify a connected reductive group, it is enough to specify a compact torus together with the collection of characters of that torus which will serve as the weights for the semisimple part of the group. The torus in turn can be described as $`\mathrm{\Lambda }U(1)`$ for some lattice $`\mathrm{\Lambda }`$, which is the form used in the statement of the lemma. (This choice of torus is dictated by physical considerations.)
To complete the specification, the weight spaces must be given. The Weierstrass model is singular along a (reducible) curve $`C`$; the general singularity over each irreducible component of $`C`$ is a rational double point . Let us consider the intersection configuration of the exceptional curves and the exceptional divisors on $`X`$. In most cases, the intersection matrix is (up to a sign) the unique Cartan matrix of a Lie algebra $`𝔤`$; here we also find the non-simply laced algebras, as the exceptional curves might undergo a monodromy transformation as they move in the exceptional divisors along the curve $`C`$. In some cases a more delicate argument is needed .
Note that the group is semisimple precisely when $`MW(X/B)`$ has rank $`0`$.
###### Corollary 1.3.
Let $`G`$ be a non-simply laced group (a local factor of the entire gauge group) associated to the singularities over a curve $`C`$ in the discriminant, as in the above proof. Then the exceptional curves in one homology class are parameterized by a curve $`C^{}`$, a finite branched cover of $`C`$. The cover is of degree $`2`$ unless $`G`$ is locally isomorphic to $`G_2`$; in the latter case, the degree of the cover is $`3`$.
###### Definition 1.4.
With the notation of the above corollary, we write $`g(C^{})=g^{}`$.
###### Remark 1.5.
If $`B`$ is ruled, from $`X`$ and the group $`G`$ we can construct a K3 surface $`S`$ with a gauge bundle $`H`$, the “heterotic dual” of $`X`$. Many of the physics predictions stated in this paper were originally derived by analyzing this duality.
###### Definition 1.6.
Let $`𝚺B`$ be the ramification locus of $`\pi `$. $`𝚺`$ is a divisor. We write
(1.2)
$$𝚺=𝚺_0_{ii}𝚺_i,$$
where $`𝚺_0`$ is the (possibly reducible) component over which the “general” singular fiber is a node (Kodaira type $`I_1`$), and each $`𝚺_i`$ is an irreducible component of $`𝚺𝚺_0`$.
###### Remark 1.7.
Since $`XB`$ is an elliptic Calabi–Yau , then:
$$𝒪_XK_X=\pi ^{}(K_B+\frac{1}{12}𝚺),\text{ and }𝚺|12K_B|.$$
The choice of the notation and of the indices in Definition 1.6 is motivated by the gauge group associated to the resolution of the general singular point of $`𝚺_i`$: we denote in fact by $`G_i`$ this group and because $`G_0`$ is trivial, the relevant groups are $`G_i,i1`$.
## 2. A first look at $`𝓡`$: when $`𝑿\mathbf{=}𝑾`$.
Let us define the fundamental invariant:
###### Definition 2.1.
$`𝓡`$$`=\frac{1}{2}\chi _{top}(X)+30K_B^2.`$
The following holds:
###### Theorem 2.2.
Under the hypothesis of Definition 1.1, suppose that in addition $`X=W`$ is a smooth Weierstrass model. Then $`𝓡`$=0.
###### Proof.
Following the algorithm provided in one can in fact show that
$$\chi _{top}(X)=60K_B^2.$$
This statement is also buried in the proof of \[28, (3.10)\]; even though the author claims it only for $`B=^2,𝔽_1,^1\times ^1`$. ∎
We should also point out that $`𝓡`$=0 is not a sufficient condition for $`W`$ to be smooth (see also Section 7). We will give an alternative proof of this theorem in Corollary 6.10; this proof will follow the mathematical ideas that arise from considering the “vanishing of the anomalies” in string theory. The following definition is also motivated by the physics literature:
###### Definition 2.3.
$`H_{ch}=`$ $`𝓡`$ $`+\mathrm{dim}(G)\mathrm{rk}(G).`$
Theorem 8.2 describes explicitly $`H_{ch}`$. We discuss in detail the motivation behind this definition and a dictionary between the geometry and the quantum field theory in .
## 3. A look from string theory: <br>Gauge theory on $`X`$ and vanishing of the anomalies
Before studying other properties of the invariant $`𝓡`$ defined in the previous section, we look at the gauge theory interpretation of our setup. From this point of view we consider a gauge theory on $`X`$ (coupled to gravity), in which certain coefficients of the curvature are required to vanish: these are the “anomaly cancellations” whose geometric counterparts are a formula for $`𝓡`$ and certain geometric constraints. We will see that a formula for $`𝓡`$ is the geometric counterpart of the first anomaly cancellation (Theorem 3.1 and ); the second anomaly cancellation and the corresponding geometric constraints will be discussed in Section 9.
### 3.1. From elliptic fibrations to gauge theory
When one of the “type II string theories” is formulated on a ten-manifold of the form $`M^{3,1}\times X`$ with $`X`$ a Calabi–Yau threefold and $`M^{3,1}`$ a flat spacetime of dimension four, the resulting theory has a low energy approximation which takes the form of a four-dimensional quantum field theory with quite realistic physical properties (depending on certain properties of the Calabi–Yau threefold).
Elliptic Calabi–Yau threefolds with a section $`\pi :XB`$ have also been used in a different way in string theory. We can ask what happens to the type IIA theory in the limit when the Calabi–Yau metric on $`X`$ is varied so that the fibers of the map $`\pi `$ shrink to zero area. It turns out that the resulting physical theory has a low energy approximation which takes the form of a six-dimensional quantum field theory. This limiting theory can also be described more directly, in terms of the periods $`\tau (b)`$ of the elliptic curves $`\pi ^1(b)`$, regarded as a multi-valued function on $`B`$. The type IIB string theory is compactified on $`B`$ with the aid of this function, using what are known as D-branes along the discriminant locus of the map $`\pi `$. (This latter approach is known as “F-theory.”)
Fact: This six-dimensional quantum field theory includes gravity as well as a gauge field theory whose gauge group is the group $`G`$ defined in Section 1; in order to be a consistent quantum theory, the “anomalies” of this theory must vanish.
### 3.2. Curvatures, anomaly polynomial and traces
Schwarz shows that these models ($`N=1`$ theories in six dimensions with a semisimple group $`G`$) are constrained by anomaly cancellation. The anomaly is characterized by an eight-form made from curvatures of the Levi–Civita connection and of the gauge connection. This eight-form is naturally defined on an auxiliary eight-manifold $`Y`$.<sup>2</sup><sup>2</sup>2$`Y`$ is a manifold with boundary, whose boundary is the product of $`S^1`$ and the original six-dimensional spacetime.
If we have a manifold $`Y`$ equipped with a principal $`G`$-bundle $`𝒢`$ (the “gauge bundle”), then the curvature $`F`$ of the gauge connection is an $`\mathrm{ad}(𝒢)`$-valued two-form, where each fiber $`\mathrm{ad}(𝒢)_x`$ of $`\mathrm{ad}(𝒢)`$ is isomorphic to the Lie algebra $`𝔤`$ of $`G`$, with $`𝒢_x`$ acting on $`\mathrm{ad}(𝒢)_x`$ via the adjoint action of $`G`$ on $`𝔤`$. Similarly, if $`Y`$ is equipped with a (pseudo-)Riemannian metric, then the curvature $`R`$ of the Levi–Civita connection is a two-form taking values in the endomorphisms of the tangent bundle.
The “anomaly polynomial” is a differential form on $`Y`$ which involves expressions like $`\mathrm{tr}R^k`$ and $`\mathrm{Tr}_\rho F^k`$, where $`\rho `$ is some representation of the Lie algebra. These expressions are to be interpreted as follows: the representation $`\rho `$ can be regarded as a homomorphism $`\rho :𝔤\mathrm{End}(V)`$ for some (complex) vector space $`V`$. As an endomorphism of $`V`$, $`\rho (F_x)`$ can be raised to the $`k^{\text{th}}`$ power; the resulting endomorphism $`\rho (F_x)^k`$ of $`V`$ has a trace, which we denote as
$$\mathrm{Tr}_\rho F^k=\mathrm{trace}_V\rho (F)^k.$$
Although this expression might have depended on the choice of isomorphism to $`𝔤`$, in fact it is invariant under the adjoint action of $`G`$ on $`𝔤`$ and so is independent of choices.
Similarly, the expressions $`\mathrm{tr}R^k`$ are evaluated with the help of the “vector” representation of the corresponding orthogonal group.
### 3.3. Vanishing of the anomalies
The first requirement is the vanishing of the coefficient of a certain curvature term, which imposes restrictions on the choice of the group and its “matter” representations which can occur. In we discuss extensively the geometric realization of this formula, when F-theory is compactified on an elliptic Calabi–Yau threefold:
###### Theorem 3.1.
The anomalies are characterized by an eight-form, made from curvatures and gauge field two-forms. One requirement is the vanishing of the coefficient of the curvature term $`\mathrm{tr}R^4`$, where $`R`$ is the curvature of the Levi–Civita connection. In our geometric set up this leads to
(3.1)
$$𝓡=H_{ch}\mathrm{dim}(G)+\mathrm{rk}(G),$$
where $`H_{ch}`$ involves the dimension of certain representations of the group $`G`$.
The representations which occur in $`H_{ch}`$ are well defined in terms of quantum field theory; this motivates our Definition 2.3. Theorem 8.2 identifies these representations in our geometric set up. We will analyze the geometric counterpart of the following statement (“The generalized Green–Schwarz mechanism”) in Section 9:
###### Theorem 3.2.
Let us assume that $`G`$ is semisimple, and so in particular that $`\mathrm{rk}MW=0`$ and that $`G`$ is locally isomorphic to $`_iG_i`$, where $`G_i`$ are simple groups. If the requirement specified in Theorem 3.1 holds then the remaining terms of the anomaly polynomial (in a suitable normalization ) are:
(3.2)
$$\frac{9n_T}{8}(\mathrm{tr}R^2)^2+\frac{1}{6}\mathrm{tr}R^2X_i^{(2)}\frac{2}{3}X_i^{(4)}+4\underset{i<j}{}Y_{ij}$$
where $`n_T`$ denotes the number of “tensor multiplets” (which coincides with $`h^{1,1}(B)1`$ in our theories), and where
$`X_i^{(n)}`$ $`=\mathrm{Tr}_{\mathrm{adj}}F_i^n{\displaystyle \underset{\rho }{}}n_\rho \mathrm{Tr}_\rho F_i^n`$
$`Y_{ij}`$ $`={\displaystyle \underset{\rho ,\sigma }{}}n_{\rho \sigma }\mathrm{Tr}_\rho F_i^2\mathrm{Tr}_\sigma F_j^2.`$
$`\mathrm{Tr}_{adj}`$ means the trace in the adjoint representation, $`\mathrm{Tr}_\rho `$ denotes the trace in the representation $`\rho `$ of the simple group $`G_i`$ (see the above subsection), $`n_\rho `$ is the multiplicity of the representation $`\rho `$ of $`G_i`$ in the matter representation,<sup>3</sup><sup>3</sup>3In the physics literature one says that there are “$`n_\rho `$ hypermultiplets in the representation $`\rho `$. and $`n_{i,j}`$ is the multiplicity of the representation $`(\rho ,\sigma )`$ of $`G_i\times G_j`$.
The Green–Schwarz cancellation mechanism (in the generalized form due to Sagnotti , see also Sadov ) says that the anomalies can be cancelled provided that (3.2) can be written in the form:
(3.3)
$$\left(𝐬\mathrm{tr}R^2+𝐭_𝐢\mathrm{tr}F_i^2\right)\left(𝐮\mathrm{tr}R^2+𝐯_𝐢\mathrm{tr}F_i^2\right),$$
where $`𝐬`$, $`𝐭_𝐢`$, $`𝐮`$, and $`𝐯_𝐢`$ are divisors on the base $`B`$ (which correspond to “tensor multiplets” in the physical theory), the product is calculated using the intersection pairing on $`B`$, and $`\mathrm{tr}F_i^2`$ is evaluated in an appropriate “fundamental” representation of $`G_i`$.
In the case that $`G`$ is not semisimple, the anomaly polynomial is also known, but it is much more complicated. In this paper, we will only consider the anomalies associated to the Green–Schwarz mechanism in the semisimple case.
## 4. About $`𝑯_{𝒄𝒉}`$: a look from the physics literature.
We state some of the physics predictions on $`H_{ch}`$, based on Schwarz’s analysis; these predictions motivated our geometric definition of $`H_{ch}`$ (see Definition 2.3). (We have only described a small number of the predictions which appear in the physics literature—others can be found in .<sup>4</sup><sup>4</sup>4Note in particular that used anomaly cancellation—as we do—to make and verify predictions about $`H_{ch}`$.)
Case 0: If $`W=X`$ is a smooth Weierstrass model, that is $`G=\{e\}`$ and $`\mathrm{dim}(G)\mathrm{rk}(G)=0`$, then the quantum field theory tells us that $`H_{ch}=0`$ and $`𝓡`$ $`=0`$, as $`H_{ch}`$ is the (sum of) dimensions of certain irreducible representations of $`G`$ (see 3.1). This is in agreement with Theorem 2.2.
Case I: If $`W`$ is singular along a single, smooth curve of genus $`g`$ of $`A_{N1}`$ singularities everywhere, we know from Section 1 that $`G=\mathrm{SU}(N)`$. The authors of show that under these hypothesis
$$H_{ch}=g(\mathrm{dim}(G)\mathrm{rk}(G)),$$
and also state that the same should hold for any isolated curve. In this case one would have:
$$𝓡=(g1)(\mathrm{dim}(G)\mathrm{rk}(G)).$$
Case II: If the group is non-simply laced (see Section 1) and $`W`$ is singular along a unique curve $`C`$ of genus $`g`$, then some of the exceptional divisors in $`X`$ mapping to $`C`$ are ruled surfaces over a curve $`C^{}`$ of genus $`g^{}`$ (see Corollary 1.3). Assume that there are $`B_1`$ branch points of the map $`C^{}C`$, and that all degenerations of the generic singular fiber occur at these branch points. The authors of show that in most such cases
$$H_{ch}=g(\mathrm{dim}(G)\mathrm{rk}(G))+(g^{}g)_0,$$
where $`_0`$ is a constant (which corresponds to the “charged dimension” of a certain representation $`\rho _0`$ of $`G`$—see Definition 8.1). In this case one would have:
(4.1)
$$𝓡=(g1)(\mathrm{dim}(G)\mathrm{rk}(G))+(g^{}g)_0.$$
In the case of $`I_{2k+1}`$ with monodromy (yielding gauge group $`G=\mathrm{Sp}(k)`$), this formula is modified to one which involves $`B_1`$ as well:
(4.2)
$$𝓡=(g1)(\mathrm{dim}(G)\mathrm{rk}(G))+(g^{}g)_0+\frac{1}{2}B_1(2k).$$
Case III: If $`W`$ is singular along a single, smooth curve of genus $`g`$, the singularities are generically of type $`A_{N1}`$ singularities, but they become of type $`A_N`$ at $`B_2`$ isolated points: we know from Section 1 that $`G=\mathrm{SU}(N)`$. The authors of and show that under these hypothesis
$$H_{ch}=g(\mathrm{dim}(G)\mathrm{rk}(G))+B_2N.$$
In this case one would have:
$$𝓡=(g1)(\mathrm{dim}(G)\mathrm{rk}(G))+B_2N.$$
In Section 8 we will prove that all of these predictions hold and give a global explanation for the above formulas; we will also derive the value of $`_0`$ (which depends on $`G`$).
## 5. Working assumptions and (most of the) notation
Our basic strategy for verifying the formula for $`𝓡`$ is as follows. On the one hand, the Euler characteristic of $`X`$ can be calculated exploiting the elliptic fibration, studying the various types of singular fibers which can occur, and assigning to each a “contribution” to the Euler characteristic. First, the generic fibers make no contributions. Second, the fibers over the curves $`𝚺_i`$ make contributions which can be accounted for in terms of the genus of $`𝚺_i`$ and of its monodromy cover as well as the type of the Kodaira fiber. This leaves the contributions from intersection points of the $`𝚺_i`$’s, or from special points along the $`𝚺_i`$’s at which the fiber becomes worse.
On the other hand, a parallel decomposition can be made of the representation theory. There are specific contributions to $`H_{ch}`$ which are associated to the various local factors $`G_i`$ of the gauge group, and depend on the genus of $`𝚺_i`$ and of its monodromy cover. If these are subtracted from our formula, what remains is a sum of contributions from the intersection points of the $`𝚺_i`$’s, or from special points along the $`𝚺_i`$’s at which the fiber becomes worse.
Thus, once the “generic” singularities have been matched up, the verification can be reduced to a local question—for each type of singular fiber, verify that its contribution to the Euler characteristic is compatible with the assignment of a factor in the representation to the fiber.
We will carry this out under some assumptions about the degenerations. To simplify matters and isolate the core of the problem, we will consider the case of a single non-abelian factor $`G_1`$ in the gauge group. We will also make some simplifying assumptions about which degenerate fibers are allowed. (The cases we consider can be extended to a more general set-up: see Remarks 6.11 and 8.8 as well as .) Our specific assumptions are as follows (see Equation (1.2)):
* The locus of enhanced gauge symmetry is over a unique smooth curve $`𝚺_1`$.
* The coefficients in the Weierstrass equation are otherwise general; following we assume that the local equations can be determined by the data in Table 1.
###### Proposition 5.1.
Under these hypothesis the group $`G`$ alone determines the multiplicity $`m`$ of $`𝚺_1`$ in $`𝚺`$ (see \[4, Table 2\] and the Tables in Appendix I):
$$𝚺_0+m𝚺_1=𝚺|12K_B|.$$
Equivalently
$$𝚺\text{ is defined by the equation }s^m\sigma _0,$$
where $`\sigma _0`$ defines $`𝚺_0`$ and $`𝚺_1`$ is defined by $`s=0.`$
###### Definition 5.2.
We denote by $`𝛍\mathbf{(}𝐟\mathbf{)}`$ and $`𝛍\mathbf{(}𝐠\mathbf{)}`$ the multiplicity of $`f`$ and $`g`$ resp. along $`𝚺_1`$, and by $`𝛍_𝐏\mathbf{(}𝐟\mathbf{,}𝐠\mathbf{)}`$ the intersection multiplicity of $`f/s^{\mu (f)}`$ and $`g/s^{\mu (g)}`$ at a point $`P𝚺_1`$.
###### Definition 5.3.
We denote by $`X_{𝚺_1}`$ the singular fiber of Kodaira type over the general point of $`𝚺_1`$.
We denote by $`\{Q_1,\mathrm{}Q_C\}`$, the singularities of $`𝚺_0`$ away from $`𝚺_1`$: these are cusps: $`C`$ is then the number of cusps of $`𝚺_0`$.
We denote by $`X_{𝚺_0}`$ the singular (nodal) fiber over the general point of $`𝚺_0`$ while $`X_C`$ is the singular (cuspidal) fiber over each point $`Q_j`$.
If $`𝚺_0`$ and $`𝚺_1`$ are disjoint, all the degenerate elliptic fibers are the ones described above; if $`𝚺_0𝚺_1\mathrm{}`$ there are other degenerate elliptic fibers, not necessarily of Kodaira type, over each intersection point. A complete classification of such degenerations is not available, except in the case of simple normal crossings , and the list of possibilities could be quite complicated.
These points (the $`P_{\mathrm{}}^i`$ below) are exactly the singularities of $`𝚺`$ along $`𝚺_1`$; the roots of $`\sigma _0mods`$ determine the intersection of $`𝚺_1`$ and $`𝚺_0`$.
In particular:
###### Proposition 5.4.
Our assumptions imply the following:
* The equation defining $`𝚺_0mods=0`$ splits in the product of at most two factors: $`\beta _1^{r_1}\beta _2^{r_2}\sigma _1mods`$. Each $`\beta _i`$ is irreducible, and together with $`r_i`$ is determined by the choice of the group $`G`$ (see the Tables in Appendix I).
* $`\beta _i`$ is smooth near $`𝚺_1`$: $`(s,\beta _i)`$ are local coordinates around each intersection point and $`r_i`$ is the intersection multiplicity of $`\beta _i`$ with $`𝚺_1`$. We write:
$$𝚺_0𝚺_1=\{P_1^1,\mathrm{}P_1^{B_1},P_2^1,\mathrm{}P_2^{B_2}\},$$
where $`B_i`$ is the number of the distinct roots of $`\beta _i`$; note that if $`r_i=1`$, the $`P_i^j`$’s are points of simple normal crossings intersections.
* (Equivalently:) $`𝚺_0𝚺_1=(12K_Bm𝚺_1)𝚺_1=r_1B_1+r_2B_2.`$
* The degenerate elliptic fiber over each point $`P_i^{\mathrm{}}`$ and the local equation around $`P_i^{\mathrm{}}`$ does not depend on $`\mathrm{}`$, but only on $`i=1,2`$: without loss of generality we write $`X_{P_i}=\pi ^1(P_{\mathrm{}}^i)`$. We write the local equation in Table 2.
* The intersection multiplicity $`\mu _{P_i^{\mathrm{}}}(f,g)`$ does not depend on $`\mathrm{}`$, but only on $`i=1,2`$; we denote it by $`\mu _i(f,g)`$.
###### Proof.
It follows from . ∎
###### Proposition 5.5.
* If $`G=\mathrm{SO}(k)`$ ($`k>8`$), and $`m`$ is as defined in Proposition 5.1 then $`\mathrm{rk}(G)+2=m`$;
* if $`G=e,\mathrm{SU}(2),\mathrm{SU}(3),\mathrm{SO}(8),E_6,E_7,E_8`$ , and $`J`$ is regular, then $`\mathrm{rk}(G)+2=m`$, $`\mu _P(f,g)=0`$;
* if $`G=\mathrm{SU}(n)`$ and $`J`$ has a pole along $`𝚺_1`$, then $`\mathrm{rk}(G)+1=m`$, $`\mu (f),\mu (g)=0`$,
Note that $`G=\{e\}`$ corresponds here to the Kodaira fiber of type $`II`$.
###### Proof.
It can be verified by inspection and explicit computations. ∎
We list the values of $`m,r_1`$, and $`r_2`$ in the Tables in Appendix I.
## 6. De constructing $`𝓡`$
In this section we set up an algorithm to compute $`𝓡`$, the fundamental invariant defined in 2.1.
We break up the contributions to $`\chi _{top}`$ as follows:
###### Lemma 6.1.
The following lines add to the topological Euler characteristic of $`X`$:
$`{\displaystyle \frac{1}{2}}\chi _{top}(_{i,\mathrm{}}\pi ^1(P_{\mathrm{}}^i))=`$ $`{\displaystyle \frac{1}{2}}\chi _{top}(X_{P_1})B_1+{\displaystyle \frac{1}{2}}\chi _{top}(X_{P_2})B_2`$
$`{\displaystyle \frac{1}{2}}\chi _{top}(\pi ^1(𝚺_1_{\mathrm{},i}P_{\mathrm{}}^i))=`$ $`\chi _{top}(X_{𝚺_1})[1g(𝚺_1){\displaystyle \frac{1}{2}}B_1{\displaystyle \frac{1}{2}}B_2]`$
$`{\displaystyle \frac{1}{2}}\chi _{top}(\pi ^1(𝚺_0_jQ_j_{\mathrm{},i}P_{\mathrm{}}^i)=`$ $`{\displaystyle \frac{1}{2}}\chi _{top}(𝚺_0_jQ_j_{\mathrm{},i}P_{\mathrm{}}^i)`$
$`{\displaystyle \frac{1}{2}}\chi _{top}(_j\pi ^1(Q_j))=`$ $`{\displaystyle \frac{1}{2}}\chi _{top}(X_C)C`$
###### Proof.
We compute the Euler characteristic of $`X`$ via the structure of elliptic fibration (the Euler characteristic of the general fiber is zero) and Mayer–Vietoris’ sequence. ∎
Now we want to effectively calculate each contribution in the above equations in terms of quantities which depend on the singularities along $`C`$ and the group $`G`$.
Note also that the singularities along $`C`$ are determined by the geometry of the discriminant locus on $`B`$; this is in turn determined by the intersections of a section of some multiple of $`K_B`$ with $`𝚺_1`$ (see Remark 1.7).
### 6.1. De constructing $`\chi _{top}(𝚺_0_jQ_j_{\mathrm{},i}P_{\mathrm{}}^i)`$:
XX
We start with the following definitions:
###### Definition 6.2.
(Defining $`\alpha _j`$.) If $`\varphi _1:B_1B`$ is the blow up of a point $`PD`$, with exceptional divisor $`E`$ and
$$D_1=\varphi _1^{}(D)\alpha _1(P)E,\text{ its strict transform.}$$
###### Corollary 6.3.
With the above notation:
$$(K_{B_1}+D_1)D_1=(K_B+D)D\alpha _1(P)(\alpha _1(P)1).$$
In particular, if $`P`$ is a smooth point of $`D`$, $`\alpha _1(P)=1`$; if $`P`$ is a cuspidal point of $`D`$, $`\alpha _1(P)=2`$.
###### Definition 6.4.
Let $`\varphi _{n(i)}^i\mathrm{}\varphi _v^i\mathrm{}\varphi _1^i`$ be the embedded resolution of $`𝚺_0`$ around the point $`P_1^i`$ and $`\{\alpha _v^1\}:=\{\alpha _v(P_1^i)\}_{v=1}^{n(i)}`$ the collection of the integers as in 6.2 ($`v`$ depends on $`i`$, but we believe the distinction is clear.) Let us define:
$$ϵ_j=[\alpha _v^j(\alpha _v^j1)\mathrm{\#}\varphi ^1(P_j^k)].$$
If $`P_1^i`$ is a smooth point, $`ϵ_1=1`$; for the cuspidal points $`Q_j`$ we have $`ϵ=1`$.
###### Corollary 6.5.
$$\chi _{top}(𝚺_0P_j)=\chi _{top}(\stackrel{~}{𝚺_0})\mathrm{\#}\varphi ^1(P_j)=(K_B+𝚺_0)𝚺_0+ϵ_j$$
###### Proposition 6.6.
With the above notation we have:
$$\begin{array}{ccc}\hfill \chi _{top}(𝚺_0_jQ_j_{\mathrm{},i}P_{\mathrm{}}^i)+1112K_B^2=& mK_B𝚺_1+2m𝚺_1𝚺_0+m^2𝚺_1^2+\hfill & \\ & +B_1ϵ_1+B_2ϵ_2+C.\hfill & \end{array}$$
$`ϵ_1`$ and $`ϵ_2`$ are defined in 6.4; they are determined by non-generic the singularities along $`C`$.
###### Proof.
From the previous corollary we have:
$$\chi _{top}(𝚺_0_jQ_j_{\mathrm{},i}P_{\mathrm{}}^i)=(K_B+𝚺_0)𝚺_0+ϵ_1B_1+ϵ_2B_2+1C.$$
Note that $`𝚺_0|12K_Bm𝚺_1|`$, which gives:
$$(K_B+𝚺_0)𝚺_0=1112K_B^2+mK_B𝚺_1+2m𝚺_1𝚺_0+m^2𝚺_1^2.$$
Substituting this in the above equation we obtain the statement of the proposition. ∎
### 6.2. De constructing $`C`$, the number of cusps:
KEJWSHT
###### Lemma 6.7.
$`f`$ and $`g`$ then have
$$(4K_B\mu (f)𝚺_1)(6K_B\mu (g)𝚺_1)=24K_B^2+\{4\mu (g)+6\mu (f)\}K_B𝚺_1+\mu (f)\mu (g)𝚺_{1}^{}{}_{}{}^{2}$$
intersection points, counted with multiplicity ($`f|4K_B|,g|6K_B|`$).
###### Proposition 6.8.
Then number of cusps $`C`$, away from $`𝚺_1`$ is:
$$C=24K_B^2+\{4\mu (g)+6\mu (f)\}K_B𝚺_1\mu _1(f,g)B_1\mu _2(f,g)B_2+\mu (f)\mu (g)𝚺_{1}^{}{}_{}{}^{2},$$
where $`\mu _i(f,g)`$ are defined in Propositions 5.2 and 5.4.
$`\mu _i(f,g),\mu (f)`$, and $`\mu (g)`$ depend on the equation (I.2) and are determined by the (non-generic and generic) singularities along $`C`$ .
###### Proof.
$`C`$ is the number of cusps away from $`𝚺_1`$; our assumptions in Section 5 imply that the cusps are determined by the common zeroes of the polynomials $`\{f=g=0\}`$ away from $`𝚺_1`$ (these are ordinary vanishing, see equation (I.1)). $`f`$ and $`g`$ might also vanish along $`𝚺_1`$, of orders $`\mu (f)`$ and $`\mu (g)`$; $`fmods`$ and $`gmods`$ might have a common zero along $`𝚺_1`$. The multiplicities of these latter zeros are measured by $`\mu _i(f,g)`$. (See Appendix I.) ∎
###### Proposition 6.9.
Using the formulas (6.6) and (6.8) derived above, we re-arrange the contribution to $`\chi _{top}(X)`$ in 6.1 as follows:
$`{\displaystyle \frac{1}{2}}\chi _{top}(_{i,\mathrm{}}\pi ^1(P_{\mathrm{}}^i))=`$ $`{\displaystyle \frac{1}{2}}\chi _{top}(X_{P_1})B_1+{\displaystyle \frac{1}{2}}\chi _{top}(X_{P_2})B_2`$
$`{\displaystyle \frac{1}{2}}\chi _{top}(\pi ^1(𝚺_1_{\mathrm{},i}P_{\mathrm{}}^i))=`$ $`(g(𝚺_1)1)(m){\displaystyle \frac{1}{2}}mB_1{\displaystyle \frac{1}{2}}mB_2`$
$`{\displaystyle \frac{1}{2}}\chi _{top}(\pi ^1(𝚺_0_jQ_j_{\mathrm{},i}P_{\mathrm{}}^i)+54K_B^2=`$ $`{\displaystyle \frac{1}{2}}mK_B𝚺_1+{\displaystyle \frac{1}{2}}m^2𝚺_1^2+m𝚺_1𝚺_0`$
$`+{\displaystyle \frac{1}{2}}ϵ_1B_1+{\displaystyle \frac{1}{2}}ϵ_2B_2`$
$`+{\displaystyle \frac{1}{2}}[\{4\mu (g)+6\mu (f)\}K_B𝚺_1]{\displaystyle \frac{1}{2}}\mu _1(f,g)B_1`$
$`{\displaystyle \frac{1}{2}}\mu _2(f,g)B_2+{\displaystyle \frac{1}{2}}[\mu (f)\mu (g)𝚺_{1}^{}{}_{}{}^{2}]`$
$`{\displaystyle \frac{1}{2}}\chi _{top}(_j\pi ^1(Q_j)24K_B^2=`$ $`\{4\mu (g)+6\mu (f)\}K_B𝚺_1\mu _1(f,g)B_1`$
$`\mu _2(f,g)B_2+\mu (f)\mu (g)𝚺_{1}^{}{}_{}{}^{2}`$
The entries in the left hand sides of the above equations add to $`𝓡`$; the coefficients on the right hand side are determined by the singularities, generic and non-generic, along $`CW`$.
###### Corollary 6.10.
If $`X`$ is a smooth Weierstrass model, then $`𝓡`$=0.
One of the aims of this paper is to show that the entries on the right hand side are a collection of dimensions of certain representations of $`G`$ (Main Theorem 8.2) which are determined by the singularities, generic and non-generic along $`CW`$. In Section 9 we will show that also the converse is true, that is, the assigned representations determine uniquely the geometry of $`W`$.
###### Remark 6.11.
Note that the formula in Proposition 6.9 admits an immediate generalization to cases in which there are more simple factors in the gauge group (corresponding to additional components of the discriminant). A somewhat more involved notation is required, to handle possibilities of singular curves $`𝚺_j`$ or intersections among several components, but the same geometric principles we used above will lead to a formula of the same general type.
## 7. A second look at $`𝓡`$: $`𝚺_1𝚺_0=0`$ ($`𝚺_1`$ is isolated).
We have considered in Section 2 the case $`X=W`$, we consider now the case when $`𝚺_1`$ does not intersect the rest of the discriminant locus: equivalently, $`W`$ is singular along a single curve $`C`$ and the singularities are uniform along $`C`$. This case was also considered in the physics literature, see Section 4.
Here the computations are simpler, and we can see clearly how by using the geometry of the base we can write $`𝓡`$ (that is, the equation in (6.9)) as a function of the singular locus and certain representations of the group $`G`$. The first implication of the hypothesis is that $`J:B^1`$ is well defined around $`𝚺_1`$. By analyzing the vanishing of the anomaly we find a curious relation between the Coxeter number and rank in the case of the “exceptional series” of Deligne.
###### Theorem 7.1.
If $`𝚺_1`$ does not intersect the other components of the discriminant locus then
$$\text{ }𝓡=(\mathrm{dim}(G)\mathrm{rk}(G))(g1).$$
Note that $`\mathrm{dim}(G)\mathrm{rk}(G)=\mathrm{dim}\mathrm{adj}_G\mathrm{dim}Ker(\mathrm{adj}_G)`$.
###### Proof.
Case I: $`J`$ is regular along $`𝚺_1`$ (simply laced groups in Deligne’s exceptional series.)
In this case $`G=\{e\},\mathrm{SU}(2),\mathrm{SU}(3),\mathrm{SO}(8),E_6,E_7,E_8`$, which (except for the trivial case) are precisely the simply laced groups in Deligne’s exceptional series. Here the singular fibers are of types $`II,III,IV,I_0^{},IV^{},III^{},II^{}`$, and $`m=\chi _{top}(X_{𝚺_1})<12`$.
By assumptions, $`B_1=B_2=0`$ and $`𝓡`$ becomes ((6.9)):
$`{\displaystyle \frac{1}{2}}\chi _{top}(\pi ^1(𝚺_1_{\mathrm{},i}P_{\mathrm{}}^i))=`$ $`(g(𝚺_1)1)(m)`$
$`{\displaystyle \frac{1}{2}}\chi _{top}(\pi ^1(𝚺_0_jQ_j_{\mathrm{},i}P_{\mathrm{}}^i)+54K_B^2=`$ $`{\displaystyle \frac{1}{2}}mK_B𝚺_1+{\displaystyle \frac{1}{2}}m^2𝚺_1^2`$
$`+{\displaystyle \frac{1}{2}}\{4\mu (g)+6\mu (f)\}K_B𝚺_1`$
$`+{\displaystyle \frac{1}{2}}[\mu (f)\mu (g)𝚺_1^2]`$
$`{\displaystyle \frac{1}{2}}\chi _{top}(_j\pi ^1(Q_j)24K_B^2=`$ $`\{4\mu (g)+6\mu (f)\}K_B𝚺_1+\mu (f)\mu (g)𝚺_1^2`$
Now we use the geometry of the singularities of $`W`$:
$`K_B𝚺_1+𝚺_1^2`$ $`=2g(𝚺_1)2`$
$`𝚺_0𝚺_1=012K_B𝚺_1`$ $`=m𝚺_1^2`$
By solving the system we have ($`m<12`$):
$`K_B𝚺_1`$ $`={\displaystyle \frac{m}{12m}}2(g1)`$
$`𝚺_1^2`$ $`={\displaystyle \frac{12}{12m}}2(g1)`$
By substituting the above equation in the right hand side of $`𝓡`$, we see that every term is a multiple of $`g1`$.
Then:
$$\text{ }𝓡=\frac{6}{12m}\{2m+[m2\mu (g)][m3\mu (f)]+m^2\}(g1)=\frac{6}{12m}m(m2)(g1),$$
in fact, by the definition of $`m`$, $`[m2\mu (g)][m3\mu (f)]=0`$ (see Appendix I).
In this case we also have $`\mathrm{rk}(G)=m2`$ (see Proposition 5.5) and thus:
$$\text{ }𝓡=\frac{6(\mathrm{rk}(G)+2)}{10\mathrm{rk}(G)}\mathrm{rk}(G)(g1).$$
Now, for these groups
$$\frac{6(\mathrm{rk}(G)+2)}{10\mathrm{rk}(G)}\mathrm{rk}(G)=h(G)=\mathrm{rk}(G)\mathrm{dim}(G),$$
where $`h(G)`$ is the Coxeter number (see the following Lemma 7.2).
This is in agreement with the expectations from physics (see Section 4) together with Corollary 2.3.
Case II: $`J`$ has a pole along $`𝚺_1`$ Since $`J:B^1`$ is well defined around $`𝚺_1`$, $`J_{\mathrm{}}𝚺_1=K_B𝚺_1=0`$. This together with the assumption $`0=𝚺_0𝚺_1=(12K_Bm𝚺_1)𝚺_1=m𝚺_1^2`$, implies $`2g2=0`$, that is, $`g=1.`$
The substitution in (6.6) gives $`𝓡`$ $`=0`$. This is again consistent with the expectations in Section 4.∎
###### Lemma 7.2.
The Coxeter numbers of the simply laced groups in Deligne’s exceptional series satisfy the relation:
$$h(G)=\frac{6(\mathrm{rk}(G)+2)}{10\mathrm{rk}(G)}.$$
###### Proof.
Case by case checking. ∎
This adds to the numerology of the exceptional series presented by Deligne in .
## 8. Another look at $`𝓡`$: the Main Theorem.
In the discussion below, we will describe the matter representation as a representation of the Lie algebra $`𝔤`$; it is in fact induced from a representation of the full gauge group $`G`$ associated to $`X`$.
###### Definition 8.1.
Let $`\rho `$ be a representation of a Lie algebra $`𝔤`$, with Cartan subalgebra $`𝔥`$. The charged dimension of $`\rho `$ is $`(\mathrm{dim}\rho )_{ch}=\mathrm{dim}(\rho )\mathrm{dim}(ker\rho |_𝔥)`$.
For example, if $`\rho `$ is the adjoint representation then
$$(\mathrm{dim}\mathrm{adj})_{ch}=\mathrm{dim}𝔤\mathrm{dim}𝔥=\mathrm{dim}G\mathrm{rk}G.$$
###### Theorem 8.2.
Notation as in Section 5. Then:
$$𝓡=(g1)\mathrm{dim}(adj)_{ch}+(g^{}g)\mathrm{dim}(\rho _0)_{ch}+\underset{P𝒜}{}\delta _P\mathrm{dim}(\rho _P)_{ch},$$
where $`𝒜=\{P𝚺_1𝚺_0`$ such that the fiber over $`P`$ is of Kodaira type $`\}`$, $`g^{}`$ is defined in 1.3, the representations $`\rho _P`$ all come from a small list of representations given in Table A, and the coefficient $`\delta _P`$ is $`\frac{1}{2}`$ if the representation is quaternionic and is $`1`$ if the representation is real or complex. (The quaternionic cases are labeled with $`\frac{1}{2}`$ in the Table.)
In Table A, we give the Kodaira type of the general hyperplane section through the singular fibers which occur under our hypotheses. For each type of singular fiber, we either list the associated representation $`\rho _j`$, or (in the case of monodromy) we separate the “non-isolated part” of the representation and call it $`\rho _0`$, listing any residual representation as $`\rho _j`$. In addition, in a few cases a representation occurs with multiplicity and (for later convenience at the end of section 9) we identify an irreducible representation $`\widehat{\rho }`$ in the Table.
###### Remark 8.3.
Our assumption of a smooth, flat elliptic fibration, imposes restrictions on the type of degenerate singular fibers that might occur:
(i) If $`\{e\}`$ is associated to the Kodaira type fiber $`II`$, there is a double point singularity in the fiber over the simple normal crossings intersection point of the two branches ($`𝚺_0`$ and $`𝚺_1`$). This is terminal but not canonical, leading to a smooth but not flat fibration and a non-minimal Calabi–Yau threefold. We assume then that such points do not occur: the curve is isolated and the theorem holds (see Theorem 7.1).
(ii) If $`G=\{e\}`$ (associated to the Kodaira type fiber $`I_1`$) or $`G=Sp(k)`$ (associated to the Kodaira type fiber $`I_{2k+1}`$), the resolution of the generic singularities leaves a double point singularity in the fiber over the simple normal crossings intersection point of the two branches ($`𝚺_0`$ and $`𝚺_1`$). In fact, if the equation is otherwise generic, then no small resolution exists. We assume here for simplicity that there are no such points.
(iii) If $`G`$ is associated to the Kodaira type fiber $`II^{}`$, or $`I_n^{},n12`$, the equation of the Weierstrass model is not minimal at the non-simple normal crossings intersection point of the two branches $`𝚺_0`$ and $`𝚺_1`$. In order to resolve this singularity we would need to blow up $`B`$ the basis of the fibration. In the resulting elliptic fibration (still flat and Calabi–Yau), the two branches of discriminant are separated. We assume then that such points do not occur.
###### Remark 8.4.
We have used the following notation in Table A:
* Cases with no small resolution are denoted “NSR”, and cases with non-minimal Weierstrass model are denoted “NM”.
* A dash denotes the trivial representation, whereas a blank entry denotes a situation in which there is no representation which belongs in that location.
* The classical groups $`\mathrm{SU}(n)`$, $`\mathrm{Sp}(n)`$ have representations on $`^n`$, $`^n`$, respectively, which are known as the fundamental representations and denoted by “$`\mathrm{fund}`$”. This representation is quaternionic in the case of $`\mathrm{Sp}(n)`$. The second exterior power of the fundamental representation is denoted by “$`\mathrm{\Lambda }^2`$”. In the case of $`\mathrm{Sp}(n)`$, $`\mathrm{\Lambda }^2`$ is reducible and its irreducible “traceless” part is denoted by “$`\mathrm{\Lambda }_0^2`$”.
* The classical group $`\mathrm{SO}(n)`$ has a representation on $`𝐑^n`$ called the vector representation and denoted by “$`\mathrm{vect}`$”. Its double cover $`\mathrm{Spin}(n)`$ has spinor representations. When $`n`$ is odd, there is one spinor representation, of dimension $`2^{(n1)/2}`$, denoted by “$`\mathrm{spin}`$”. When $`n`$ is even, there are two half-spinor representations, each of dimension $`2^{(n2)/2}`$, denoted by “$`\mathrm{spin}_+`$” and “$`\mathrm{spin}_{}`$”. Note that the spinor or half-spinor representations are real if $`n0,\pm 1mod8`$, complex if $`n\pm 2mod8`$, and quaternionic if $`n\pm 1,4mod8`$.
* In the case of the exceptional groups, we label representations by their dimension (given in boldface type).
###### Proof.
As we have already seen in Section 7, the intersection numbers of the various parts of the discriminant in $`B`$ determine the geometry of $`W`$ and the choice of the group $`G`$ and vice versa. Following Section 4, we write all the terms in $`𝓡`$ in Proposition 6.9, as coefficients of $`g(C)`$, the genus of the curve of singularities, the number of points where the singularities are non-generic, and $`g(C^{})=g^{}`$, when the groups are non-simply laced, and then interpret the results. The coefficients in 6.9 are determined by the group and the local geometry (the degeneration of the general rational double point) and are listed in Appendix I. We divide the proof in 3 steps.
$``$ Step I (8.1):
We show how the geometry suggests the appropriate substitutions for $`𝚺_0𝚺_1`$, $`K_B𝚺_1`$, $`𝚺_{1}^{}{}_{}{}^{2}`$ and also $`B_1`$ if the group has monodromy branched at $`B_1`$ points.
If $`B_1=B_2=0`$, then the substitutions are uniquely determined (see Section 7).
In section 9 we show how these substitutions are equivalent to certain representation-theoretic facts. If $`G\mathrm{Sp}(k)`$ or $`\mathrm{SO}(m)`$, then after the substitutions we obtain the data in Table B. That is, the resulting formula for $`𝓡`$ can be written as a sum of local terms, associated to various points $`P`$, which can be collected into a formula of the form
(8.1)
$$𝓡=(g1)(\mathrm{dim}(G)\mathrm{rk}G)+(g^{}g)_0+\underset{j=1}{\overset{2}{}}B_j_j.$$
The local contributions $`_j`$ are recorded in Table B.
In the cases $`G=\mathrm{Sp}(k)`$, $`G=\mathrm{SO}(m)`$, there are choices in making the substitutions but if a careful choice is made we can again write things in the form (8.1) (see also Section 9 for a better interpretation).
As we will point out in Remark 8.8 below, the substitutions can be formulated in a very general way which allows them to be applied in cases beyond the specific ones considered here .
$``$ Step II (8.2): We show how we can naturally interpret the entries in Table B as charged dimensions of certain representations (multiplied by the coefficient $`\delta `$), given in Table A. That is, $`_j=\delta _j\mathrm{dim}(\rho _j)_{ch}`$. If $`p`$ is not a branch point, then the (resolution of the) general elliptic surface through $`P`$ can be associated to a group $`G^{}`$ containing $`G`$, and the representation is obtained via the branching rules for the adjoint representation of $`G^{}`$.
If $`G`$ is non-simply laced, then we consider $`GG^{}`$, $`G^{}`$ simply laced, and we use again the branching rules. (This gives the representation-theoretic interpretation of the number “$`_0`$” from equations (4.1), (4.2).)
$``$ Step III (8.3). Finally we show how the number $`\delta `$ can be derived from the geometry of the degeneration of the general double point to the singularity over $`p`$. ∎
### 8.1. Step I: The substitutions
###### Proposition 8.5.
Assume that resolution of the curve of singularities $`C`$ leads to a non-simply laced group $`G`$, as in 1.3. Namely, some of the exceptional divisors are ruled over a curve $`C^{}`$, which is a finite cover of $`C`$ of degree $`d=2,3`$ ($`3`$ if and only if $`G=G_2`$), ramified at $`B_1`$ points. Write $`g^{}=g(C^{})`$, then:
$$B_1=2(g^{}g)(2d2)(g1)$$
###### Proof.
The statement follows from Hurwitz’s formula. ∎
###### Proposition 8.6.
Following the notation in Section 5, we have:
$$𝚺_1𝚺_0=r_1B_1+r_2B_2.$$
If the group $`G`$ is non-simply laced, then
$$𝚺_1𝚺_0=r_2B_2+2r_1(g^{}g)(2d2)r_1(g1),$$
if there are $`B_1`$ branch points of: $`C^{}C`$.
###### Proposition 8.7.
The appropriate substitutions for $`K_B𝚺_1`$ and $`𝚺_1^2`$ are the ones given in Table C.
###### Proof.
* When $`J`$ is finite and there is no monodromy, i.e., cases $`II`$, $`III`$, $`IV`$, $`I_0^{}`$, $`IV^{}`$, $`III^{}`$, $`II^{}`$ corresponding to the simply laced groups $`\{e\}`$, $`\mathrm{SU}(2)`$, $`\mathrm{SU}(3)`$, $`\mathrm{SO}(8)`$, $`E_6`$, $`E_7`$, $`E_8`$, in Deligne’s exceptional series, then the local geometry is given by the following equations:
$`K_B𝚺_1+𝚺_1^2`$ $`=2g(𝚺_1)2`$
$`(12K_Bm𝚺_1)𝚺_1`$ $`=r_1B_1+r_2B_2(𝚺_1𝚺_0=r_1B_1+r_2B_2),`$
which can be solved since $`m<12`$:
$$K_B𝚺_1=\frac{2m(g1)+r_1B_1+r_2B_2}{12m};𝚺_{1}^{}{}_{}{}^{2}=\frac{24(g1)+r_1B_1+r_2B_2}{12m}.$$
* When $`J`$ is finite and there is monodromy, i.e., cases $`IV`$, $`I_0^{}`$, $`I_0^{}`$, $`IV^{}`$ corresponding to groups $`\mathrm{Sp}(1)`$, $`G_2`$, $`\mathrm{SO}(7)`$, $`F_4`$ which includes the remainder of Deligne’s exceptional series, the local geometry is the same but we also use Proposition 8.5 to eliminate $`B_1`$ in favor of $`g^{}g`$:
$$K_B𝚺_1=\frac{(2m2r_1(d1))(g1)+2r_1(g^{}g)+r_2B_2}{12m};$$
$$𝚺_{1}^{}{}_{}{}^{2}=\frac{(242r_1(d1))(g1)+2r_1(g^{}g)+r_2B_2}{12m}.$$
* If $`G=\mathrm{SU}(n)`$, $`n3`$, Table 1 in Appendix I tells us that
(8.2)
$$B_1=K_B𝚺_1,$$
where $`B_1`$ is the number of non-simple normal crossings intersections. The genus formula then says that
(8.3)
$$𝚺_{1}^{}{}_{}{}^{2}=2(g1)+B_1.$$
(The case of $`\mathrm{SU}(2)`$ is similar, using $`B_1=2K_B𝚺_1`$.)
* If $`G=\mathrm{Sp}([\frac{n}{2}])`$, $`n3`$, coming from $`I_n`$ with monodromy, then $`B_1=2K_B𝚺_1`$ so that
$$K_B𝚺_1=(g^{}g)(g1).$$
Combining this with the genus formula yields
$$𝚺_{1}^{}{}_{}{}^{2}=(g^{}1)+(g1).$$
* Finally, if $`G=\mathrm{SO}(2n+7)`$ or $`\mathrm{SO}(2n+8)`$ coming from $`I_n^{}`$, $`n1`$, then $`B_2=2K_B𝚺_1𝚺_{1}^{}{}_{}{}^{2}`$ which can be combined with the genus formula and solved to give:
$$K_B𝚺_1=2(g1)+B_2;𝚺_{1}^{}{}_{}{}^{2}=4(g1)+B_2.$$
Step I now proceeds as follows: use the data in Tables 3 and 4 in Appendix I to evaluate the “local” contributions to the Euler characteristic, in the formula for $`𝓡`$ given in Proposition 6.9. Then make the substitutions given in Propositions 8.6 and 8.7 (supplementing them with Proposition 8.5 if there is monodromy) into the resulting formula; in all but a few cases (detailed below) this yields a formula of the form
$$𝓡=(g1)(\mathrm{dim}(G)\mathrm{rk}G)+(g^{}g)_0+\underset{j=1}{\overset{2}{}}B_j_j$$
with the local contributions $`_j`$ recorded in Table B. (For simplicity of notation, we define $`_0=0`$ when there is no monodromy.)
The exceptional cases are $`I_{2k+1}`$ with monodromy, and $`I_n^{}`$. In the case of $`I_{2k+1}`$ with monodromy, the formula should be written with a term $`kB_1`$ to which the substitution from Proposition 8.5 is not applied.<sup>5</sup><sup>5</sup>5We are choosing to do this in order to more easily present the formula as agreeing with a calculation in representation theory; of course, the version of this formula in which all $`B_1`$ terms have been eliminated is also perfectly valid.
In the case of $`I_n^{}`$, the term $`m𝚺_1𝚺_0`$ in the formula for $`𝓡`$ should be broken into two parts, using the substitution from Proposition 8.6 to evaluate a term of the form $`(m2)𝚺_1𝚺_0`$, but evaluating the remaining term $`2𝚺_1𝚺_0`$ as
$$2𝚺_1𝚺_0=2(12K_B𝚺_1m𝚺_{1}^{}{}_{}{}^{2})=(488m)(g1)+(242m)B_2$$
(using Proposition 8.7 for the last step).
The results of all of these manipulations are recorded in the coefficients given in Table B.
###### Remark 8.8.
It is worth observing, for possible generalizations to other cases , that the substitutions we have used can be formulated intrinsically without reference to assumptions about the particular types of degenerate fibers which occur. This is clear for the substitutions given in Propositions 8.5 and 8.6. In the case of Proposition 8.7, when $`J`$ is finite the substitution only depends on the discriminant locus. If $`J=\mathrm{}`$ and we have type $`I_n`$ along $`𝚺_1`$, consider the Weierstrass equation
(8.4)
$$y^2=x^3+fx+g$$
(which is intrinsically associated to the elliptic fibration) and note that neither $`f`$ nor $`g`$ vanishes identically along $`𝚺_1`$. The location of the singularity is given by either $`x=3g/2f`$ or (equivalently) $`x=2f^2/9g`$. There is then a divisor $`\beta `$ on $`𝚺_1`$ (in the class $`2𝚺_iB`$) represented by $`\mathrm{div}_{𝚺_1}(g)\mathrm{div}_{𝚺_1}(f)`$ or by $`2\mathrm{div}_{𝚺_1}(f)\mathrm{div}_{𝚺_1}(g)`$. In our case, this divisor coincides with the divisor $`B_1`$ (when there is monodromy) or $`2B_1`$ (when there is no monodromy) which we used in Proposition 8.7.
Similarly, if $`J=\mathrm{}`$ and we have type $`I_n^{}`$ then neither $`f/s^2`$ nor $`g/s^3`$ vanishes identically along $`𝚺_1`$. The divisor $`\beta `$ on $`𝚺_1`$, which coincides with the divisor $`B_2`$ which we used in Proposition 8.7, is represented by $`\mathrm{div}_{𝚺_1}(g/s^3)\mathrm{div}_{𝚺_1}(f/s^2)`$ or by $`2\mathrm{div}_{𝚺_1}(f/s^2)\mathrm{div}_{𝚺_1}(g/s^3)`$.
Note that this same computation could just as easily be carried out in the case of multiple components of the discriminant. The starting point would be a straightforward generalization of the equation in Proposition 6.9. Then for each component of the discriminant, one would use the corresponding substitution (according to the singularity type along that component) and manipulate the substituted formula precisely as above. The result is a division into “non-local” terms associated to the various factors of the gauge group (taking precisely the same form as above), and “local” terms associated to isolated points along the discriminant locus. We will explore this generalization further in .
### 8.2. Step II: Branching rules
In this subsection and the next, we explain how to systematically determine representations $`\rho _j`$, associated to monodromy covers and to degeneration points, whose charged dimensions reproduce the numbers $`_j`$ which were calculated in Table B.
Let $`𝔥𝔤`$ be a subalgebra of a Lie algebra. Given an irreducible representation $`\rho :𝔤\mathrm{GL}(N,)`$, a natural question is how $`\rho `$ decomposes under $`𝔥`$. The answer can be obtained by following the “branching rules” (see for example ).
The representation $`\rho _0`$
In the case of non-simply laced groups, according to the representation $`\rho _0`$ is determined by the branching rules for $`𝔤_0𝔤`$, where $`𝔤_0`$ is the non-simply laced algebra and $`𝔤`$ is the corresponding simply laced algebra (whose Dynkin diagram covers that of $`𝔤_0`$). In each such case, $`𝔤_0`$ is the fixed subalgebra of some outer automorphism of $`𝔤`$ of finite order.
###### Proposition 8.9 ().
The following branching rules hold (using the notation for representations established in Remark 8.4):
* $`\mathrm{Sp}(k)\mathrm{SU}(2k)`$ (involutive outer automorphism):
$$\mathrm{adj}\mathrm{SU}(2k)=\mathrm{adj}\mathrm{Sp}(k)\mathrm{\Lambda }_0^2$$
* $`\mathrm{Sp}(k)\mathrm{SU}(2k+1)`$ (outer automorphism):
$$\mathrm{adj}\mathrm{SU}(2k+1)=\mathrm{adj}\mathrm{Sp}(k)\mathrm{\Lambda }^2\mathrm{fund}\mathrm{fund}$$
* $`\mathrm{SO}(2k1)\mathrm{SO}(2k)`$ (involutive outer automorphism):
$$\mathrm{adj}\mathrm{SO}(2k)=\mathrm{adj}\mathrm{SO}(2k1)\mathrm{vect}$$
* $`G_2\mathrm{SO}(8)`$ (outer automorphism):
$$\mathrm{adj}\mathrm{SO}(8)=\mathrm{adj}G_2\mathrm{𝟕}\mathrm{𝟕}$$
* $`F_4E_6`$ (involutive outer automorphism):
$$\mathrm{adj}E_6=\mathrm{adj}F_4\mathrm{𝟐𝟔}$$
In the involutive cases, we have $`\rho _0`$ given as the $`(1)`$-eigenspace of the involution, i.e., the complement of $`\mathrm{adj}𝔤_0`$ within $`\mathrm{adj}𝔤`$. Thus $`\rho _0`$ coincides with $`\mathrm{\Lambda }_0`$, $`\mathrm{vect}`$, and $`\mathrm{𝟐𝟔}`$ in the first, third, and fifth cases above, respectively.
In the case of $`\mathrm{Sp}(k)\mathrm{SU}(2k+1)`$ , although the automorphism of $`𝔤`$ has order $`4`$, the monodromy action is only order $`2`$, and $`\rho _0`$ is again given by the complement of $`\mathrm{adj}𝔤_0`$ within $`\mathrm{adj}𝔤`$, i.e., $`\rho _0=\mathrm{\Lambda }^2\mathrm{fund}\mathrm{fund}`$.
In the case of $`G_2`$, the order $`3`$ monodromy action leads to the representation $`\rho _0`$ occuring with multiplicity two in the complement of $`\mathrm{adj}𝔤_0`$. (These two copies correspond to the eigenspaces for the monodromy action with eigenvalues $`e^{\pm 2\pi i/3}`$.) Thus, in this case $`\rho _0=\mathrm{𝟕}`$.
Note that in all cases, the charged dimension of the representation $`\rho _0`$ agrees with the number $`_0`$ calculated in Table B.
The representations $`\rho _j`$
Representations associated to the points $`p`$ can also be determined via branching rules, using a method pioneered by Katz and Vafa . If the general surface section through $`p`$ has a rational double point associated to $`G^{}G`$, then the representation associated to $`p`$ is determined by the corresponding branching rule (modulo a few subtleties to be discussed in the next subsection).
###### Proposition 8.10 ().
The following branching rules hold (still using the notation from Remark 8.4):
* $`\mathrm{SU}(n)\mathrm{SU}(n+1)`$:
$$\mathrm{adj}\mathrm{SU}(n+1)=\mathrm{adj}\mathrm{SU}(n)\mathrm{fund}\overline{\mathrm{fund}}\mathrm{𝟏}$$
* $`\mathrm{SU}(n)\mathrm{SO}(2n)`$:
$$\mathrm{adj}\mathrm{SO}(2n)=\mathrm{adj}\mathrm{SU}(n)\mathrm{\Lambda }^2\overline{\mathrm{\Lambda }^2}\mathrm{𝟏}$$
* $`\mathrm{SO}(2k)\mathrm{SO}(2k+2)`$:
$$\mathrm{adj}\mathrm{SO}(2k+2)=\mathrm{adj}\mathrm{SO}(2k)\mathrm{vect}\mathrm{vect}\mathrm{𝟏}$$
* $`\mathrm{Spin}(10)E_6`$:
$$\mathrm{adj}E_6=\mathrm{adj}\mathrm{Spin}(10)\mathrm{Spin}_+\mathrm{Spin}_{}\mathrm{𝟏}$$
* $`\mathrm{Spin}(12)E_7`$:
$$\mathrm{adj}E_7=\mathrm{adj}\mathrm{Spin}(12)\mathrm{Spin}_+\mathrm{Spin}_{}\mathrm{𝟏}\mathrm{𝟏}\mathrm{𝟏}$$
* $`E_6E_7`$:
$$\mathrm{adj}E_7=\mathrm{adj}E_6\mathrm{𝟐𝟕}\overline{\mathrm{𝟐𝟕}}\mathrm{𝟏}.$$
* $`E_7E_8`$:
$$\mathrm{adj}E_8=\mathrm{adj}E_7\mathrm{𝟓𝟔}\mathrm{𝟓𝟔}\mathrm{𝟏}\mathrm{𝟏}\mathrm{𝟏}.$$
(There are also non-standard embeddings of $`D_4`$ into $`D_5`$ which lead to branching rules involving the $`\mathrm{Spin}_+`$ or $`\mathrm{Spin}_{}`$ representations of $`\mathrm{SO}(8)`$ rather than the vector representation.)
Each of these branching rules takes the form
(8.5)
$$\mathrm{adj}𝔤=\mathrm{adj}𝔤_0\rho \overline{\rho }\mathrm{𝟏}$$
for some representation $`\rho `$; it is $`\rho `$ which determines the matter representation.
For example, when the general fiber of type $`\mathrm{SU}(2k)`$, degenerates to $`\mathrm{SU}(2k+1)`$, then we use the branching rule corresponding to the inclusion $`\mathrm{SU}(2k)\mathrm{SU}(2k+1)`$ to determine the correct representation “$`\mathrm{fund}`$” appearing as $`\rho `$ in the statement of the theorem.
The matter representations $`\rho _j`$ for non-simply laced groups at non-branch points can be inferred by looking at the representation of the corresponding simply laced group.
The cases $`\mathrm{SO}(12)E_7`$ and $`E_7E_8`$ (as well as the fundamental representation of $`\mathrm{Sp}(k)`$) lead to quaternionic representations and follow a somewhat different pattern, as we will explain in the next subsection. In all other cases, the representation $`\rho _j`$ determined by these branching rules has a charged dimension which agrees with the number $`_j`$ calculated in Table B.
### 8.3. Step III: Resolutions of non-generic singularities, deformation theory, complex and quaternionic representations
Up to this point, we have described the “matter” representation as a complex representation of the group G (as is customary in the physics literature<sup>6</sup><sup>6</sup>6In the physics literature, one refers to “hypermultiplets taking values in a complex representation” or, equivalently, “half-hypermultiplets taking values in a quaternionic representation.”). However, the representation we need is more accurately described as a quaternionic representation, that is, a representation into $`\mathrm{GL}(^n)`$. Given a complex representation $`\rho `$, the representation $`\rho \overline{\rho }`$ is automatically quaternionic—this is how one passes from complex to quaternionic in many cases. However, some quaternionic representations cannot be described as the sum of a complex representation with its complex conjugate. This explains the presence of the factor $`\delta =\frac{1}{2}`$ in certain terms of the formula for $`𝓡`$, since in all cases we are actually counting $`1/2`$ of the quaternionic dimension of the representation.
How do these complex and quaternionic representations show up in the geometry? Consider again the general elliptic surface passing through $`p`$. In all the cases we are considering, this surface has a rational double point singularity, which can be associated to a simply laced group $`G^{}`$. Deforming to a nearby surface we again find a rational double point, this time the one associated to the group $`G`$.
There are three possibilities for a one-parameter family of rational double points: (1) it fails to admit a simultaneous resolution of singularities, (2) it is a base-change of a family of type (1) which admits a simultaneous resolution of singularities, or (3) it admits a simultaneous resolution of singularities, and is not the base-change of a family which failed to admit such a resolution. When analyzed carefully, the Katz–Vafa prescription operates differently in these cases, depending on whether or not simultaneous resolution is possible.<sup>7</sup><sup>7</sup>7We are grateful to Sheldon Katz for correspondence on this point. It is possible to explicitly compute whether or not this is possible in each instance, using the formulas in . (One calculates the equation of the family after performing the base-change which ensures that simultaneous resolution is possible; the fact that a base-change has been performed can then be recognized from the dependence of all coefficients on $`t^k`$ rather than $`t`$, for some integer $`k`$ which represents the degree of the base-change map. See , where many of these calculations have been carried out.)
Of the branching rules described in Proposition 8.10, the first one ($`\mathrm{SU}(n)\mathrm{SU}(n+1)`$) falls in case (3), and all others fall in cases (1) and (2) (depending on whether $`\rho _j`$ is being treated as a representation of a simply laced or a non-simply laced group). There is a further distinction that can be made in case (1): making a base-change to produce a simultaneous resolution, the base-change group will act on the set of roots, and this action may or may not induce monodromy on the Dynkin diagram.
In case (1), if we perform a finite base-change, a simultaneous resolution becomes possible and the branching rules determine the representations which are involved. However, the covering group for the base-change acts on these representations, and only the invariant representation appears in the original family. In four of the branching rules from Proposition 8.10, there is monodromy on the Dynkin diagram and we have already analyzed the corresponding representations from that point of view. The representation $`\rho `$ (whose weights are represented by holomorphic curves) is mapped to the representation $`\overline{\rho }`$ (whose weights are represented by anti-holomorphic curves) with the upshot being that each ramification point on the parameter curve is associated to $`1/2`$ of the full representation. (Of course, we are not counting this as a contribution to the local representation at the branch point—this part of the representation theory is non-local, and is accounted for by the representation $`\rho _0`$.)
Note that these same four branching rules also occur in the context of case (2) families, where there is no monodromy. In these cases, the entire branching rule plays a rôle, and the quaternionic representation associated to such a point is $`\rho \overline{\rho }`$ (corresponding to the complex representation $`\rho `$). Note that the singularity is fully resolved in these cases, as is reflected in the Euler characteristic computations in Table 4 in Appendix I.
The remaining two types of branching rules, $`\mathrm{SO}(12)E_7`$ and $`E_7E_8`$, only occur in the context of case (1) in our setup, and there is no monodromy on the Dynkin diagram. In these cases, the action of the covering group similarly maps $`\rho `$ to $`\overline{\rho }`$, but in these cases the representation is quaternionic and $`\rho \overline{\rho }`$. The upshot is that the “complex representation” associated to each such point is $`1/2`$ of the quaternionic representation $`\rho `$. (Note that the covering group acts as $`1`$ on the “$`\mathrm{𝟏}`$” summands in the branching rule, so that these do not contribute as they are not invariant.) In both of these cases, the singularity of the surface is not fully resolved, as is reflected in the Euler characteristic computations in Table 4 in Appendix I.
The multiplicities of these points are slightly different in the cases of Kodaira fibers of types II and III, but the same representations occur. See , where these cases are worked out in detail.
## 9. Another look at the substitutions: <br>Representation theory.
We have seen how the degenerations of the general singularity determine certain representations of the group $`G`$; here we show that the converse also holds: once one chooses the representations which might occur, the geometry of the Calabi–Yau is completely determined by some relations in representation theory. We will at the same time verify the additional anomaly cancellations stated in Theorem 3.2.
We will verify that the generalized Green–Schwarz anomaly cancellation mechanism works in the way that was proposed by Sadov .<sup>8</sup><sup>8</sup>8We have corrected some minor numerical errors in . The factored form (3.3) is taken to be
(9.1)
$$\frac{1}{2}\left(\frac{1}{2}K_B\mathrm{tr}R^2+2𝚺_i\mathrm{tr}F_i^2\right)\left(\frac{1}{2}K_B\mathrm{tr}R^2+2𝚺_i\mathrm{tr}F_i^2\right).$$
The anomaly cancellation requirements are deduced by comparing this with equation (3.2). The coefficients of $`(\mathrm{tr}R^2)^2`$ agree due to the relation $`9n_T=K_B^2`$, which follows from Noether’s theorem on the surface $`B`$ (since $`\chi (𝒪_B)=1`$). The remaining coefficients lead to equations
$`6K_B𝚺_1(\mathrm{tr}F_i^2)`$ $`=\mathrm{Tr}_{\mathrm{adj}}F_i^2+{\displaystyle \underset{\rho }{}}n_\rho \mathrm{Tr}_\rho F_i^2`$
$`3𝚺_{i}^{}{}_{}{}^{2}(\mathrm{tr}F_i^2)^2`$ $`=\mathrm{Tr}_{\mathrm{adj}}F_i^4+{\displaystyle \underset{\rho }{}}n_\rho \mathrm{Tr}_\rho F_i^4`$
$`𝚺_i𝚺_j(\mathrm{tr}F_i^2)(\mathrm{tr}F_j^2)`$ $`={\displaystyle \underset{\rho ,\sigma }{}}n_{\rho \sigma }\mathrm{Tr}_\rho F_i^2\mathrm{Tr}_\sigma F_j^2`$
which must be evaluated using the relations in the ring of $`G`$-invariant functions. Note that in our case there is a single local factor $`G_i`$ of the gauge group $`G`$, and we can suppress the subscript $`i`$ and denote its adjoint curvature by $`F`$.
We must also specify, for each type of group, a “fundamental representation” in which to evaluate the trace $`\mathrm{tr}`$ on the left-hand side of the equations. We take $`\mathrm{tr}=\mathrm{Tr}_{\mathrm{fund}}`$ to be the trace in the usual fundamental representation for $`\mathrm{SU}(n)`$ and $`\mathrm{Sp}(k)`$, we take $`\mathrm{tr}=\frac{1}{2}\mathrm{Tr}_{\mathrm{vect}}`$ to be one-half of the trace in the vector representation for $`\mathrm{Spin}(m)`$, and we take $`\mathrm{tr}`$ to be the trace in the smallest representation of the group in the case of the exceptional groups.
Note that if we were to replace $`\mathrm{tr}`$ by some multiple of it, say $`\lambda \mathrm{tr}`$, then we would multiply $`6K_B\mathrm{\Sigma }_1`$ by $`\lambda `$ and $`3\mathrm{\Sigma }_1^2`$ by $`\lambda ^2`$. Making the geometry match the representation theory completely constrains our choice of $`\lambda `$, and we express everything below in terms of the “correct” trace for each group.
Having specified the fundamental representation, $`\mathrm{tr}F^2`$ will correspond to a basis of Casimir operators of second order, and $`(\mathrm{tr}F^2)^2`$ will be one of the basis elements for Casimir operators of the fourth order; when there is a second independent fourth-order Casimir, the second basis element can be taken to be $`\mathrm{tr}F^4`$. Traces taken in other representations can be expressed in terms of these. We have collected the data of this sort that we need (mostly taken from Erler ) in Table D (in which we use the notation $`\mathrm{spin}_{}`$ to denote either $`\mathrm{spin}`$ or $`\mathrm{spin}_\pm `$).
It is now a straightforward matter to verify the remaining anomaly cancellations. We illustrate the procedure in the case of $`G=\mathrm{SU}(n)`$, $`n4`$, with a matter representation in which the adjoint representation has multiplicity $`g`$, the fundamental representation has multiplicity $`B_2`$, and $`\mathrm{\Lambda }^2`$ has multiplicity $`B_1`$ (as specified in Theorem 8.2).
From Table D, we read off the facts which must hold in order for the gauge and mixed anomalies to cancel:
$`6K_B\mathrm{\Sigma }_1`$ $`=2n(g1)+B_2+(n2)B_1`$
$`3\mathrm{\Sigma }_1^2`$ $`=6(g1)+0B_2+3B_1`$
$`0`$ $`=2n(g1)+B_2+(n8)B_1.`$
(Note that there are two equations coming from the quartic anomaly, since there are two independent fourth order Casimirs.)
To verify these, we use the geometric relations which characterize $`g`$, $`B_1`$, and $`B_2`$, namely
$`B_1`$ $`=K_B𝚺_1`$
$`B_2`$ $`=(8K_Bn𝚺_1)𝚺_1`$
$`g`$ $`=({\displaystyle \frac{1}{2}}K_B+{\displaystyle \frac{1}{2}}𝚺_1)𝚺_1.`$
When these are substituted into the right-hand side of the proposed anomaly relations,
$`2n({\displaystyle \frac{1}{2}}K_B+{\displaystyle \frac{1}{2}}𝚺_1)+(8K_Bn𝚺_1)+(n2)(K_B)`$ $`=6K_B`$
$`6({\displaystyle \frac{1}{2}}K_B+{\displaystyle \frac{1}{2}}𝚺_1)+0(8K_Bn𝚺_1)+3(K_B)`$ $`=3𝚺_1`$
$`2n({\displaystyle \frac{1}{2}}K_B+{\displaystyle \frac{1}{2}}𝚺_1)+(8K_Bn𝚺_1)+(n8)(K_B)`$ $`=0,`$
the relations are verified.
A similar verification can be carried out in all cases. It is convenient to supplement the geometric formulas for $`g`$ and $`B_j`$’s with a formula for for $`g^{}g`$ in the case of monodromy, and to compute a quantity $`\widehat{B}`$ in a few cases (in order to match the representation $`\widehat{\rho }`$ in the representation theory, as determined in Theorem 8.2). We summarize the data in Table E. (We have omitted the relation $`g=(\frac{1}{2}K_B+\frac{1}{2}𝚺_1)𝚺_1`$, which always holds.) Carrying out the verification is then a simple exercise in combining Tables D and E, as we have done in the case of $`\mathrm{SU}(n)`$ above.
###### Remark 9.1.
As in Remark 8.8, we can express this part of the verification of the anomaly cancellation in terms which are somewhat more intrinsic. We defer the details of this to , but observe here how this can be carried out in the case of $`\mathrm{SU}(n)`$, $`n4`$.
The intrinsic geometric quantities we need are the divisor $`𝚺_0𝚺_1`$, the arithmetic genus $`p_a(𝚺_1)`$, and the divisor $`\beta `$ from Remark 8.8 (which is the intrinsic version of $`2B_1`$). We derive from these an intrinsic version of $`B_2`$, represented as $`𝚺_0𝚺_12\beta `$. Then in the anomaly cancellation requirements, we can represent the coefficient of $`\mathrm{tr}F_i^2`$ as
$$2n(p_a(𝚺_1)1)+(𝚺_0𝚺_12\beta )+\frac{n2}{2}\beta $$
and the coefficient of $`(\mathrm{tr}F_i^2)^2`$ as
$$6(p_a(𝚺_1)1)+\frac{3}{2}\beta .$$
## Appendix I: How to compute $`𝓡`$ <br>(the coefficients in Proposition 6.9 and other things)
In this section we study the local equations and the geometric data for each group and their generic degenerations.
Following we analyze the local equations in Tables 1 and 2. In Tables 3 and 4 we list, for each group, the coefficients of the right hand side of the equation defining $`𝓡`$, in Proposition 6.9. The entries of Table 1 are taken from , those of Table 3 are well known; to compute the others we need the affine equations of (I.2) and (I.1). We will work out the details for the case $`G=\mathrm{SU}(2k)`$ in Appendix II.
We need to use a more general form of the Weierstrass equation (1.1), namely
(I.2)
$$y^2+a_1xy+a_3y=x^3+a_2x^2+a_4x+a_6.$$
Since $`W`$ is assumed to be a Calabi–Yau $`a_j|jK_B|`$.
###### Definition I.1.
It is convenient to use the following:
$`b_2=a_{1}^{}{}_{}{}^{2}+4a_2`$, $`b_4=a_1a_3+2a_4`$, $`b_6=a_3^2+4a_6`$,
$`b_8=(a_{1}^{}{}_{}{}^{2}+a_2)a_6a_{4}^{}{}_{}{}^{2}`$
The coefficients in (1.1) are now:
$$f=\frac{1}{48}(b_{2}^{}{}_{}{}^{2}24b_4),g=\frac{1}{864}(b_{2}^{}{}_{}{}^{3}+36b_2b_4216b_6).$$
If $`a_j`$ (resp. $`b_j`$) vanish along $`𝚺_1`$ of order $`k`$, then we write
$$a_{j,k}=\frac{a_j}{s^k}(\text{resp. }b_{j,k}=\frac{b_j}{s^k}).$$
Table 1 is mostly taken from : the first two columns list the Kodaira fiber and the associated group (see Section 1); in the middle columns we write the order of vanishing of each $`a_i`$ along $`𝚺_1`$. Recall that our hypothesis (a flat Calabi–Yau fibration) imposes some restriction on the self-intersection of the ramification divisor (see the Remark after the Main Theorem 8.2). In the last column, we exhibit how the equation for $`𝚺_0mods`$ breaks into factors; the power $`r_j`$ which gives the multiplicity of the factor $`\beta _j`$ is indicated in the factorization in each case.
We have incorporated some necessary corrections to the Table from . First, the entry for $`I_{2k+1}`$, $`k1`$, with gauge group $`\mathrm{SU}(2k+1)`$ corresponds to the Weierstrass equation
$$y^2+a_1xy+a_{3,k}s^ky=x^3+a_{2,1}sx^2+a_{4,k+1}s^{k+1}x+a_{6,2k+1}s^{2k+1},$$
which has discriminant
$$\frac{1}{16}a_{1}^{}{}_{}{}^{4}(a_{1}^{}{}_{}{}^{2}a_{6,\mathrm{\hspace{0.17em}2}k+1}a_1a_{3,k}a_{4,k+1}+a_{3,k}^{}{}_{}{}^{2}a_{2,\mathrm{\hspace{0.17em}1}})s^{2k+1}\frac{1}{16}a_1^3a_{3,k}^3s^{3k}+O(s^{2k+2}).$$
Thus, the correct leading term in the local equation of $`𝚺_0`$ in this case (the “residual discriminant”) takes the form
$$a_1^3(a_1b_{8,2k+1}a_{3,k}^3),\text{if}k=1,$$
and
$$a_1^4b_{8,2k+1},\text{if}k>1$$
(not $`a_1^6a_{6,2k+1}`$ as was written in ).
Second, the residual discriminant in the case $`IV`$ (with gauge group $`\mathrm{SU}(3)`$) should read $`27a_{3,1}^4`$ rather than $`27a_{3,2}^4`$.
###### Remark I.2.
Following (see also Proposition 5.4) we see that if
$$𝚺_0𝚺_1=\{P_1^1,\mathrm{}P_1^{B_1},P_2^1,\mathrm{}P_2^{B_2}\},$$
then the local equation of $`𝚺_0`$ around $`P_i^{\mathrm{}}`$ does not depend on $`\mathrm{}`$, but only on $`i=1,2`$.
In Table 2 we list the local equation (l.e.) of $`𝚺_0`$ around $`P_1`$ and $`P_2`$. As usual, we denote by $`s=0`$ the divisor $`𝚺_1`$; $`t`$ is a convenient coordinate vanishing at $`P_i`$ and $`\gamma _i`$ is a suitable invertible function near $`\{s=t=0\}`$.
Our assumption on the existence of a smooth Calabi–Yau resolution imposes of $`𝚺_1`$ and $`𝚺_0`$ We write “NM” or “NSR” if the intersection type, as stated in Table 1 is not compatible with our hypothesis due to the singularities being non-minimal or having no small resolution.
In Table 3, $`h`$ denotes the Coxeter number of the group $`G`$, $`m`$ the multiplicity of $`𝚺_1`$ in the discriminant, and $`\mu (f)`$ (resp. $`\mu (g)`$) the vanishing of $`f`$ (resp. $`g`$) in equation (1.1) along $`𝚺_1`$ (see also Section 5).
In Table 4 we write, for each Kodaira type fiber and associated group, the coefficients needed to compute $`𝓡`$, as in Proposition 6.9. The general Kodaira type fiber over $`𝚺_1`$ degenerates over both $`P_i`$ at the intersection with $`𝚺_0`$. As in Table 2 we write “NM” or “NSR” if the intersection type, as stated in Table 1 is not compatible with our hypothesis. We describe the degenerate singular fibers: if they are of Kodaira type we use Kodaira’s notation. Note that these are not necessarily the Kodaira type of the general Weierstrass surface passing through the degenerate fiber; for example in the case of $`G=E_7`$ ($`III^{}`$), the degenerate fiber is again of type $`III^{}`$, but the general Weierstrass surface has a $`II^{}`$ singularity (see also Section 8.3). These distinctions are important in computing $`𝓡`$ as in Theorem 8.2.
The fibers of non-Kodaira type are the branch points of an outer automorphism of the group; we denote these with “br.”.
## Appendix II: The entries in the above Tables for $`G=\mathrm{SU}(2k),k2`$ and $`I_{2k}`$ fiber type.
We illustrate the pattern of computations needed to compile the Tables in Appendix I with the specific example $`G=\mathrm{SU}(2k)`$.
The generalized Weierstrass equation has the form:
$$y^2+a_1xy=x^3+a_2sx^2+a_4s^kx+a_6s^{2k}.$$
$`b_2=a_{1}^{}{}_{}{}^{2}+4a_2s`$, $`b_4=2a_4s^k`$, $`b_6=4a_6s^{2k}`$,
$`b_8=[(a_{1}^{}{}_{}{}^{2}+4a_2s)a_6a_{4}^{}{}_{}{}^{2}]s^{2k},b_{8,2k}=(a_{1}^{}{}_{}{}^{2}+4a_2s)a_6a_{4}^{}{}_{}{}^{2}`$
$`f=\frac{1}{48}(b_{2}^{}{}_{}{}^{2}24b_4),g=\frac{1}{864}(b_{2}^{}{}_{}{}^{3}+36b_2b_4216b_6).`$
$`𝚺:`$ $`s^{2k}\{(a_{1}^{}{}_{}{}^{4}+16a_{2}^{}{}_{}{}^{2}s^2+8a_{1}^{}{}_{}{}^{2}a_2s)[(a_{1}^{}{}_{}{}^{2}+4a_2s)a_6a_{4}^{}{}_{}{}^{2}]+`$
$`8s^k[8a_{4}^{}{}_{}{}^{3}+272a_{6}^{}{}_{}{}^{2}s^k9a_4a_6(a_{1}^{}{}_{}{}^{2}+4a_2s)],`$
$`𝚺_0:`$ $`(a_{1}^{}{}_{}{}^{4}+16a_{2}^{}{}_{}{}^{2}s^2+8a_{1}^{}{}_{}{}^{2}a_2s)[(a_{1}^{}{}_{}{}^{2}+4a_2s)a_6a_{4}^{}{}_{}{}^{2}]+`$
$`8s^k[8a_{4}^{}{}_{}{}^{3}+272a_{6}^{}{}_{}{}^{2}s^k9a_4a_6(a_{1}^{}{}_{}{}^{2}+4a_2s)]\}.`$
At the points of intersections of $`𝚺_0`$ and $`𝚺_1`$, either $`a_1=0`$ ($`P_1^{\mathrm{}}`$) or $`b_{8,2k}=0`$ ($`P_2^{\mathrm{}}`$). (In the notation of Section 5, $`a_1=\beta _1`$.)
###### Remark II.1.
$`r_1=4`$ and $`r_2=1`$; there are $`B_1=K_B𝚺_1`$ points of $`P_1`$ type, and $`(8K_B2k𝚺_1)𝚺_1=B_2`$ points of $`P_2`$ type. The second condition follows from the first one, as $`𝚺_1𝚺_0=4B_1+B_1`$.
### II.1. Computing $`ϵ_1`$:
Let $`t=:a_1,s`$ be the local coordinates around a point $`P_{\mathrm{}}^1`$. (In the notation of Section 5, $`a_1=\beta _1`$.)
Then
$$𝚺_0:\gamma _0t^4+\gamma _1s^2+\gamma _2t^2s+\gamma _3s^k,$$
where $`\gamma _i`$ is invertible at $`s=t=0`$.
We can write
$$𝚺_0:\gamma _0t^{2k}+\gamma _1s^2=0$$
which defines an $`A_{2k1}`$ curve singularity. Since the blowup of an $`A_{2k1}`$ curve singularity yields an $`A_{2k3}`$ singularity, we have
$$(\mathrm{\#}\varphi ^1(P_1);\{\alpha _v^1\})=(2;2,\mathrm{}2),(k\text{ times});\text{ then }ϵ_1=2k2.$$
### II.2. Computing $`ϵ_2`$:
Since $`𝚺_0`$ is smooth around each point $`P_2`$
$$(\mathrm{\#}\varphi ^1(P_2);\{\alpha _v^2\})=(1;1);ϵ_2=1.$$
### II.3. Computing $`\mu (f,g)`$:
From the equations we see that $`f`$ and $`g`$ have a common zero along $`𝚺_1`$ when $`b_2=0`$, and there $`2K_B𝚺_1`$ such points. Now set
$$g^{}=\frac{b_2}{18}f+g=\frac{1}{72}(b_2b_4)\frac{1}{12}b_6.$$
Then $`\mu (f,g)=\mu (f,g^{})`$ \[14, Section 1\].
From the equation above we see that $`P𝚺_1`$ is a common zero of $`f`$ and $`g`$ if and only if $`a_1=0`$. As in II.1 we take $`t:=a_1,s`$ as the local coordinates around $`P`$.
$`\mu (f,g)`$ $`=\mathrm{dim}_{}[[s,t]]/(f,g^{}),\text{ where}`$
$`f`$ $`t^4+\gamma _1t^2s+\gamma _2s^2+\gamma _3s^k`$
$`g^{}`$ $`24s^k\{\gamma _4t+s^k)\},`$
for suitable invertible functions $`\gamma _i`$ (around $`s=t=0`$). Then \[14, Ex. 1.2.5\]
$$\mu (f,g)=6k.$$
### II.4. Computing $`\chi _{top}(X_{P_1})`$.
After $`\mathrm{}`$ blowups the Weierstrass equation becomes: $`y^2+a_1xy=x^3s^{\mathrm{}}+a_2sx^2+a_4xs^k\mathrm{}+a_6s^{2k2\mathrm{}}`$, and there are isolated singular points (nodes) on the fiber at $`P_1=0`$. These points can be blown up with small resolutions: the fiber over the points $`P_1`$ is of Kodaira type $`D_{2k}`$ and $`\chi _{top}(X_{P_1})=2k+2.`$
### II.5. Computing $`\chi _{top}(X_{P_2})`$.
$`𝚺_0`$ and $`𝚺_1`$ intersect transversally at $`P_2`$, and it is easy to see that the corresponding fiber $`X_{P_2}`$ is of type $`I_{2k+1}`$ and $`\chi _{top}(X_{P_2})=2k+1.`$
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# 1 Introduction
## 1 Introduction
According to the recently formulated Equivalence Principle (EP), all physical systems are equivalent under coordinate transformations 1.–5. It has been shown that the implementation of such a principle unequivocally leads to the Quantum Hamilton–Jacobi Equation (QHJE). The latter was first analyzed independently by Floyd in a series of remarkable papers 6. In 5. the formulation of Quantum Mechanics (QM) from the EP was extended to higher dimensions and to the relativistic case as well. This approach suggests that QM and General Relativity (GR) are two facets of the same medal 1.–5. In this letter we will argue that QM and GR are intimately related. In particular, we suggest that gravitation has a purely quantum mechanical origin.
An outcome of the formulation of QM based on the EP is that the term $`𝒲VE`$, with $`V`$ the potential and $`E`$ the energy of the system, corresponds to the inhomogeneous term in the transformation properties of the state with $`𝒲=𝒲^00`$ (see Refs. 1.–5.). It turns out that this term is of a purely quantum nature. A related aspect concerns the appearance of fundamental constants in the QHJE. In particular, the implementation of the EP leads to the introduction of universal length scales. This has an important consequence once we take into account that the quantum potential is always non–trivial. This is a result which follows from a rigorous analysis of the QHJE. Here and throughout this paper it is important to distinguish between the quantum potential arising in the approach adopted here and that in the Bohm theory of QM 7. 8. The two are not in general the same (see Refs. 6. 9. 1.-5.). In particular, it turns out that even in the case of $`𝒲^0`$, the corresponding quantum potential is far from being trivial. This key point is due to the fact that the quantum reduced action, or quantum Hamiltonian characteristic function, is always non–trivial. In particular we have
$$𝒮_0cnst,$$
(1)
which follows as a direct consequence of the EP 1.–5.
## 2 The equivalence postulate
Before proceeding, let us analyze further the EP. To understand the basic motivation for its formulation, let us consider, in the classical framework, two particles of mass $`m_A`$ and $`m_B`$ with relative velocity $`v`$ 4. For an observer at rest with respect to the particle $`A`$, the two systems have reduced actions $`𝒮_0^{clA}(q_A)=cnst`$ and $`𝒮_0^{clB}(q_B)=m_Bvq_B`$. In this case, setting $`𝒮_0^{clA}(q_A)=𝒮_0^{clB}(q_B)`$ defines a highly singular coordinate transformation. However, when the same system is described by an observer in a frame which is not at rest with respect to either $`A`$ and $`B`$, we have that equating the two reduced actions does not lead to such a singularity. Thus, this strong singularity disappears if the frame one uses to describe the systems $`A`$ and $`B`$ has a non–zero velocity with respect to both. For example, if the observer is in a frame moving with constant acceleration $`a`$ with respect to the systems $`A`$ and $`B`$, then
$$\stackrel{~}{𝒮}_0^{clA}(Q_A)=\frac{m_A}{3a}(2aQ_A)^{\frac{3}{2}},\stackrel{~}{𝒮}_0^{clB}(Q_B)=\frac{m_B}{3a}(v^2+2aQ_B)^{\frac{3}{2}},$$
(2)
where $`Q_A`$ ($`Q_B`$) is the coordinate of particle $`A`$ ($`B`$) in the accelerated frame. If in describing particle $`B`$ in the accelerated frame one uses the coordinate $`Q_A`$ defined by $`\stackrel{~}{𝒮}_0^{clA}(Q_A)=\stackrel{~}{𝒮}_0^{clB}(Q_B)`$, then the resulting dynamics coincides with the one of particle $`A`$, that is
$$\stackrel{~}{𝒮}_0^{clB}(Q_B(Q_A))=\stackrel{~}{𝒮}_0^{clA}(Q_A),$$
(3)
which shows that the system $`B`$, described in terms of the coordinate $`Q_A`$, coincides with the system $`A`$. Hence, in Classical Mechanics (CM), the equivalence under coordinate transformations requires choosing a frame in which no particle is at rest. The existence of a distinguished frame, the one at rest, seems peculiar as on general grounds what is equivalent under coordinate transformations in all frames should remain so even in the one at rest. This leads to postulate that it is always possible to connect two systems by a coordinate transformation. In other words, it is natural to require that given two systems with reduced actions $`𝒮_0`$ and $`𝒮_0^v`$, there always exists the “$`v`$–map” $`qq^v`$ defined by 1.–5.
$$𝒮_0^v(q^v)=𝒮_0(q).$$
(4)
The above example shows that the HJ formalism provides the natural setting to describe physical systems. The equivalence under coordinate transformations is somehow in the spirit of general relativity. This property of HJ theory is still more evident if one notes that the classical HJ equation itself is obtained by looking for the canonical transformation of the conjugate variables that leads to the trivial Hamiltonian $`H=0`$. Thus, according to CM, since all states are equivalent, in the sense of the canonical transformations, to the trivial one, there is a sort of EP. Our view is slightly different from the one considered in the framework of canonical transformations. Actually, it is just the above example, in which equivalence under coordinate transformations always exist except that in the case one considers the particle at rest, which suggests the following stronger concept of equivalence. Let us denote by $``$ the space of all possible states $`𝒲`$. The equivalence postulate reads 1.–5.
For each pair $`𝒲^a,𝒲^b`$, there exists a $`v`$–transformation such that
$$𝒲^a(q)𝒲_{}^{a}{}_{}{}^{v}(q^v)=𝒲^b(q^v).$$
(5)
It has been shown in 1.–5. that the implementation of the EP unequivocally leads to the quantum HJ equation in any dimension and in the relativistic case as well.
## 3 Fundamental constants and the quantum potential
Due to the structure of the QHJE we have that the quantum potential will in general depend on fundamental constants. Let us show how these constants arise. We first focus on the one–dimensional Quantum Stationary Hamilton–Jacobi Equation (QSHJE). This reads
$$\frac{1}{2m}\left(\frac{𝒮_0(q)}{q}\right)^2+V(q)E+\frac{\mathrm{}^2}{4m}\{𝒮_0,q\}=0,$$
(6)
where $`\{𝒮_0,q\}=𝒮_0^{\prime \prime \prime }/𝒮_0^{}3(𝒮_0^{\prime \prime }/𝒮_0^{})^2/2`$ denotes the Schwarzian derivative and $`Q=\frac{\mathrm{}^2}{4m}\{𝒮_0,q\}`$ is the quantum potential. The general real solution of (6) has the form
$$e^{\frac{2i}{\mathrm{}}𝒮_0\{\delta \}}=e^{i\alpha }\frac{w+i\overline{\mathrm{}}}{wi\mathrm{}},$$
(7)
where $`w=\psi ^D/\psi R`$ and $`(\psi ^D,\psi )`$ are two real linearly independent solutions of the associated Schrödinger equation. Furthermore, we have $`\delta =\{\alpha ,\mathrm{}\}`$, with $`\alpha R`$ and $`\mathrm{}=\mathrm{}_1+i\mathrm{}_2`$ some integration constants ($`\overline{\mathrm{}}`$ denoting the complex conjugate of $`\mathrm{}`$). Observe that $`\mathrm{}_10`$, which is equivalent to having $`𝒮_0cnst`$, is a necessary condition to define the term $`\{𝒮_0,q\}`$ in the QSHJE.
There is a simple reason why fundamental constants should be hidden in $`\mathrm{}`$. To see this, consider the Schrödinger equation in the trivial case $`𝒲^0(q^0)0`$, that is $`_{q^0}^2\psi =0`$. Two linearly independent solutions are $`\psi ^D=q^0`$ and $`\psi =1`$. Now a basic aspect of the formulation is manifest duality between real pairs of linearly independent solutions 1.–5. This is a fact which is strictly related to the Legendre duality first observed in 10. and further investigated in 11.–14. Thus, whereas in the standard approach one usually considers only one solution of the Schrödinger equation, i.e. the wave–function itself, in the present formulation both $`\psi ^D`$ and $`\psi `$ enter the relevant formulas. This leads to expressions containing linear combinations of $`\psi ^D`$ and $`\psi `$, typically $`\psi ^D+i\mathrm{}\psi `$ that for $`\psi ^D=q^0`$ and $`\psi =1`$ reads $`q^0+i\mathrm{}_0`$, so $`\mathrm{}_0\mathrm{}`$ should have the dimensions of a length. The fact that $`\mathrm{}`$ has the dimensions of a length is true for any state. This follows from the observation that the ratio $`w=\psi ^D/\psi `$ is a Möbius transformation of the trivializing map transforming any state to $`𝒲^0`$ 1.–5. Hence $`w`$, and therefore $`\mathrm{}`$, has the dimensions of a length.
Since $`\mathrm{}_0`$ enters the QSHJE with $`𝒲^00`$, the system does not provide any dimensionful quantity. This implies that we have to introduce some fundamental lengths. Let us show this in some detail. The reduced action $`𝒮_0^0`$ corresponding to the state $`𝒲^0`$ is
$$e^{\frac{2i}{\mathrm{}}𝒮_0^0\{\delta \}}=e^{i\alpha }\frac{q^0+i\overline{\mathrm{}}_0}{q^0i\mathrm{}_0},$$
(8)
and the conjugate momentum $`p_0=_{q^0}𝒮_0^0`$ has the form
$$p_0=\pm \frac{\mathrm{}(\mathrm{}_0+\overline{\mathrm{}}_0)}{2|q^0i\mathrm{}_0|^2}.$$
(9)
A property of $`p_0`$ is that it vanishes only for $`q^0\pm \mathrm{}`$. Furthermore, $`|p_0|`$ reaches its maximum at $`q^0=\mathrm{Im}\mathrm{}_0`$
$$|p_0(\mathrm{Im}\mathrm{}_0)|=\frac{\mathrm{}}{\mathrm{Re}\mathrm{}_0}.$$
(10)
Since $`\mathrm{Re}\mathrm{}_00`$, $`p_0`$ is always finite. Thus, $`\mathrm{Re}\mathrm{}_00`$ provides a sort of ultraviolet cutoff. This is a property which extends to arbitrary states. Actually, the conjugate momentum has the form
$$p=\frac{\mathrm{}W(\mathrm{}+\overline{\mathrm{}})}{2\left|\psi ^Di\mathrm{}\psi \right|^2},$$
(11)
where $`W=\psi ^{}\psi ^D\psi ^D^{}\psi `$ is the Wronskian. Since $`W`$ is a non–vanishing constant, it follows that $`\psi ^D`$ and $`\psi `$ cannot have common zeroes, and by $`\mathrm{Re}\mathrm{}0`$ we see that $`p`$ is finite $`qR`$. Therefore, the EP implies an ultraviolet cutoff on the conjugate momentum.
In Refs. 2. and 4. it has been shown that fundamental constants also arise in considering the classical limit. In particular, one first considers
$$\underset{\mathrm{}0}{lim}p_0=0,$$
(12)
and note that $`\mathrm{Im}\mathrm{}_0`$ in (9) can be absorbed by a shift of $`q^0`$. Hence, in (12) we can set $`\mathrm{Im}\mathrm{}_0=0`$ and distinguish the cases $`q^00`$ and $`q^0=0`$. From (12)
$$p_0{}_{\stackrel{}{\mathrm{}0}}{}^{}\{\begin{array}{cc}\mathrm{}^{\gamma +1},\hfill & q_00,\hfill \\ \mathrm{}^{1\gamma },\hfill & q_0=0,\hfill \end{array}$$
(13)
where $`1<\gamma <1`$ with $`\gamma `$ defined by $`\mathrm{Re}\mathrm{}_0{}_{\stackrel{}{\mathrm{}0}}{}^{}\mathrm{}_{}^{\gamma }`$. The are not many fundamental lengths in nature. In particular, we note that a fundamental length satisfying this condition on the power of $`\mathrm{}`$ is the Planck length $`\lambda _p=\sqrt{\mathrm{}G/c^3}`$, while the Compton length is excluded by the condition $`\gamma <1`$. Also, as we will see in considering the $`E0`$ and $`\mathrm{}0`$ limits for the free particle of energy $`E`$, the natural choice is just the Planck length. With this choice of $`\mathrm{Re}\mathrm{}_0`$ the maximum of $`|p_0|`$ is
$$|p_0(\mathrm{Im}\mathrm{}_0)|=\sqrt{\frac{c^3\mathrm{}}{G}}.$$
(14)
Setting $`\mathrm{Im}\mathrm{}_0=0`$ and $`\mathrm{Re}\mathrm{}_0=\lambda _p`$, the quantum potential associated to the trivial state $`𝒲^0`$ is
$$Q^0=\frac{\mathrm{}}{4m}\{𝒮_0^0,q^0\}=\frac{\mathrm{}^3G}{2mc^3}\frac{1}{|q^0i\lambda _p|^4}.$$
(15)
There are two basic aspects in this expression. Firstly, the gravitational constant $`G`$ results from ensuring consistency with the classical limit. We saw that this arises naturally as a consistency condition. Furthermore, $`Q^0`$ is negative definite. Thus, even if we are still considering a one–dimensional system, we are starting to see some motivation for the emergence of the gravitational interaction. In particular, note that this analysis is essentially the same of the one in three dimensional space as in the case of a free particle we can consider a reduced action of the form $`𝒮_0(x)+𝒮_0(y)+𝒮_0(z)`$, which can always be chosen as a possible solution of the QHSJE when the potential has the form $`V(x,y,z)=V_1(x)+V_2(y)+V_3(z)`$. We also note that the fact that we are in the framework of non–relativistic QM, does not exclude the appearance of $`c`$ in the integration constants of the QSHJE (an example of the appearance of $`c`$ in QM is the non–relativistic treatment of an electron in a magnetic field).
The appearance of fundamental constants can also be seen by considering the $`\mathrm{}0`$ and $`E0`$ limits 2. for the conjugate momentum of a free particle of energy $`E`$
$$p_E=\pm \frac{\mathrm{}(\mathrm{}_E+\overline{\mathrm{}}_E)}{2\left|k^1\mathrm{sin}(kq)i\mathrm{}_E\mathrm{cos}(kq)\right|^2},$$
(16)
where $`k=\sqrt{2mE}/\mathrm{}`$ and $`\mathrm{}_E`$ is the integration constant of the QSHJE. We should require that (see Refs. 2. and 4.)
$$\underset{\mathrm{}0}{lim}p_E=\pm \sqrt{2mE},$$
(17)
and
$$\underset{E0}{lim}p_E=p_0=\pm \frac{\mathrm{}(\mathrm{}_0+\overline{\mathrm{}}_0)}{2|qi\mathrm{}_0|^2}.$$
(18)
However, we see that the term $`\mathrm{}_E\mathrm{cos}(kq)`$ in Eq.(16) is ill–defined in the $`\mathrm{}0`$ limit, a problem which has been recently considered also by Floyd 15. Thus, the existence of the classical limit implies some condition on $`\mathrm{}_E`$. In particular, in order to reach the classical value $`\sqrt{2mE}`$ in the $`\mathrm{}0`$ limit, the quantity $`\mathrm{}_E`$ should depend on $`E`$. In Refs. 2. and 4. it has been shown that
$$\mathrm{}_E=k^1e^{\alpha (x_p^1)}+e^{\beta (x_p)}\mathrm{}_0,$$
(19)
where $`x_p=k\lambda _p=\sqrt{2mEG/\mathrm{}c^3}`$ and
$$\alpha (x_p^1)=\underset{k1}{}\alpha _kx_p^k,\beta (x_p)=\underset{k1}{}\beta _kx_p^k.$$
(20)
It follows that
$$p_E=\pm \frac{2k^1\mathrm{}e^{\alpha (x_p^1)}+\mathrm{}e^{\beta (x_p)}(\mathrm{}_0+\overline{\mathrm{}}_0)}{2\left|k^1\mathrm{sin}(kq)i\left(k^1e^{\alpha (x_p^1)}+e^{\beta (x_p)}\mathrm{}_0\right)\mathrm{cos}(kq)\right|^2}.$$
(21)
The function $`\alpha (x_p^1)`$ is constrained by the conditions
$$\underset{\mathrm{}0}{lim}e^{\alpha (x_p^1)}=1,\underset{E0}{lim}E^{1/2}e^{\alpha (x_p^1)}=0,$$
(22)
whereas for $`\beta (x_p)`$ we have
$$\underset{\mathrm{}0}{lim}\mathrm{}^1e^{\beta (x_p)}\mathrm{}_0=0.$$
(23)
One of the conditions we used to derive the above formulas concerns the existence of the classical limit. In this context we should observe that it may be that the classical expressions themselves may contain further terms which do not depend on $`\mathrm{}`$. As an example, observe that two free particles should always contain the gravitational potential as intrinsic interaction. This seems to be connected with the residual indeterminacy discussed in 15. The aim of the present paper is to investigate the possibility that such interaction may be a consequence of the quantum potential.
The appearance of the Planck scale in the hidden constants has been considered in Refs. 2. and 4. This seems related to the ’t Hooft’s approach 16. Possible connections have been considered by Floyd 15. and in 5.
## 4 The cocycle condition and the quantum nature of interactions
Let us further consider the nature of the EP itself. We introduce the notation
$$J_{ki}=\frac{q_i}{q_k^v},$$
(24)
and
$$(p^v|p)=\frac{\underset{k}{}p_k^{v^2}}{_kp_k^2}=\frac{p^tJ^tJp}{p^tp}.$$
(25)
The only possibility to reach any other state $`𝒲^v0`$ starting from $`𝒲^0`$ is that it transforms with an inhomogeneous term 1.–5.
$$𝒲^v(q^v)=(p^v|p^a)𝒲^a(q^a)+(q^a;q^v),$$
(26)
and
$$Q^v(q^v)=(p^v|p^a)Q^a(q^a)(q^a;q^v),$$
(27)
where $`(q^a;q^v)`$ denotes a still undefined function which depends on $`q^a`$ and $`q^v`$. Let us denote by $`a,b,c,\mathrm{}`$ a set of different $`v`$–transformations. Comparing
$$𝒲^b(q^b)=(p^b|p^a)𝒲^a(q^a)+(q^a;q^b)=(q^0;q^b),$$
(28)
with the same formula with $`q^a`$ and $`q^b`$ interchanged we have
$$(q^b;q^a)=(p^a|p^b)(q^a;q^b).$$
(29)
More generally, comparing
$$𝒲^b(q^b)=(p^b|p^c)𝒲^c(q^c)+(q^c;q^b)=(p^b|p^a)𝒲^a(q^a)+(p^b|p^c)(q^a;q^c)+(q^c;q^b),$$
(30)
with (28) we obtain the basic cocycle condition
$$(q^a;q^c)=(p^c|p^b)\left[(q^a;q^b)(q^c;q^b)\right],$$
(31)
which is the essence of the EP. In particular, this condition unequivocally leads to determine the correction to the CHJE. In doing this, one shows that Eq.(31) implies a basic Möbius invariance of $`(q^a;q^b)`$. The $`𝒲^00`$ state plays a special role. Setting $`𝒲^a=𝒲^0`$ in Eq.(26) yields $`𝒲^v(q^v)=(q^0;q^v)`$. Thus, in general
$$𝒲(q)=(q^0;q),$$
(32)
so that, according to the EP (5), all states correspond to the inhomogeneous part in the transformation of the $`𝒲^0`$ state induced by some $`v`$–map. Since the inhomogeneous part has a purely quantum origin, we conclude that the Equivalence Postulate implies that interactions have a purely quantum origin.
The role of the quantum potential as responsible for interactions can be made clearer from the observation that the EP implies
$$𝒲^v(q^v)+Q^v(q^v)=(p^v|p)\left(𝒲(q)+Q(q)\right).$$
(33)
Then, taking $`𝒲=𝒲^00`$ and omitting the superscript $`v`$, we have
$$𝒲(q)=(p|p^0)Q^0(q^0)Q(q),$$
(34)
showing that any potentials can be expressed in quantum terms.
In 5. it has been observed that there is a hidden antisymmetric tensor in QM which arises from the continuity equation. We also note that in the one–dimensional case, the freedom deriving from the underlying hidden tensor one meets in the higher dimensional case reflects itself in the appearance of the integration constants. These are related to the $`SL(2,C)`$ symmetry
$$e^{2i𝒮_0/\mathrm{}}\frac{Ae^{2i𝒮_0/\mathrm{}}+B}{Ce^{2i𝒮_0/\mathrm{}}+D},$$
(35)
of the equation
$$\{e^{2i𝒮_0/\mathrm{}},q\}=4m𝒲/\mathrm{}^2,$$
(36)
which is equivalent to the QSHJE (6). In particular, as we said, there is a complex integration constant $`\mathrm{}`$ which is missing in the Schrödinger equation. Changing this constant corresponds to a Möbius transformation (35). While this leaves $`𝒲`$ unchanged, it mixes the quantum potential and the kinetic term. Thus the quantum potential is essentially parameterized by $`SL(2,C)`$ transformations in which the constants $`A,B,C`$ and $`D`$ depend, by dimensional analysis and consistency of relevant limits considered above, on fundamental constants. We may expect that these constants and the above Möbius transformations in three and four dimension (for the relativistic generalization) should be related to fundamental interactions.<sup>1</sup><sup>1</sup>1It is worth mentioning that the geometrical building block of string theory, which also explains why string theory includes gravity, is the thrice punctured Riemann sphere. The latter can be characterized just by the basic $`SL(2,C)`$ Möbius symmetry related to the arbitrariness of the position of the punctures.
## 5 The two–particle model
The above investigation suggests considering that the quantum potential $`Q`$ is at the origin of the interactions. Thus, it may be that the constants defining $`Q`$ depend on the intrinsic properties of the particles. This would lead to different possible forms of $`Q`$ and therefore of the admissible interactions. As we saw, there are subtle questions concerning the classical limit. Similarly, one should consider the relativistic case as it may lead to results which may remain hidden if considered directly in the non–relativistic case. Similarly, at least in the case of gravitational interaction, one should consider the analysis of macroscopic objects to take into account possible collective effects. Nevertheless, the above suggestions indicate that it is worth studying the case of two free particles and then looking at the possible structure of the quantum potential.
In the case of two free particles of energy $`E`$ and masses $`m_1`$ and $`m_2`$, the QSHJE reads
$$\frac{1}{2m_1}(_1𝒮_0)^2+\frac{1}{2m_2}(_2𝒮_0)^2E\frac{\mathrm{}^2}{2m_1}\frac{\mathrm{\Delta }_1R}{R}\frac{\mathrm{}^2}{2m_2}\frac{\mathrm{\Delta }_2R}{R}=0.$$
(37)
The continuity equation is
$$\frac{1}{m_1}_1(R^2_1𝒮_0)+\frac{1}{m_2}_2(R^2_2𝒮_0)=0.$$
(38)
Next, we set
$$r=r_1r_2,r_{c.m.}=\frac{m_1r_1+m_2r_2}{m_1+m_2},m=\frac{m_1m_2}{m_1+m_2},$$
(39)
where $`r_1`$ and $`r_2`$ are the ray vectors of the two particles. With respect to the new variables the equations (37) and (38) have the form
$$\frac{1}{2(m_1+m_2)}(_{r_{c.m.}}𝒮_0)^2+\frac{1}{2m}(𝒮_0)^2E\frac{\mathrm{}^2}{2(m_1+m_2)}\frac{\mathrm{\Delta }_{r_{c.m.}}R}{R}\frac{\mathrm{}^2}{2m}\frac{\mathrm{\Delta }R}{R}=0,$$
(40)
$$\frac{1}{m_1+m_2}_{r_{c.m.}}(R^2_{r_{c.m.}}𝒮_0)+\frac{1}{m}(R^2𝒮_0)=0,$$
(41)
where $``$ ($`_{r_{c.m.}}`$) and $`\mathrm{\Delta }`$ ($`\mathrm{\Delta }_{r_{c.m.}}`$) are the gradient and Laplacian with respect to the components of the vector $`r`$ ($`r_{c.m.}`$). These equations can be decomposed into the equations for the center of mass $`r_{c.m.}`$ and those for the relative motion. We will concentrate on the latter. It satisfies the QSHJE
$$\frac{1}{2m}(𝒮_0)^2E\frac{\mathrm{}^2}{2m}\frac{\mathrm{\Delta }R}{R}=0,$$
(42)
and the continuity equation
$$(R^2𝒮_0)=0.$$
(43)
In 5 it has been stressed that the continuity equation implies
$$R^2_i𝒮_0=ϵ_i^{i_2\mathrm{}i_D}_{i_2}F_{i_3\mathrm{}i_D},$$
(44)
where $`F`$ is a $`(D2)`$–form. In the 3D case $`R^2_i𝒮_0`$ is the curl of a vector that we denote by $`B`$
$$𝒮_0=R^2\times B.$$
(45)
The QSHJE (42) reduces to the “canonical form”
$$j^2=\mathrm{}^2R^3\mathrm{\Delta }R+2mER^4,$$
(46)
where $`j^2j_kj^k`$ with
$$j=\times B.$$
(47)
Using the identity $`(a\times b)(c\times d)=(ac)(bd)(ad)(bc)`$, Eq.(46) reads
$$\mathrm{\Delta }B^2(B)^2=\mathrm{}^2R^3\mathrm{\Delta }R+2mER^4.$$
(48)
It is worth stressing that $`j`$ resembles the usual current. However, besides the mass term, we stress again that here $`R`$ and $`𝒮_0`$ are not in general the ones one obtains identifying $`Re^{\frac{i}{\mathrm{}}𝒮_0}`$ with the wave–function. Nevertheless, by construction we have that $`\psi =Re^{\frac{i}{\mathrm{}}𝒮_0}`$ solves the Schrödinger equation. Thus, we have
$$j=\frac{i\mathrm{}}{2}(\psi \overline{\psi }\overline{\psi }\psi ),$$
(49)
that, in the case in which $`\psi `$ is the wave–function, coincides, upon dividing by $`m`$, with the usual quantum mechanical current.
We have seen that two free particles have a non–trivial quantum potential whose structure depends on the field $`B`$. In the following we will use the above results in order to investigate whether this potential may in fact have a gravitational leading behavior.
## 6 Gravitational interaction and quantum potential
After summarizing the main results so far, we will write down the differential equation that the quantum potential should satisfy in order to obtain the gravitational interaction. We will then investigate in some detail such an equation.
The aim of the previous sections was to show the main aspects suggesting that the quantum potential is at the origin of fundamental interactions. Even if these aspects have been discussed in detail, it is useful to collect them together before formulating the hypothesis and then deriving the relevant equations.
1. The EP implies that the reduced action is always non–trivial. In particular, this is true also for the free particle of vanishing energy. Furthermore, if $`\psi \overline{\psi }`$, such as in the case of the wave–function for bound states, then $`\psi =R\left(Ae^{\frac{i}{\mathrm{}}𝒮_0}+Be^{\frac{i}{\mathrm{}}𝒮_0}\right)`$, with $`\psi \overline{\psi }`$ giving $`|A|=|B|`$. Thus there is no track of the condition $`𝒮_0=cnst`$. On the other hand, this cannot be a solution of the QSHJE and would give an inconsistent classical limit. Remarkably, this answers the objections concerning the classical limit posed by Einstein. He just noticed that for a particle in a box the identification of the wave–function with $`Re^{\frac{i}{\mathrm{}}𝒮_0}`$ gives $`𝒮_0=cnst`$ and this cannot reproduce, in the $`\mathrm{}0`$ limit, the non–trivial $`𝒮_0^{cl}`$. This result has been previously derived by Floyd in a series of important papers 6. Related aspects have also been considered in the interesting papers by Reinisch<sup>2</sup><sup>2</sup>2I am grateful to G. Reinisch who informed me that the argument about the unphysical $`\mathrm{}0`$ limit was explicitly used by Einstein (see pg.243 of Holland’s book 8.). 9.
2. This property of the reduced action implies the existence of an intrinsic potential energy which, like the rest mass of special relativity, is universal. In particular, the quantum potential is always non–trivial. This is different from the standard approach where there are examples in which $`Q=0`$ so that the QHJE would coincide with the classical one.
3. The existence of the classical limit implies that the quantum potential depends, through the hidden initial conditions coming from the QSHJE, on fundamental length scales which in turn depend on $`\mathrm{}`$. It is a basic fact that these initial conditions are missing in the Schrödinger equation. In particular, the emergence of the Planck length, and therefore of Newton’s constant, arises from considering the classical limit for the free particle of vanishing energy.
4. It can be seen in the formulation that the quantum potential provides particle’s response to an external perturbation. For example, in the case of tunnelling, the attractive nature of the quantum potential guarantees the reality of the conjugate momentum and therefore of the velocity field $`v=1/_Epp/m`$ (see Refs. 6. and 4.). More precisely, inside the barrier the quantum potential decreases its value in such a way that $`(_q𝒮_0)^2`$ remains positive definite. As a consequence, the role of this internal energy, which is a property of all forms of matter, should manifest itself through effective interactions depending on the above fundamental constants.
5. The fundamental implication of the EP is the cocycle condition (31). In particular, from this condition, one obtains an expression for the interaction terms which is purely of quantum origin.
6. The fact that QM arises from an EP which is reminiscent of Einstein’s EP strongly indicates a deep relation between gravitation and QM itself.
The most characteristic property of the quantum potential is its universal nature: it is a property possessed by all forms of matter. On the other hand, we know that such a property is the one characterizing gravity. Therefore, if we write down the classical equations of motion for a pair of particles, we should always include, already at the classical level, the gravitational interaction. Furthermore, the quantum potential for a free particle is negative definite. This should be compared with the attractive nature of gravity.
## 7 The quantum potential with the gravitational potential as a leading term
The above remarks suggest formulating the hypothesis that the quantum potential is in fact at the origin of gravitation. Thus we look for solutions of the QSHJE leading to the classical HJ equation for the gravitational interaction. In particular, we should investigate whether in the case of two free particles the quantum potential
$$Q=\frac{\mathrm{}^2}{2m}\frac{\mathrm{\Delta }R}{R},$$
(50)
admits the form
$$Q=V_G,$$
(51)
with $`V_G`$ reducing to the Newton potential in the $`\mathrm{}0`$ limit
$$\underset{\mathrm{}0}{lim}V_G=G\frac{m_1m_2}{r}.$$
(52)
If such a solution exists then, in the limit $`\mathrm{}0`$, Eq.(42) corresponds to the HJ equation for the gravitational potential
$$\frac{1}{2m}(𝒮_0^{cl})^2G\frac{m_1m_2}{r}E=0.$$
(53)
Summarizing, the above problem corresponds to finding all the possible $`R`$ satisfying the equation
$$\frac{\mathrm{}^2}{2m}\frac{\mathrm{\Delta }R}{R}=V_G=G\frac{m_1m_2}{r}+𝒪(\mathrm{}),$$
(54)
where the higher order terms $`𝒪(\mathrm{})`$ will generally depend on $`r`$, such that $`R`$ and $`𝒮_0`$ satisfy Eqs.(42) and (43). Let us consider the set $`=\{R|sol.of(\text{54})\}`$. The above problem is equivalent to find the set $`=\{B|sol.of(\text{48})withR\}`$ (recall that if $`R`$ and $`B`$ solve Eq.(48), then $`𝒮_0=R^2\times B`$ is solution of the QSHJE and of the continuity equation). It follows that the set of possible potentials with gravitational behavior $`r^1`$ is given by
$$𝒱_G=\left\{\frac{\mathrm{}^2}{2m}\frac{\mathrm{\Delta }R}{R}\right|R_G\},$$
(55)
where $`_G=\{R|R,Bexists\}`$. In other words, we have to find all the possible $`R`$ satisfying (54) and then restricting to those for which there exists a field $`B`$ satisfying (48). This would fix the set of admissible potentials $`𝒱_G`$ to be investigated. Note that the fact that the higher order terms in (54) are not fixed implies that $``$ has infinitely many elements. This set identifies infinitely many equations of the kind (48), one for each $`R`$. Thus, on general grounds, one should expect that the set $``$, and therefore $`𝒱_G`$, be non–trivial.
## 8 The spherical case
While an adequate treatment of the above problem will be considered in a future publication, here we consider some related preliminary aspects. By introducing the $`B`$ field we saw that it should be possible to a find a solution to the two–particle model. However, a more effective way of considering such a problem seems to reformulate it as follows. First we note that by (42) and (54) we have that $`𝒮_0`$ should satisfy the equation
$$\frac{1}{2m}(𝒮_0)^2=E+G\frac{m_1m_2}{r}+𝒪(\mathrm{}).$$
(56)
Thus, instead of finding first the possible $`R`$, it seems convenient to solve Eq.(56) which looks simpler than Eq.(54). A general solution of this equation would involve terms depending also on $`\theta `$ and $`\varphi `$. However, the simplest situation is when $`𝒮_0`$ is a function of $`r`$. In this case $`𝒮_0=\widehat{r}_r𝒮_0(r)`$, where $`\widehat{r}`$ is the unit vector along $`r`$. Eq.(56) becomes
$$\frac{1}{2m}(_r𝒮_0)^2=E+G\frac{m_1m_2}{r}+𝒪(\mathrm{}),$$
(57)
and the continuity equation reads $`(R^2\widehat{r}_r𝒮_0)=0`$, giving
$$R=\frac{1}{r\sqrt{_r𝒮_0}}.$$
(58)
Since the radial part of the Laplacian is $`r^1_r^2r`$, we have that the QSHJE (42) becomes
$$\frac{1}{2m}(_r𝒮_0)^2E+\frac{\mathrm{}^2}{4m}\{𝒮_0,r\}=0.$$
(59)
Formally this equation is the one–dimensional QSHJE for a free particle on the non–negative part of the real axis. Therefore, by Eq.(16) we have
$$_r𝒮_0=\pm \frac{\mathrm{}(\mathrm{}_E+\overline{\mathrm{}}_E)}{2\left|k^1\mathrm{sin}(kr)i\mathrm{}_E\mathrm{cos}(kr)\right|^2}.$$
(60)
To establish the right asymptotic we should handle the indeterminacy discussed above and eliminated by a suitable choice of the constant $`\mathrm{}_E`$. In particular, while in the previous case the structure of $`\mathrm{}_E`$ was fixed by requiring that $`p_E\pm \sqrt{2mE}`$ as $`\mathrm{}0`$, we should now investigate the full functional structure of the right hand side of (60) at the different scales defined by the parameters $`\mathrm{}_E`$, $`\mathrm{}`$, $`m`$ and $`E`$.
We now consider the general case by adding to $`𝒮_0`$ the dependence on $`\theta `$ and $`\varphi `$ and then studying the possible appearance of the $`r^1`$ term. Setting $`\psi =Re^{\frac{i}{\mathrm{}}𝒮_0}`$, which is a solution of the Schrödinger equation, we have $`𝒮_0=\frac{\mathrm{}}{2i}\mathrm{ln}(\psi /\overline{\psi })`$, so that
$$(𝒮_0)^2=\frac{\mathrm{}^2}{4|\psi |^4}\underset{j=1}{\overset{3}{}}(\overline{\psi }_j\psi \psi _j\overline{\psi })^2,$$
(61)
where $`_1=_x`$, $`_2=_y`$ and $`_3=_z`$. Since $`\psi `$ solves the free Schrödinger equation, we have
$$\psi =\underset{l=0}{\overset{\mathrm{}}{}}\underset{m=l}{\overset{l}{}}\underset{j=1}{\overset{2}{}}c_{lmj}R_{klj}(r)Y_{lm}(\theta ,\varphi ),$$
(62)
where the $`Y_{lm}(\theta ,\varphi )`$ denote the spherical harmonics and
$$R_{kl1}=(1)^l2\frac{r^l}{k^l}\left(\frac{1}{r}_r\right)^l\frac{\mathrm{sin}kr}{r},R_{kl2}=(1)^l2\frac{r^l}{k^l}\left(\frac{1}{r}_r\right)^l\frac{\mathrm{cos}kr}{r}.$$
(63)
These are linearly independent solutions of the radial part of the Schrödinger equation
$$R_{klj}^{\prime \prime }+\frac{2}{r}R_{klj}^{}+\left[k^2\frac{l(l+1)}{r^2}\right]R_{klj}=0.$$
(64)
We are studying this equation at $`r>0`$. In this respect note that the singularity of $`R_{kl2}`$ at $`r=0`$ would give a term $`\delta (r)`$ in the right hand side of (64). In spherical coordinates we have
$$(𝒮_0)^2=\frac{\mathrm{}^2}{4|\psi |^4}\left[(\overline{\psi }_r\psi \psi _r\overline{\psi })^2+\frac{1}{r^2}(\overline{\psi }_\theta \psi \psi _\theta \overline{\psi })^2+\frac{1}{r^2\mathrm{sin}^2\theta }(\overline{\psi }_\varphi \psi \psi _\varphi \overline{\psi })^2\right].$$
(65)
The properties of this expression will be considered elsewhere. However, as a preliminary step, we consider the first term in the square bracket in (65). Note that
$$=(_r,\frac{1}{2r}(e^{i\varphi }l_+e^{i\varphi }l_{}),\frac{i}{r\mathrm{sin}\theta }l_z),$$
(66)
where $`l_\pm =l_x\pm il_y=e^{\pm i\varphi }(_\theta +i\mathrm{cot}\theta _\varphi )`$, with $`l_x`$, $`l_y`$ and $`l_z`$ denoting the components of the angular momentum operator. Since $`l_+Y_{lm}=a_{lm}Y_{lm+1}`$ and $`l_{}Y_{lm}=a_{lm1}Y_{lm1}`$, with $`a_{lm}\sqrt{(l+m+1)(lm)}`$, we have
$$\psi =\underset{\{lmj\}}{}(c_{lmj}R_{klj}^{}Y_{lm},\frac{1}{2r}c_{lmj}R_{klj}(e^{i\varphi }a_{lm}Y_{lm+1}e^{i\varphi }a_{lm1}Y_{lm1}),\frac{i}{r\mathrm{sin}\theta }c_{lmj}R_{klj}mY_{lm}),$$
(67)
where $`_{\{lmj\}}_{l=0}^{\mathrm{}}_{m=l}^l_{j=1}^2`$. Finally, using $`R_{klj}^{}=lr^1R_{klj}kR_{kl+1j}`$, we have
$$(\overline{\psi }_r\psi \psi _r\overline{\psi })^2=4\mathrm{}^2\left(\underset{\{lmj\}}{}\underset{\{l^{}m^{}j^{}\}}{}\overline{c}_{l^{}m^{}j^{}}R_{kl^{}j^{}}\overline{Y}_{l^{}m^{}}c_{lmj}(lr^1R_{klj}kR_{kl+1j})Y_{lm}\right).$$
(68)
## 9 Conclusions
Let us conclude by observing that the aim of the present investigation is to propose a possible quantum origin of the gravitational interaction. In particular, we made a preliminary investigation of the problem of finding the set $`𝒱_G`$ of potentials with gravitational leading term originated by the quantum potential of two free particles. The general solution seems to be of mathematical interest and will be considered in a future publication. In this context we would like to mention the Schrödinger–Newton equation 17. which concerns a problem reminiscent of the one introduced in this paper. Finally, we would like to mention that geometrical aspects related to the quantum HJ equation have been considered also in 18. and references therein.
Acknowledgements. It is a pleasure to thank the anonymous Referees for interesting comments and D. Bellisai, G. Bertoldi, R. Carroll, A.E. Faraggi, E.R. Floyd, J.M. Isidro, P.A. Marchetti, M. Mariotti, L. Mazzucato, P. Pasti, G. Reinisch, P. Sergio and M. Tonin, for stimulating discussions. Work supported in part by the European Commission TMR programme ERBFMRX–CT96–0045.
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# Two-particle correlation from a relativistic fluid with a first order phase transition
## 1 INTRODUCTION
In the physics of relativistic heavy ion collisions aiming at studying the QCD phase transition of hot and dense matter, two-particle correlation is widely used to investigate the characteristics of the space-time evolution in the collision processes . It is well known as Hanbury-Brown Twiss effect that the two-particle intensity correlation function contains the information on the “size” of the source of emitted particles due to a quantum symmetry. In recent experiments, the sizes are estimated from the correlation function through Gaussian fitting and the physical meaning of the sizes depends on the fitting function. For an azimuthally symmetric source, the Yano-Koonin-Podgoretskiĭ(YKP) parametrization is used as one of the fitting functions. In this parametrization, three of fitting parameters can be regarded as source sizes directly. This advantage gives us a chance to observe a signature of quark-gluon matter because large emission duration is expected in the case of a first order phase transition . In the case of the relativistic heavy ion collisions, the reaction is highly dynamic; the particle source cannot be considered to be static one. So we analyze the size parameters for pions quantitatively based on a hydrodynamic model which takes into account the longitudinal and transverse expansion and a first order phase transition .
## 2 FORMALISM
From a hydrodynamic point of view, we consider the pion source as a thermalized one and a completely chaotic source. Hence, the two-particle correlation function is given as
$$C(𝐤_\mathrm{𝟏},𝐤_\mathrm{𝟐})=1+\frac{|I(𝐤_\mathrm{𝟏},𝐤_\mathrm{𝟐})|^2}{(dN/d^3𝐤_\mathrm{𝟏})(dN/d^3𝐤_\mathrm{𝟐})},$$
(1)
where $`𝐤_\mathrm{𝟏}`$ and $`𝐤_\mathrm{𝟐}`$ are momenta of a emitted pion pair and $`I(𝐤_\mathrm{𝟏},𝐤_\mathrm{𝟐})`$ is the interference term.
As usual, we introduce the relative momentum $`q^\mu =k_1^\mu k_2^\mu `$ and the average momentum $`K^\mu =(k_1^\mu +k_2^\mu )/2`$ which can be put as $`K^\mu =(K^0,K_T,0,K_L)`$ by virtue of the cylindrical symmetry. Due to the on-shell condition, $`q_\mu K^\mu =0`$, only three components of the relative momenta are independent. In the YKP parametrization, $`q_{}=\sqrt{q_x^2+q_y^2}`$, $`q_{}=q_z`$, and $`q^0`$ are used as the independent components. Then, the fitting function is given by
$$\begin{array}{c}C(q^\mu ,K^\mu )=1+\lambda \mathrm{exp}\{R_{}^2(K^\mu )q_{}^2R_{}^2(K^\mu )[q_{}^2(q^0)^2]\hfill \\ \hfill [R_0^2(K^\mu )+R_{}^2(K^\mu )]\left[q_\mu u^\mu (K^\mu )\right]^2\},\end{array}$$
(2)
where $`u^\mu `$ is a four-velocity which has only longitudinal component. The size parameters to be determined by a fit are $`R_i`$s $`(i=,,0)`$. We may choose a special reference frame where the velocity parameter $`v`$ vanishes (YKP frame), because the three size parameters are invariant under longitudinal boosts. Introducing the source function,
$$S(x^\mu ,K^\mu )=_\mathrm{\Sigma }\frac{U_\mu (x^{})d\sigma ^\mu (x^{})}{(2\pi )^3}\frac{U_\nu (x^{})k^\nu }{\mathrm{exp}(U_\rho (x^{})k^\rho /T)1}\delta ^4(xx^{}),$$
(3)
and the weighted average,
$$A(x^\mu )=\frac{d^4xA(x^\mu )S(x^\mu ,K^\mu )}{d^4xS(x^\mu ,K^\mu )},$$
(4)
the size parameters in this frame (YKP frame) are expressed as
$`R_{}^2(K^\mu )`$ $`=`$ $`(\mathrm{\Delta }y)^2,`$ (5)
$`R_0^2(K^\mu )`$ $`=`$ $`(\mathrm{\Delta }t)^2{\displaystyle \frac{2}{\beta _{}}}\stackrel{~}{x}\stackrel{~}{t}+{\displaystyle \frac{1}{\beta _{}^2}}[(\mathrm{\Delta }x)^2(\mathrm{\Delta }y)^2],`$ (6)
$`R_{}^2(K^\mu )`$ $`=`$ $`(\mathrm{\Delta }z)^2{\displaystyle \frac{2\beta _{}}{\beta _{}}}\stackrel{~}{x}\stackrel{~}{z}+{\displaystyle \frac{\beta _{}^2}{\beta _{}^2}}[(\mathrm{\Delta }x)^2(\mathrm{\Delta }y)^2],`$ (7)
where $`\stackrel{~}{x}=xx`$, $`\mathrm{\Delta }x=\sqrt{x^2x^2}`$, $`\beta _{}=K_T/K^0`$ and $`\beta _{}=K_L/K^0`$.
In order to regard the above $`R_i`$s as naive source sizes directly, the second term and the third term in the right hand side of Eqs.(6) and (7) must be small compared to the first term of each equation. Though it has been shown that this condition holds within a class of thermal model , it is not trivial in other cases. The case of a hydrodynamic model as a more realistic one is the point of the present paper. In the following, we calculate three quantities: HBT radii extracted from the fit to Eq.(1), the space-time extensions, Eqs.(5)-(7), and the source sizes, the first terms in Eqs.(6) and (7).
## 3 DISCUSSION
Fig. 1 $`K_T`$ dependence of YKP radii. Open and closed circles show the experimental data and HBT radii, respectively. Solid lines stand for space-time extensions (5)-(7). Dotted lines stand for source sizes. ($`\mathrm{\Delta }z`$ for $`R_{}`$ and $`\mathrm{\Delta }t`$ for $`R_0`$.)
Figure 1 shows $`K_T`$ dependence of the pion HBT radii. Close circles stand for the HBT radii. Solid lines are space-time extensions which are expected to agree with the HBT radii when only thermal pions are considered. Dotted lines stand for the source sizes ($`\mathrm{\Delta }z`$ for $`R_{}`$ and $`\mathrm{\Delta }t`$ for $`R_0`$). As far as $`R_{}`$ and $`R_{}`$ are concerned, those quantities seem to be consistent with one another and well reproduce the NA49 experimental results (open circles). On the other hand, $`R_0`$ is smaller than the experimental one in spite of the fact our calculation includes a first order phase transition. The time durations $`\mathrm{\Delta }t`$ agrees with $`R_0`$ at small $`K_T`$. However, some deviations from $`R_0`$ can be seen at large $`K_T`$ region, that means a failure of interpretation of $`R_0`$ as the time duration. We find that this difference occurs from the third term in Eq.(6).
In Eq.(6), the two additional terms means correlation between $`x`$ and $`z`$ (the second term) and the thickness of the source (the third term). Though contribution from the second term is small in our model, the third term becomes negative in same order with $`(\mathrm{\Delta }t)^2`$ . Such a kind of the source is called an “opaque source”. The “opaque” source emits particles dominantly from the surface and the “transparent” source emits from whole region. Figure 2 shows the source function (3) projected onto the $`xy`$ plane, $`\stackrel{~}{S}_T(x,y)=dzdtS(x^\mu ,K^\mu )`$ being normalized as $`𝑑x𝑑y\stackrel{~}{S}_T(x,y)=1`$. Note that the average transverse momentum of the emitted pions is parallel to $`x`$ direction in the figure. The left figure shows clearly surface emission; pions are emitted from the crescent region and the source is thinner in the $`x`$ direction than in the $`y`$ direction. This is the typical property of the opaque source which appears in the factor $`(\mathrm{\Delta }x)^2(\mathrm{\Delta }y)^2`$ in Eq.(6). Note that this factor also appears in the longitudinal components of the size, Eq.(7), but the prefactor $`\beta _{}^2`$ makes the contribution to $`R_{}^2`$ small. If we neglected the transverse flow by artificially putting $`U^r=0`$, the source function becomes as shown in the right figure in Fig. 2. In this case, the source function is proportional to the space-time volume of freeze-out hypersurface. As a consequence of restoration of the azimuthal symmetry which is shown clearly in the figure, the above measure of source opacity vanishes. Although the particle emission takes place almost from the thin surface, the source is not opaque in this sense. When the transverse flow exists, the source function is deformed by thermal Boltzmann factor $`\mathrm{exp}(K_T\mathrm{sinh}Y_T\mathrm{cos}\varphi /T_f)`$ with $`Y_T`$ almost proportional to $`r`$. Consequently, the number of emitted particles increases in the region where $`x>0`$ (i.e., $`\mathrm{cos}\varphi >0`$) and decreases in the region where $`x<0`$ ($`\mathrm{cos}\varphi <0`$). This flow effect deforms the surface dominant volume distribution (right figure) to the crescent shape (left figure).
Fig. 2 The source functions as functions of transverse coordinates $`x`$ and $`y`$ for $`K_T=450`$ MeV and $`Y_{\pi \pi }=4.15`$ ($`Y_{\pi \pi }=\frac{1}{2}\mathrm{ln}\frac{K^0+K_L}{K^0K_L}`$). The left figure denotes the source function with the transverse flow. The right figure denotes the source function without transverse flow.
In summary, we analyze $`K_T`$ dependence of source parameters of the YKP parametrization based on the relativistic hydrodynamics for the CERN-SPS 158 GeV/A Pb+Pb collisions. We obtain the results which are almost consistent with the experiment. However, the temporal source parameter, $`R_0`$, shows smaller value than the experiment. The source opacity makes the interpretation of $`R_0`$ as the time duration doubtful at large $`K_T`$ region. We find that the source opacity is caused by the transverse flow and the characteristics of the freeze-out hypersurface, the surface dominant freeze-out. The deviations of our results from the experiment would be improved by including the resonance decay and other effects.
The authors are indebted to Professor I. Ohba and Professor H. Nakazato for their helpful comments. This work was partially supported by a Grant-in-Aid for Science Research, Ministry of Education, Science and Culture, Japan (Grant No. 09740221) and Waseda University Media Network Center.
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# Magnetic, thermal, and transport properties on single crystals of antiferromagnetic Kondo-lattice Ce2PdSi3
## I Introduction
There has been considerable interest in understanding the interplay among the crystalline-electric-field (CEF) effect, the indirect exchange \[Ruderman-Kittel-Kasuya-Yoshida (RKKY)\] interaction among the 4$`f`$ magnetic moments and the Kondo effect in Ce compounds, since these are the decisive factors of the physical properties in these compounds. It is therefore worthwhile to carry out careful investigation in new Ce compounds. With this motivation, we report here the results of magnetic susceptibility ($`\chi `$), magnetization ($`M`$), heat-capacity ($`C`$), electrical-resistivity ($`\rho `$), thermoelectric-power ($`S`$), and Hall-coefficient ($`R_H`$) measurements on single-crystalline Ce<sub>2</sub>PdSi<sub>3</sub>, grown for the first time.
This compound has been reported to form in an AlB<sub>2</sub>-derived hexagonal crystal structure and to exhibit Kondo effect. The intensity of investigation in the RE$`{}_{2}{}^{}X`$Si<sub>3</sub> (RE $`=`$ rare earth, $`X`$ $`=`$ transition metal) series, crystallizing in the above-mentioned structure, increased only in the recent years and these compounds have been reported to exhibit many unusual features in the magnetic, thermal, and transport properties (see, for instance, Refs. 1–10 and references therein). Gd<sub>2</sub>PdSi<sub>3</sub> exhibits Kondo-lattice-like anomalies, e.g., a resistivity minimum above $`T_N`$ accompanied by a large negative magnetoresistance. These features, presumably due to a novel mechanism, are not common to Gd compounds. Ce<sub>2</sub>CoSi<sub>3</sub> is a mixed-valent compound, a small La substitution for Ce induces a non-Fermi-liquid behavior in $`\rho `$. Eu<sub>2</sub>PdSi<sub>3</sub> exhibits two distinct magnetic transitions, with the possibility of quasi-one-dimensional magnetism for the high-temperature transition and unusual magnetic characteristics. Particularly considering that the Gd-based compound in the Pd series has been found to show many interesting anisotropic features, it is tempting to carry out detailed studies on the single-crystalline Ce<sub>2</sub>PdSi<sub>3</sub> as well. Previous magnetic, electrical-resistance, and heat-capacity measurements on this compound were performed only in the polycrystalline form and no clear magnetic ordering could be detected. Thus the present studies extended to much lower temperatures, particularly on single crystals, serve as a first thorough characterization of the bulk properties of this compound.
## II Experimental details
Single crystals of Ce<sub>2</sub>PdSi<sub>3</sub> have been prepared by the Czochralski pulling method using a tetra-arc furnace in an argon atmosphere. The single-crystalline nature has been confirmed by back-reflection Laue technique. The magnetic measurements were carried out with a Quantum Design superconducting quantum interference device (SQUID) magnetometer. The heat capacity was measured by a quasiadiabatic heat- pulse method using a dilution refrigerator. The electrical-resistivity and Hall-effect measurements have been performed by a conventional dc four-probe method in the temperature interval of 0.5–300 K. The thermoelectric-power data have been taken by the differential method using a Au-Fe (0.07%)-chromel thermocouple.
## III Results and discussions
Figure 1(a) shows the temperature dependence of the inverse magnetic susceptibility $`\chi `$$`{}_{}{}^{1}(T)`$, measured in a magnetic field $`H=`$ 1 kOe for both $`H//[10\overline{1}0]`$ and $`H//[0001]`$. There is a large difference in the absolute values of $`\chi `$ for two geometries, apparently due to CEF effect. The effective magnetic moment ($`\mu _{eff}`$) and the paramagnetic Curie temperature $`\mathrm{\Theta }_P`$ are estimated from the high-temperature linear region to be about 2.60$`\mu _B`$/Ce
and 3.6 K for $`H//[10\overline{1}0]`$ and 2.65$`\mu _B`$/Ce and $``$107 K for $`H//[0001]`$, respectively. These values of $`\mu _{eff}`$ are very close to that expected for a free trivalent Ce ion (2.54$`\mu _B`$). The large negative value of $`\mathrm{\Theta }_P`$ for $`H//[0001]`$ and a small positive value for $`H//[10\overline{1}0]`$ are presumably due to CEF effect. The large anisotropy in $`\mathrm{\Theta }_P`$ may also indicate the existence of anisotropy in the exchange interaction which depends on the CEF level scheme. The expanded view of the temperature dependence of $`\chi `$ is shown in the inset of Fig. 1(a). There is no difference between the field-cooled and the zero-field-cooled measurements of $`\chi `$ down to 1.9 K within the experimental accuracy, indicating the absence of any spin-glass-like behavior. This fact is in contrast to the formation of spin-glass state in the isostructural U<sub>2</sub>PdSi<sub>3</sub>. There is a peak in $`\chi (T)`$ at around 2.8 K for $`H//[10\overline{1}0]`$ and at 2.5 K for $`H//[0001]`$. These peaks can be ascribed to the antiferromagnetic ordering. The occurrence of the peak at slightly different temperatures for two directions might be due to the anisotropic field dependence of $`T_N`$ for two directions, since the heat-capacity measurement in absence of a magnetic field shows a peak at around 3 K as shown below. Recent neutron-diffraction data on polycrystals also suggests the occurrence of antiferromagnetic (AF) ordering below 2.5 K in a sinusoidally modulated AF structure.
We have tried to analyse the $`\chi (T)`$ data using a CEF model, considering hexagonal site symmetry of Ce in Ce<sub>2</sub>PdSi<sub>3</sub>. According to the Hutchings’ notation, the CEF Hamiltonian for $`J=5/2`$ ion with the hexagonal point symmetry is given by
$$=B_2^0O_2^0+B_4^0O_4^0,$$
(1)
where $`B_n^m`$ and $`O_n^m`$ represent the CEF parameters and the Steven’s equivalent operators, respectively. The results of this CEF analysis using Eq. (1), employing the $`\chi (T)`$ data at the paramagnetic region, leads to $`B_2^0`$ 10.1 K, and $`B_4^0`$ 0.11 K. This set of parameters corresponds to a crystal-field level scheme with the three doublets $`\pm \frac{1}{2}`$, $`\pm \frac{3}{2}`$ and $`\pm \frac{5}{2}`$ at around 0, 28 K ($`\mathrm{\Delta }_1`$), and 130 K ($`\mathrm{\Delta }_2`$), respectively. Accordingly, the calculated values of $`\chi `$<sup>-1</sup> are shown by the solid lines in Fig. 1(a), which indicates that the anisotropy in $`\chi (T)`$ is mainly induced by the CEF effect. However, there is a deviation of the calculated $`\chi `$<sup>-1</sup> from the experimental values, particularly at low temperatures. The following explanations can be offered to this deviation: According to Szytula et al, Pd and Si atoms are random in this crystal structure (space group $`P6/mmm`$). Therefore this randomness or disorder between Pd and Si sites may produce local modification of the CEF effects due to the distribution of the CEF parameters. Alternatively, if Pd and Si are well ordered (space group $`P6_3/mmc`$, see Ref. 10), there are two different crystallographic environments for Ce ions, in which case the CEF effect may be different for these two sites. All these factors are neglected in the present CEF calculations.
The isothermal magnetization at 2 K is shown in Fig. 1(b). $`M`$ varies distinctly in different ways with the applied magnetic field for $`H//[10\overline{1}0]`$ and $`H//[0001]`$. This anisotropy in $`M`$ is presumably due to the CEF effect. The magnetic moments at $`H=`$ 50 kOe for two directions are $``$ 1.18 and $``$ 0.25$`\mu _B`$/Ce for $`H//[10\overline{1}0]`$ and $`H//[0001]`$, respectively. The larger magnetic moment for $`H//[10\overline{1}0]`$ indicates the $`a`$-$`b`$ plane as the easy plane of magnetization, and the $`\pm \frac{1}{2}`$ doublet as the ground state, in agreement with the CEF analysis from $`\chi (T)`$ described above. These two facts are also consistent with the proposed magnetic structure of Ce<sub>2</sub>PdSi<sub>3</sub> based on neutron-diffraction experiment, i.e., the Ce magnetic moments lie in the $`a`$-$`b`$ plane.
Figure 2 shows the temperature dependence of heat capacity, $`C(T)`$, for Ce<sub>2</sub>PdSi<sub>3</sub>. A peak in $`C(T)`$ at around 3 K indicates the existence of the antiferromagnetic ordering below $`T_N`$ 3 K, supporting the conclusion from the $`\chi `$$`(T)`$ behavior. The transition is, however, not very sharp, which might be due to the presence of site disorders. The inset shows the $`C/T`$ vs $`T^2`$ plot at temperatures far below $`T_N`$ where phonon contribution is negligible. The linear coefficient of specific heat ($`\gamma `$), named as the Sommerfeld coefficient, estimated from this plot at such low temperatures is 108 mJ/K<sup>2</sup> mol Ce. At this temperature range the specific heat can be expressed as $`C=\gamma T+\beta T^3`$ with $`\beta `$ 862 mJ/K<sup>4</sup> mol Ce and the parameter $`\beta `$ mainly comes from the contribution of an antiferromagnetic-magnon part. For Y<sub>2</sub>PdSi<sub>3</sub>, the non-magnetic compound serving as a reference for phonon contribution, the $`\gamma `$ value is 4.5 mJ/K<sup>2</sup> molY. The moderately large value of $`\gamma `$ for Ce<sub>2</sub>PdSi<sub>3</sub> suggests that even in the magnetically-ordered state Ce<sub>2</sub>PdSi<sub>3</sub> may be classified as a heavy fermion. The solid line in Fig. 2 shows the magnetic entropy ($`S_{4f}`$) estimated from the 4$`f`$ contribution ($`C_m`$) to $`C`$. $`C_m`$ is obtained by employing the $`C`$ values of Y<sub>2</sub>PdSi<sub>3</sub> (Ref. 1) as a reference for the lattice contribution taking into account the difference of Debye temperatures of the two compounds by using the procedure suggested in Ref. 14. At $`T_N`$, $`S_{4f}`$ 2.3 J/K mol Ce is $``$ 40% of $`R`$ ln 2 expected for a complete removal of the two fold degeneracy of a CEF ground-state doublet. This reduced entropy value might be due to the substantial Kondo-derived reduction of the Ce moments and/or a presence of short-range correlations above $`T_N`$. Tentatively assuming a Kondo-derived reduction, the magnetic entropy $`S_{4f}`$ ($``$ 0.4 $`R`$ ln 2) at $`T_N`$ yields the Kondo temperature $`T_K`$ 8 K, according to the Bethe-ansatz for a spin-$`\frac{1}{2}`$ Kondo model (see Refs. 16 and 17). The theoretical calculations of $`C(T)`$ using $`T_K`$ 8 K, however, cannot reproduce the experimental curve above $`T_N`$. This deviation indicates the presence of short-range-AF correlations above $`T_N`$; the presence of these correlations is, in fact, detected in a recent neutron-scattering experiment, and a possible contribution from an excited CEF level (28 K) even at the measured temperature range.
The temperature dependence of resistivity $`\rho (T)`$ for Ce<sub>2</sub>PdSi<sub>3</sub> with the current $`J//[10\overline{1}0]`$ and $`J//[0001]`$ as well as in polycrystalline Y<sub>2</sub>PdSi<sub>3</sub> is shown in Fig. 3. $`\rho `$ for both current directions gradually decreases with decreasing temperature down to about 20 K showing a broad hump around 100 K; below 20 K, there is a weak upturn giving rise to a minimum at around 20 K, followed by a drop below 8 K. The $`\rho (T)`$ in Y<sub>2</sub>PdSi<sub>3</sub> shows usual metallic behavior, however there is a drop in $`\rho `$ below 6 K, presumably attributable to the presence of traces of the superconducting phase YPdSi. As known for many other Ce alloys, the broad hump in Ce<sub>2</sub>PdSi<sub>3</sub> can be ascribed to the combined effect of CEF and Kondo effect. The magnetic contribution to the resistivity $`\rho _m`$$`=\rho `$(Ce<sub>2</sub>PdSi<sub>3</sub>)$`\rho `$(Y<sub>2</sub>PdSi<sub>3</sub>), obtained by using $`\rho (T)`$ data for $`J//[0001]`$ in Ce<sub>2</sub>PdSi<sub>3</sub>, is shown in the inset; the ratio of the slopes of the $`\rho _m`$ vs ln $`T`$ plot at high to low temperature turns out close to 35/3 confirming that the ground state is a doublet (see Ref. 19). It clearly reveals the presence of a peak at around 50 K which might be related with $`T_K`$ enhanced by CEF effect as proposed by Hanzawa et al. Correspondingly, there is also a broad hump in $`S(T)`$ (see Fig. 4). The minimum in $`\rho (T)`$ around 20 K is either due to the Kondo effect or a consequence of magnetic-precursor effect as noted for some Gd alloys. It may be noted that the drop in $`\rho `$ sets in at 8 K, much above the $`T_N`$, similar to CePt<sub>2</sub>Ge<sub>2</sub>,
and such a feature in magnetically ordered Kondo lattices arises from a combination of indirect exchange interaction and the Kondo effect. There is a small difference in absolute values for two geometries, which might be due to the combined effect of anisotropy in the Fermi surface and preferably oriented microcracks. It may also be remarked that the residual resistivity is large even for the single crystal, which might be due to a presence of crystallographic (Pd-Si) disorder or a combined effect of disorder and a dominance of Kondo contribution even in the magnetically ordered state.
Figure 4 shows the temperature dependence of thermoelectric power ($`S`$) in Ce<sub>2</sub>PdSi<sub>3</sub> as well as in Y<sub>2</sub>PdSi<sub>3</sub> (polycrystal). In Ce<sub>2</sub>PdSi<sub>3</sub>, for the temperature gradient $`\mathrm{\Delta }`$$`T//[10\overline{1}0]`$, $`S`$ is positive at room temperature, then it gradually increases with decreasing temperature and shows a broad hump around 100 K. There is a change of sign around 50 K with a minimum around 17 K. On the other hand, for $`\mathrm{\Delta }`$$`T//[0001]`$, $`S`$ has a large negative value at room temperature and decreases with decreasing
temperature showing a broad hump around 100 K. $`S(T)`$ for this direction also shows a minimum, however, at a temperature slightly higher than that for $`\mathrm{\Delta }`$$`T//[10\overline{1}0]`$. Thus $`S(T)`$ in Ce<sub>2</sub>PdSi<sub>3</sub> is highly anisotropic with the directions of the thermal gradient. The anisotropy in the Fermi surface might be one reason behind this anisotropy, since $`S(T)`$ in the isostructural Gd<sub>2</sub>PdSi<sub>3</sub> is also anisotropic. For nonmagnetic Y<sub>2</sub>PdSi<sub>3</sub>, $`S(T)`$ has a large negative value at room temperature and decreases gradually with temperature. In Ce<sub>2</sub>PdSi<sub>3</sub>, the broad hump can be attributed to the interplay between CEF and Kondo effect. The origin of the minimum at $`T_{min}`$ $``$ 17 K is not clear yet, however, the possible explanation may be the Kondo scattering in the CEF ground state or the growth of AF correlations as in the case of CeAuAl<sub>3</sub>. Tentatively assuming the Kondo-derived origin, $`T_K`$ $``$ 8 K would be obtained using the relation $`T_K`$ $``$ $`\frac{1}{2}T_{min}`$ that holds for CeAl<sub>2</sub> and CeCu<sub>2</sub> (see Ref. 15). However, the temperature dependence of $`C`$ above $`T_N`$ suggests that this estimation of $`T_K`$ is a rough one, indicating that both Kondo effect and AF correlations may play a role for this minimum. For $`\mathrm{\Delta }`$$`T//[10\overline{1}0]`$, $`S(T)`$ is similar to the behavior in the typical magnetically ordered heavy-Kondo compounds, e.g., CeCu<sub>2</sub> and CeAl<sub>2</sub>, though the overall temperature dependence of $`S`$ in Ce<sub>2</sub>PdSi<sub>3</sub> is rather weaker. If the crystallographic disorder is the dominant origin of the large residual resistivity, the Kondo contribution to $`S(T)`$ could be suppressed. Since, according to the Gorter-Nordheim rule, the thermoelectric power for more than one scattering mechanisms can be expressed as $`S_{alloy}=[\rho _1S_1+\rho _2S_2]/\rho `$, where the subscripts 1 and 2 correspond to different scattering mechanism: In the present case, 1 represents the Kondo scattering and 2 represents the other scatterings. Therefore a large $`\rho _2`$ can suppress the Kondo contribution $`S_1`$. The large negative thermoelectric power in Y<sub>2</sub>PdSi<sub>3</sub> might arise from the 4$`d`$ band of Pd, as in the case of the 3$`d`$ band of Co in YCo<sub>2</sub>.
The $`S(T)`$ curve for $`\mathrm{\Delta }T`$// at high temperatures is almost parallel to that of Y<sub>2</sub>PdSi<sub>3</sub>, therefore a significant effect of the Pd-4$`d`$ band on $`S(T)`$ even in Ce<sub>2</sub>PdSi<sub>3</sub> cannot be ruled out. The temperature dependence of Hall coefficient ($`R_H`$) for $`H=`$ 15 kOe, shown in Fig. 5(a), also reflects the anisotropic nature of this material. For both geometries (as labeled in the figure) $`R_H`$ is positive at room temperature and increases gradually with decreasing temperature. At low temperatures $`R_H`$ becomes highly anisotropic and shows a positive peak for both geometries, in the vicinity of $`T_N`$ (see the inset). The large anisotropy observed in $`R_H`$ is also reflected in the $`\chi (T)`$ data taken at $`H=`$ 15 kOe (not shown), indicating that the anisotropy in $`R_H`$ is of magnetic origin. The positive value of $`R_H`$ at all temperatures and a positive peak at the vicinity of $`T_N`$ are similar to those in the antiferromagnetic Kondo-lattice compound CeAl<sub>2</sub>. In contrast, $`R_H`$ ($``$ 0.9$`\times `$10<sup>-10</sup> m<sup>3</sup>/C at 300 K) in Y<sub>2</sub>PdSi<sub>3</sub> is almost temperature independent. Clearly, there is a dominant 4$`f`$ contribution in Ce<sub>2</sub>PdSi<sub>3</sub>. The Hall coefficient in magnetic materials like those in Ce compounds is generally a sum of two terms; an ordinary Hall coefficient ($`R_0`$) due to Lorentz force and an anomalous part arising from magnetic scattering (skew scattering); in the paramagnetic state $`R_H=R_0+A\rho \chi `$, where $`A`$ is a constant. Using this relation, $`R_0`$ is estimated by plotting $`R_H`$ versus $`\rho \chi `$ \[Fig. 5(b)\]. From Fig. 5(b), it is obvious that the plot is linear for both $`H//[10\overline{1}0]`$ and $`H//[0001]`$ in the paramagnetic state with a value of $`R_0`$3.2$`\times `$ 10<sup>-10</sup> m<sup>3</sup>/C and 1.0$`\times `$10<sup>-9</sup> m<sup>3</sup>/C, and $`A`$ 7.4$`\times `$10<sup>-16</sup> mol/C and 4.8$`\times `$10<sup>-16</sup> mol/C, respectively. This linear behavior indicates the presence of dominant skew scattering in Ce<sub>2</sub>PdSi<sub>3</sub>. In the vicinity of $`T_N`$, however, the data deviate from the high-temperature linear variation. $`R_0`$ of different sign with the anisotropic values for two geometries indicates the presence of anisotropy in the Fermi surface, in agreement with the $`S(T)`$ data.
## IV Summary
Summarizing, we have investigated the magnetic behavior of recently synthesized Ce<sub>2</sub>PdSi<sub>3</sub> in the single-crystalline form and the results show strong anisotropic behavior of the measured properties. The paramagnetic Curie temperature for $`H//[10\overline{1}0]`$ is positive, however, the value is negative for $`H//[0001]`$. The sign of the thermoelectric power is different for the two measured crystallographic orientations at high temperatures. Distinct features due to an interplay between CEF and Kondo effect have also been observed in the thermoelectric power and resistivity data. The ordinary Hall-coefficient is anisotropic with opposite sign for the two measured geometries. The results establish that this compound is an antiferromagnetic Kondo lattice with $`T_N`$= 3 K. The magnitude of $`T_K`$ is also estimated to be of same order as $`T_N`$ and this fact suggests a delicate competition between the Kondo effect and indirect exchange interaction. Therefore it would be of interest to investigate this compound under high pressure.
ACKNOWLEDGMENT
This work has been partially supported by a Grant-in-Aid for Scientific Research from the Minstry of Education, Science, Sports and Culture of Japan.
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# 1 Introduction
## 1 Introduction
One of the most fascinating challenges in particle physics is the unification of the forces of nature. Grand unified theories (GUTs) have been advocated as an extension of the standard model (SM) to unify the gauge interactions. Their immediate consequence is the explanation of charge quantization. The scale at which GUTs must replace the SM, however, must be very high, $`M_{\mathrm{GUT}}>10^{1516}`$ GeV, in order to avoid higher-dimensional operators that would lead to a large proton decay rate. Furthermore, a large $`M_{\mathrm{GUT}}`$ is also needed for gauge-coupling unification. Since the 4D gauge couplings evolve logarithmically with the energy scale, a large $`M_{\mathrm{GUT}}`$ is needed to allow the different gauge couplings of the SM to get closer and, eventually, unify. This is the case in the supersymmetric extension of the standard model (MSSM) where gauge couplings actually unify at $`M_{\mathrm{GUT}}10^{16}`$ GeV. Since supersymmetry is also needed for the stability of the weak scale versus the large GUT scale, supersymmetric GUTs provide a very appealing framework for physics beyond the SM. Nevertheless, we must face the fact that, since GUT fields will only appear at very high energies $`M_{\mathrm{GUT}}10^{16}`$ GeV, GUTs will never be tested in a direct way. These theories predict a big “desert” from the weak to the GUT scale.
Here we will present a GUT scenario without the “desert”. It has GUT fields of masses of the order of a TeV. This scenario has been proposed in Refs. and is based on an extension of the MSSM to a five-dimensional theory with the extra dimension compactified as in the Randall–Sundrum model . Differently from Ref. , however, we will consider the MSSM gauge sector and its GUT extension living in 5D with matter localized on a 4D boundary. In this letter we will show that, even though the theory is five-dimensional, gauge couplings get logarithmic corrections at the one-loop level. Therefore the theory predicts, as in the 4D MSSM case, the right values of the gauge couplings at low energies. We will also show that proton decay operators are suppressed by the high scale $`M_{\mathrm{GUT}}10^{16}`$ GeV.
Different attempts to obtain theories that, while predicting gauge-coupling unification as in the MSSM, do not have a “desert” between the weak and the GUT scale, can be found in Refs. .
## 2 The set-up
Our set-up is based on the Randall–Sundrum 5D model , where the bulk is a slice of AdS<sub>5</sub>. This corresponds to a 5D non-factorizable geometry with a the fifth dimension $`y`$ compactified on an orbifold, $`S^1/_2`$, of radius $`R`$ with $`0y\pi R`$. The orbifold fixed points at $`y^{}=0`$ and $`y^{}=\pi R`$ are 4D boundaries of the five-dimensional space-time. The metric is given by
$$ds^2=e^{2ky}\eta _{\mu \nu }dx^\mu dx^\nu +dy^2,$$
(1)
where $`1/k`$ is the AdS curvature radius and $`\eta _{\mu \nu }=\mathrm{diag}(1,1,1,1)`$ with $`\mu =1,\mathrm{},4`$. The fundamental scale in the 5D theory, $`M_5`$, is related with the 4D Planck mass, $`M_P`$, by $`M_5^3kM_P^2`$ (for $`R>1/k`$). We assume that all the scales are of roughly the same order of magnitude $`kM_5M_P`$, with the radius of the extra dimension slightly larger, $`R11/k`$. The effective scales on the two boundaries are very different. On the $`y^{}=0`$ boundary the effective scale is $`M_P`$, while on the $`y^{}=\pi R`$ boundary this is given by $`ke^{k\pi R}`$ TeV (for the assumed value $`R11/k`$). We will hence call these two boundaries the $`M_P`$-boundary and the TeV-boundary respectively.
## 3 One-loop corrections to the gauge propagator <br>at low energies
Let us consider a 5D gauge boson , $`A_M(x,y)`$, with $`M=(\mu ,5)`$, living in the slice of AdS<sub>5</sub> described above. We want to calculate the one-loop corrections to the gauge propagator at low energies. For simplicity, we will consider a 5D scalar QED theory. We will work in the gauge $`A_5(x,y)=0`$, so we only have to consider $`A_\mu (x,y)`$. At energies below the Kaluza–Klein (KK) masses, only the massless zero-mode of the photon is relevant. This is given by
$$A_\mu (x,y)=\frac{1}{\sqrt{\pi R}}A_\mu ^{(0)}(x)+\mathrm{}.$$
(2)
This corresponds to a 4D massless state with a $`y`$-independent wave-function. Contrary to the graviton case, this gauge mode is not localized by the AdS metric and has the interesting property that it couples to the two boundaries of the orbifold with equal strength.
We want to calculate the one-loop corrections to the propagator of this massless photon generated by a 5D scalar $`\varphi `$ with charge $`1`$ and even under the $`_2`$. We will regularize this theory with a 5D Pauli–Villars (PV) field $`\mathrm{\Phi }`$ of mass $`\mathrm{\Lambda }`$. This mass corresponds to the cut-off of the theory that we will take to be $`\mathrm{\Lambda }<k`$. Let us decompose the 5D scalar fields $`\varphi `$ and $`\mathrm{\Phi }`$ in KK modes. This has been done in Refs. . For a 5D scalar particle of mass $`M`$, the approximate KK mass spectrum for $`M<k`$ is given in Table 1.
We have defined the $`n=0`$ mode as the mode that becomes massless in the limit $`M0`$. For $`M`$ of order $`k`$, this mode becomes heavier than some of the KK states since $`m_{\mathrm{KK}}=\pi ke^{k\pi R}k`$ for $`R>1/k`$. This is very different from compactifications in a slice of flat space where the $`n=0`$ mode (defined as explained above) is always the lightest state. This is the effect of the AdS<sub>5</sub> curvature that lower the masses of the KK spectrum but not the mass of the zero mode. For the scalar $`\varphi `$ and the PV field $`\mathrm{\Phi }`$, we can obtain the KK spectrum using Table 1 with the following values for $`\alpha `$:
$$\alpha ^\varphi =2,\alpha ^\mathrm{\Phi }=\sqrt{4+\frac{\mathrm{\Lambda }^2}{k^2}}2+\frac{\mathrm{\Lambda }^2}{4k^2}.$$
(3)
From this KK decomposition we can already infer the magnitude of the quantum corrections. For each KK mode of the field $`\varphi `$ there is a KK mode of the PV field whose mass acts as a cut-off scale. Since the masses of the KK modes of $`\varphi `$ and $`\mathrm{\Phi }`$ are of the same order of magnitude, we do not expect large corrections from them. Nevertheless, the zero mode of the PV field is very heavy, $`𝒪(\mathrm{\Lambda })`$, in contrast with the zero mode of $`\varphi `$ that is massless. Therefore we expect a large correction coming from this large mass splitting of the zero modes that will reproduce the quantum corrections of an ordinary 4D theory.
To see this explicitly, let us now calculate the one-loop contribution from the scalar $`\varphi `$ to the propagator of the photon $`A_\mu ^{(0)}`$. Defining the photon self-energy by $`\mathrm{\Pi }_{\mu \nu }(q)=[q^2\eta _{\mu \nu }q_\mu q_\nu ]\mathrm{\Pi }(q)`$, we find, at zero momentum, that the one-loop contribution is given by
$$\mathrm{\Pi }(0)=\frac{b_0}{8\pi ^2}\mathrm{ln}\frac{M_0^\varphi }{M_0^\mathrm{\Phi }}+\frac{b_{\mathrm{KK}}}{8\pi ^2}\underset{n=1}{\overset{e^{k\pi R}}{}}\mathrm{ln}\frac{M_n^\varphi }{M_n^\mathrm{\Phi }},$$
(4)
where $`M_0^\varphi `$ and $`M_n^\varphi `$ are respectively the masses of the zero mode and $`n`$-KK mode of the field $`\varphi `$, and similarly for $`\mathrm{\Phi }`$. We denote by $`b_0`$ and $`b_{\mathrm{KK}}`$ the beta-function coefficients of the zero mode and KK modes respectively. In the example here we have $`b_0=b_{\mathrm{KK}}=1/3`$. Using Table 1 with Eq. (3), we obtain
$$\mathrm{\Pi }(0)\frac{b_0}{8\pi ^2}\mathrm{ln}\frac{\mu }{\mathrm{\Lambda }}+\frac{b_{\mathrm{KK}}}{8\pi ^2}_1^{e^{k\pi R}}𝑑n\mathrm{ln}\left(\frac{n+1/4}{n+1/4+\mathrm{\Lambda }^2/(8k^2)}\right),$$
(5)
where as usual we have introduced an infra-red cut-off $`\mu `$, and we have replaced the sum over $`n`$ by an integral. Evaluating the integral and considering that $`\mathrm{\Lambda }<k`$, we obtain
$$\mathrm{\Pi }(0)\frac{b_0}{8\pi ^2}\mathrm{ln}\frac{\mu }{\mathrm{\Lambda }}+\frac{b_{\mathrm{KK}}}{64\pi ^2}\frac{\mathrm{\Lambda }^2}{k^2}\mathrm{ln}\frac{5m_{\mathrm{KK}}}{4\pi k}.$$
(6)
For $`\mu \mathrm{\Lambda }<k`$ the KK contribution is small and can be neglected. We then obtain that the contribution to the gauge boson propagator is dominated by the zero mode and gives exactly the same contribution as in 4 dimensions:
$$\mathrm{\Pi }(0)\frac{b_0}{8\pi ^2}\mathrm{ln}\frac{\mu }{\mathrm{\Lambda }}.$$
(7)
This is the main result of this paper. It shows that in these 5D theories the contribution to the massless-mode gauge-boson propagator depends logarithmically on the high-energy cut-off $`\mathrm{\Lambda }`$, provided that $`\mu \mathrm{\Lambda }<k`$.
It is important to see if the result of Eq. (7) can be understood as a running of the gauge coupling similar to the 4D case. Of course, this cannot be the case if matter is localized on the TeV-boundary, since on that boundary our effective scale (cut-off scale) is TeV, and above this energy effects must be considered from the fundamental (string) theory. Nevertheless, if matter is localized on the $`M_P`$-boundary, where the effective scale is $`M_P`$, we will show in the next section that the effective gauge couplings can be considered to run logarithmically with the energy similarly to the 4D case.
## 4 The 5D gauge propagator at high energies
In order to understand what is the behavior of the theory at energies above the TeV, we will derive here the 5D propagator of the gauge boson in the AdS<sub>5</sub> slice. Since we are interested in the propagator at high energies from the point of view of the $`M_P`$-boundary, we will consider the limit $`R\mathrm{}`$. In this limit we do not need to do KK decomposition and we can work directly in 5D.
Let us take the gauge $`A_5(x,y)=0`$ and consider only the transverse part of $`A_\mu (x,y)`$, i.e. we impose $`^\mu A_\mu (x,y)=0`$. It is shown in Ref. that the transverse part decouples from the non-transverse part in the equations of motion. Moreover, only the propagator of the transverse part is relevant to sources localized on the $`M_P`$-boundary since the current there is transverse, $`^\mu J_\mu =0`$. Following similar steps to those in the graviton case , we want to calculate the Green function for the gauge boson defined as
$$A_\mu (x,y)=d^4x^{}dy^{}\sqrt{g}G(x,y,;x^{},y^{})e^{2ky^{}}J_\mu (x^{},y^{}),$$
(8)
with $`^\mu J_\mu =0`$. It can be shown that $`e^{k(y+y^{})}G(x,y;x^{},y^{})`$ is also the Green function of a scalar with mass $`3k^2+2k\delta (y)`$ . Let us change the extra dimensional coordinate to $`z=e^{ky}/k`$. Taking the 4D Fourier transform of the Green function
$$G(x,z;x^{},z^{})=\frac{d^4p}{(2\pi )^4}e^{ip(xx^{})}G_p(z,z^{}),$$
(9)
we have that $`G_p(z,z^{})`$ must satisfy the equation
$$\left[_z^2\frac{1}{z}_zp^2\right]G_p(z,z^{})=zk\delta (zz^{}).$$
(10)
Solving Eq. (10) with the Neumann boundary condition on the $`M_P`$-boundary, we find
$$G(x,z;x^{},z^{})=\frac{i\pi k}{2}zz^{}\frac{d^4p}{(2\pi )^4}e^{ip(xx^{})}\left[\frac{J_0(ip/k)}{H_0^{(1)}(ip/k)}H_1^{(1)}(ipz)H_1^{(1)}(ipz^{})J_1(ipz_<)H_1^{(1)}(ipz_>)\right],$$
(11)
where $`H_\nu ^{(1)}=J_\nu +iY_\nu `$ is the Hankel function of order $`\nu `$, $`J_\nu `$ and $`Y_\nu `$ are Bessel functions, and we have defined $`z_>`$ ($`z_<`$) as the greater (lesser) of $`z`$ and $`z^{}`$. In the case where the coordinate $`z^{}`$ is on the $`M_P`$-boundary, $`z^{}=1/k`$, the Green function simplifies to
$$G(x,z;x^{},\frac{1}{k})=ikz\frac{d^4p}{(2\pi )^4}e^{ip(xx^{})}\frac{1}{p}\frac{H_1^{(1)}(ipz)}{H_0^{(1)}(ip/k)}.$$
(12)
We are interested in the limit $`r=|xx^{}|1/k`$ where the Green function is dominated by the small-momentum values of the Fourier transform ($`pk`$) and the Hankel function $`H_0^{(1)}(ip/k)`$ is approximately $`H_0^{(1)}(ip/k)=i\frac{2}{\pi }[\mathrm{ln}\frac{p}{2k}+\gamma ]+𝒪(\frac{p^2}{k^2})`$.
Let us now study the Green function of Eq. (12) in two different limits. First let us consider the propagator at large $`z`$, $`zr`$. In this case we find a falloff of the Green function
$$G(x,zr;x^{},\frac{1}{k})\frac{k}{z^2\mathrm{ln}(kz)},$$
(13)
that is similar to the graviton case . Eq. (13) implies that at large momentum $`p>1/z=ke^{ky}`$ the $`M_P`$-boundary decouples from the TeV-boundary. This is, in fact, a check of consistency of the theory that shows that the two different effective scales on the boundaries can coexist.
Let us now consider the opposite limit, $`rz`$, that also corresponds to the case when $`z=1/k`$, i.e. the gauge propagator on the $`M_P`$-boundary. In this case we have,
$$G(x,\frac{1}{k};x^{},\frac{1}{k})k\frac{d^4p}{(2\pi )^4}e^{ip(xx^{})}\frac{1}{p^2\mathrm{ln}(p/k)}.$$
(14)
From Eq. (14) we can derive the static potential on the $`M_P`$-boundary:
$$V(r)\frac{k}{r}\frac{1}{\mathrm{ln}(kr)}.$$
(15)
We see that it differs from the Coulomb potential in 4D by a logarithmic term. It means that the gauge coupling on the $`M_P`$-boundary grows, at the tree-level, logarithmically with the energy. This is in contrast with 5D theories in flat space where at the classical level the coupling grows linearly with the energy. In a theory with finite $`R`$ this “running” will be present at energies above the mass of the first KK mode $`m_{\mathrm{KK}}=\pi ke^{k\pi R}`$ (below $`m_{\mathrm{KK}}`$ we have a single massless gauge boson as in 4D). We then have
$$g^2(p>m_{\mathrm{KK}})g^2(pm_{\mathrm{KK}})\frac{\pi kR}{\mathrm{ln}k/p}.$$
(16)
This mild logarithmic evolution of the gauge coupling allows us to go to high energies without entering in the strong coupling regime. We must stress that this “running” of the gauge coupling is a tree-level effect, not a quantum one. As a consequence, it will be universal for the different groups of the SM and will not affect gauge-coupling unification.
From the tree-level behaviour of the propagator in Eq. (11), we learn that the theory remains weakly coupled for $`p<1/z`$. This suggests that the theory can be renormalized as long as we keep our cut-off scale below $`1/z`$, i.e. $`\mathrm{\Lambda }<ke^{ky}`$. Notice that this cut-off scale depends on the position in the extra dimension. This should be expected in a theory with the metric (1), since the effective scale of the 4D space-time at the position $`y`$ is given by $`ke^{ky}`$. Using the cut-off $`\mathrm{\Lambda }=ke^{ky}`$, we can calculate quantum corrections in a very simple way. We just need the 5D propagators for $`r>z`$. For the case of the 5D massless scalar discussed in the previous section, the propagator behaves as in 4D , $`kd^4p/(2\pi )^4(1/p^2)`$, giving then the same one-loop correction to the gauge coupling as in 4D.
## 5 GUTs in a slice of AdS<sub>5</sub>
Let us now proceed to show that theories with gauge bosons in a slice of AdS<sub>5</sub> can have gauge-coupling unification. We will take a top–down approach. We will assume that we have a supersymmetric GUT in the slice of AdS<sub>5</sub> and show that this theory, when broken to the MSSM group, leads to a successful prediction for the gauge couplings at low energies.
As a toy example, let us consider an SU(5) theory. Due to the $`_2`$ orbifold symmetry, the massless gauge sector of this theory consists of a $`N=1`$ vectormultiplet . They contain the SM gauge bosons plus the GUT gauge bosons, $`X`$ and $`Y`$, that complete the SU(5) representation. The KK spectrum consists of $`N=2`$ vectormultiplets with masses $`(n\frac{1}{4})m_{\mathrm{KK}}`$. Let us now consider that on the $`M_P`$-boundary we have a chiral supermultiplet, in the 24 representation of SU(5), whose scalar gets a vacuum expectation value (VEV) equal to $`M_{\mathrm{GUT}}10^{16}`$ GeV (slightly below $`kM_P`$) breaking the SU(5) group down to the MSSM. This can be achieved in the same way as in ordinary 4D SU(5) theories, since our theory on the boundary is 4D $`N=1`$ supersymmetric. It is easy to calculate the KK spectrum of the resulting theory. The $`n=0`$ MSSM gauge bosons remain massless, while the GUT gauge bosons, $`X`$ and $`Y`$, have masses $`M_{X,Y}M_{\mathrm{GUT}}`$ <sup>1</sup><sup>1</sup>1 More precisely, we find that $`M_{X,Y}`$ is determined by the equation $`M_{X,Y}(\mathrm{ln}(M_{X,Y}/2k)+\gamma +1/2)+M_{\mathrm{GUT}}^2/M_{X,Y}=0`$.. The KK mass spectrum ($`n1`$), however, is not modified by the VEV of the $`\mathrm{𝟐𝟒}`$ (up to corrections of $`𝒪(M_n^2/k^2)`$) and therefore the KK modes approximately respect the SU(5) symmetry. Consequently, only the zero modes (as we claimed before) will give a relative one-loop contribution to the SM gauge couplings which, at energies $`\mu `$, is given by
$$\frac{1}{\alpha _i}=\frac{1}{\alpha _j}+\frac{b_ib_j}{2\pi }\mathrm{ln}\frac{M_{\mathrm{GUT}}}{\mu },$$
(17)
where $`b_i`$ is the contribution of the massless modes to the beta-function coefficient of the MSSM gauge group $`i`$. Therefore, in order to have the same predictions for the gauge coupling as in 4D supersymmetric GUTs, we must just demand that the massless states of the theory be those of the MSSM. This will be the case of the gauge sector, as we already explained. For the Higgs sector we can, as usual, assume that they arise from a $`\mathrm{𝟓}`$ and $`\overline{\mathrm{𝟓}}`$ of SU(5). Since we need to have only the SU(2)<sub>L</sub>-doublet light, we will need a mechanism that provides a doublet–triplet mass splitting inside the $`\mathrm{𝟓}`$ and $`\overline{\mathrm{𝟓}}`$. Several mechanisms exist in the literature for 4D. It is not clear if these mechanisms can also work in 5D. Nevertheless, we can just rely on these mechanisms by assuming that the Higgs live on the $`M_P`$-boundary.
Finally, we must implement the matter sector. Since they form complete SU(5) multiplets, $`\overline{\mathrm{𝟓}}`$ and $`\mathrm{𝟏𝟎}`$, they are irrelevant to gauge-coupling unification (they will not contribute at the one-loop level to the relative corrections to the gauge couplings). The matter sector, however, must satisfy important constraints from proton decay. Since there are very light KK modes of the $`X`$ and $`Y`$ bosons ($`m_{\mathrm{KK}}`$ TeV), we must worry about proton-decay operators induced by these modes. If we analyze the $`y`$-dependent wave-function of these modes, however, we find that they are peaked on the TeV-boundary . Therefore, proton-decay constraints can be satisfied by just placing the matter sector on the $`M_P`$-boundary. In this case, even if we sum over the full KK tower of the $`X`$ and $`Y`$ bosons, we obtain that the strength of the dimension-six proton-decay operator is given by
$$\underset{n}{}g_n^2\frac{1}{m_n^2}g^2\frac{\pi kR}{M_{\mathrm{GUT}}^4}_0^{M_{\mathrm{GUT}}}m𝑑m=g^2\frac{\pi kR}{2M_{\mathrm{GUT}}^2},$$
(18)
where $`g_n`$ and $`m_n`$ are respectively the coupling to the $`M_P`$-boundary and the mass of the KK of the $`X,Y`$ bosons that can be derived following Ref. . We see that the result is similar to that in a 4D theory where one finds $`g^2/M_{\mathrm{GUT}}^2`$. The operator (18) is, however, slightly larger in our theory than in 4D theories because of the factor $`\pi kR`$ in Eq. (18). This enhancement is due to the fact that the gauge coupling grows (at tree-level) with the energy according to Eq. (16). Notice that, at the scale $`M_{\mathrm{GUT}}`$, the theory is close to the strong coupling regime. This is why we expect in these GUTs a proton decay rate for $`p\pi e`$ closer to the experimental limit than in 4D GUTs.
Up to now we have just assumed that $`m_{\mathrm{KK}}`$ (approximately the mass of the lightest KK state) is an independent parameter of the theory that we have taken to be close to the weak scale by choosing $`R11/k`$. Nevertheless, it would be interesting to relate $`m_{\mathrm{KK}}`$ with the weak scale. One way to do this is by associating the supersymmetry-breaking scale with $`m_{\mathrm{KK}}`$. A realization of this is given in Ref. . By assuming different boundary conditions for bosons and fermions on the TeV-boundary, we can get a fermion–boson mass splitting of $`𝒪(m_{\mathrm{KK}})`$. This breaks supersymmetry and induces a Higgs mass of $`𝒪(m_{\mathrm{KK}})`$ . If this mass is negative, this will trigger electroweak symmetry breaking. This scenario therefore links the scale $`m_{\mathrm{KK}}`$ with the weak scale.
### 5.1 GUT physics at TeV energies
Although this theory resembles the ordinary 4D supersymmetric GUT, it has very different implications at TeV-energies. While 4D supersymmetric GUTs predict that only the MSSM fields have masses of the order of the weak scale, with a big “desert” up to the GUT scale, our theory has plenty of new physics at the TeV. There are the KK states not only of the SM but also of the GUT fields and graviton. It has been shown in Ref. that the KK modes of the SM gauge bosons have sizeable couplings to the SM fermions living on the $`M_P`$-boundary ($`0.2g`$) and therefore they could be seen as resonances in TeV colliders. On the other hand, the KK modes of the GUT fields couple very weakly to the $`M_P`$-boundary (this is why the proton decay rate is small). These modes, however, can be produced at TeV energies by processes mediated by virtual SM gauge bosons (that live in the 5D bulk and propagate between the two boundaries). At these energies, also graviton KK modes can be produced. In fact, since the effective scale on the TeV-boundary is $`ke^{k\pi R}`$ TeV, quantum gravity or string effects can be important and possible to test.
## 6 Conclusion
We have shown that, in theories with gauge bosons propagating in the 5D bulk of the Randall–Sundrum model (a slice of AdS<sub>5</sub>), the gauge coupling gets logarithmic corrections similar to those in 4D. These theories contain light (TeV) KK excitations, but only the (massless) zero modes contributes to the renormalization of the gauge coupling.
We have proposed a GUT where the gauge bosons live in the 5D bulk, while matter is localized on the $`M_P`$-boundary. The theory has the massless spectrum of the MSSM and predicts the right value for the gauge couplings at low energies. On the other hand, we find, at the TeV scale, KK modes of the GUT fields. These modes are very weakly coupled to the fermions of the SM and consequently proton decay rates are suppressed. Nevertheless, they couple to the SM gauge bosons with sizeable couplings, providing the possibility to test GUTs at TeV colliders.
## Acknowledgements
We thank Jaume Garriga, Tony Gherghetta and Riccardo Rattazzi for very useful discussions. This work has been supported in part by the Spanish CICYT contract AEN99-0766.
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# Estimates of Weil-Petersson volumes via effective divisors
## 1. Introduction and Statement of Results
The total free energy of two dimensional gravity, which is a generating function for certain intersection numbers on the compactified moduli spaces $`\overline{}_{g,n}`$ of stable $`n`$-pointed curves, was conjectured by Witten (and proved by Kontsevich) to satisfy certain KdV equations. This gave new insight in the geometry of those moduli spaces \[WI, KO, DIJ\].
The Mumford class $`\kappa _1`$ on $`\overline{}_{g,0}`$ was shown to be proportional to the cohomology class $`[\omega _{WP}]`$ of the Weil-Petersson form by Wolpert in \[WO\]. Furthermore he showed that the restriction of this class to any component of the compactyfying divisor coincides with the corresponding Weil-Petersson class. Arbarello and Cornalba introduced classes $`\kappa _1`$ on $`\overline{}_{g,n}`$, proved a similar restriction property for these and concluded proportionality on all $`\overline{}_{g,n}`$ \[A-C 2\]. So Weil-Petersson volumes $`\mathrm{vol}(_{g,n})`$ are up to a normalizing factor the intersection numbers
$$V_{g,n}=_{_{g,n}}\kappa _1^{3g3+n}.$$
Recently, in papers by Kaufmann, Manin, Zagier, and Zograf \[K-M-Z, M-Z, ZO3\] a generating function was introduced for intersection numbers of Mumford’s tautological classes \[MU2\], and shown to be equal to the above generating function up to a change of variables. Previously, Zograf had computed the volumes for genus 0,1, and 2 explicitly \[ZO1, ZO2\]. Manin and Zograf also gave estimates of the volume for fixed genus and $`n\mathrm{}`$.
The aim of this note is to study the Weil-Petersson volume of the moduli spaces $`_{g,n}`$ for fixed $`n`$ and large $`g`$. Introducing the decorated Teichmüller space, Penner \[PE\] gave a technique how to integrate top degree differential forms on $`_{g,n}`$, which led to an estimate of the volumes of $`_{g,1}`$ from below with respect to $`g\mathrm{}`$. With these methods, Grushevski \[GR\] recently proved an upper bound for the volume of $`_{g,n}`$ for fixed $`n>0`$ and large $`g`$. For $`n=1`$ his upper estimate has the same order of growth as Penner’s lower estimate.
However, for the classical moduli spaces $`_{g,0}`$ the asymptotics of the volume for $`g\mathrm{}`$ have not been treated.
Here, we give a different approach, which does not require the existence of punctures. On one hand we use the known push-pull type formulas in the spirit of Arbarello and Cornalba \[A-C 1, A-C 2\] to estimate the volume of $`_{g,n+1}`$ from below in terms of the volume of $`_{g,n}`$ for any given $`g`$ and $`n0`$. On the other hand, we base our estimates on the fact that $`\kappa _1`$ is ample and the above restriction property.
In this way it is possible to estimate the volume of the moduli space $`_{g,0}`$ from below in terms of the volume of moduli spaces of lower genus.
We set $`V_{(0,3)}=0`$. The values $`V_{(0,4)}=1`$, $`V_{(0,5)}=5`$, $`V_{(1,1)}=\frac{1}{24}`$, and $`V_{(1,2)}=\frac{1}{8}`$ are known. We prove the following Theorems.
###### Theorem 1.
Let $`2g2+n>0`$ and $`(g,n)(0,4),(1,1)`$. Then
(1)
$$V_{g,n+1}\frac{1}{2}(3g2+n)(7g7+3n)V_{g,n}+\frac{1}{24^gg!}$$
###### Theorem 2.
Let $`g>1`$. Then
(2)
$$V_{g,0}\frac{1}{28}V_{g1,2}+\frac{1}{672}V_{g1,1}+\frac{1}{14}\underset{j=2}{\overset{[g/2]}{}}V_{j,1}V_{gj,1}\frac{1}{28}(V_{\frac{g}{2},1})^2,$$
with $`V_{\frac{g}{2},1}=0`$, if $`g`$ is odd.
Together with the results of Penner and Grushevsky these imply the existence of constants $`0<c<C`$, independent of $`n`$ such that
(3)
$$c^g(2g)!\frac{V_{g,n}}{(3g3+n)!}C^g(2g)!$$
for all fixed $`n0`$ and large $`g`$.
In particular for all $`n0`$
(4)
$$\underset{g\mathrm{}}{lim}\frac{\mathrm{log}\frac{V_{g,n}}{(3g3+n)!}}{g\mathrm{log}g}=2$$
## 2. Proof of the estimates
For any $`g`$ and $`n`$ with $`2g2+n>0`$ the map $`_{g,n+1}_{g,n}`$ forgetting the last puncture is known to extend holomorphically to a map
$$\pi _{n+1}:\overline{}_{g,n+1}\overline{}_{g,n}.$$
For $`n>0`$ it possesses natural sections $`\sigma _j`$; $`j=1,\mathrm{},n`$ (cf. \[A-C 2, sect. 1\]) with corresponding divisors $`D_j`$. Denote the relative dualizing sheaf $`\omega _{\overline{}_{g,n+1}/\overline{}_{g,n}}`$ by $`\omega _{n+1}`$. Let
$$\psi _j:=c_1(\sigma _j^{}\omega _{n+1}),$$
and
$$K:=c_1(\omega _{n+1}(D))H^2(\overline{}_{g,n+1},),$$
where $`D=D_1+\mathrm{}D_n`$.
Finally
$$\kappa _j:=\pi _{n+1}(K^{j+1})H^{2j}(\overline{}_{g,n+1},)$$
for $`j=0,\mathrm{},3g3+n`$.
For $`n=0`$ these are equal to the Mumford classes (denoted also by $`\kappa _j`$). Moreover $`\kappa _1`$ is ample on $`\overline{}_{g,n}`$.
According to Mumford \[MU2\] the classes $`\kappa _j`$ are numerically effective on $`\overline{}_{g,0}`$ for $`j=1,\mathrm{},3g3`$ in the sense that for any complete subvariety $`W\overline{}_{g,0}`$ of dimension $`j`$
$$_W\kappa _j0$$
holds. Also, for $`j_1+\mathrm{}+j_k=3g3`$
(5)
$$_{\overline{}_{g,0}}\kappa _{j_1}\mathrm{}\kappa _{j_k}0$$
following \[HA2\].
We shall need the following formulas for cohomology classes to be found in \[A-C 2, (1.7), (1.9), and (1.10)\].
(6)
$$\pi _n(\psi _1^{a_1}\mathrm{}\psi _{n1}^{a_{n1}}\psi _n^{a_n+1})=\psi _1^{a_1}\mathrm{}\psi _{n1}^{a_{n1}}\kappa _{a_n}$$
$$\text{for}a_j0$$
(7)
$$\pi _n(\psi _1^{a_1}\mathrm{}\psi _{n1}^{a_{n1}})=$$
$$\underset{j;a_j>0}{}\psi _1^{a_1}\mathrm{}\psi _{j1}^{a_{j1}}\psi _j^{a_j1}\psi _{j+1}^a\mathrm{}\psi _{n1}^{a_{n1}}$$
(8)
$$\kappa _a=\pi _{n+1}^{}(\kappa _a)+\psi _{n+1}^a\text{on}\overline{}_{g,n+1}$$
(9)
$$\kappa _0=2g2+n.$$
###### Lemma 1.
Let $`m_j0`$ be integers such that $`_{j=1}^kjm_j=3g3+n`$ with $`n0`$. Then
(10)
$$_{_{g,n}}\kappa _1^{m_1}\mathrm{}\kappa _k^{m_k}0.$$
###### Proof.
The above equation (5) is the statement for $`n=0`$. We proceed by induction on $`n`$. Assume (10) for some $`n0`$. Then
$$_{_{g,n+1}}\kappa _1^{m_1}\mathrm{}\kappa _k^{m_k}=$$
$$_{_{g,n}}\pi _{n+1}((\pi _{n+1}^{}\kappa _1+\psi _{n+1})^{m_1}(\pi _{n+1}^{}\kappa _2+\psi _{n+1}^2)^{m_2}$$
$$\mathrm{}(\pi _{n+1}^{}\kappa _k+\psi _{n+1}^k)^{m_k})$$
Since
$$_{_{g,n}}\pi _{n+1}\left(\pi _{n+1}^{}(\kappa _1^{j_1}\mathrm{}\kappa _k^{j_k})\psi _{n+1}^{\mathrm{}+1}\right)=_{_{g,n}}\kappa _1^{j_1}\mathrm{}\kappa _k^{j_k}\kappa _{\mathrm{}},$$
the above integral can be expressed as a sum of non-negative terms. ∎
Proof of Theorem 1. We first note that for $`g>0`$
$$_{_{g,n+1}}\psi _{n+1}^{3g2+n}=_{_{g,n+1}}\psi _1^{3g2+n}=_{_{g,n}}\psi _1^{3g3+n}=\mathrm{}=_{_{g,1}}\psi _1^{3g2}$$
by (7). These integrals are known to be equal to $`1/(24^gg!)`$ (cf. \[F-P\]). For $`g=0`$
$$_{_{0,n+1}}\psi _{n+1}^{n2}=_{_{0,4}}\psi _1=_{_{0,3}}\kappa _0=1.$$
By Lemma 1, for $`j<3g2+n`$
$$_{_{g,n+1}}\pi _{n+1}^{}(\kappa _1)^j\psi _{n+1}^{3g2+nj}=_{_{g,n}}\kappa _1^j\kappa _{3g3+nj}0,$$
hence
$$_{_{g,n+1}}\kappa _1^{3g2+n}=_{_{g,n+1}}\left(\pi _{n+1}^{}(\kappa _1)+\psi _{n+1}\right)^{3g2+n}=$$
$$\underset{j=0}{\overset{3g2+n}{}}\left(\genfrac{}{}{0pt}{}{3g2+n}{j}\right)_{_{g,n+1}}(\pi _{n+1}^{}(\kappa _1))^j\psi _{n+1}^{3g2+nj}$$
$$(3g2+n)_{_{g,n}}\kappa _1^{3g3+n}\kappa _0+$$
$$\frac{1}{2}(3g2+n)(3g3+n)_{_{g,n}}\kappa _1^{3g4+n}\kappa _1+$$
$$_{_{g,n+1}}\psi _{n+1}^{3g2+n}=$$
$$((3g2+n)(2g2+n)+\frac{1}{2}(3g2+n)(3g3+n))_{_{g,n}}\kappa _1^{3g2+n}+$$
$$_{_{g,n+1}}\psi _{n+1}^{3g2+n}=$$
$$\frac{1}{2}(3g2+n)(7g7+3n)_{_{g,n}}\kappa _1^{3g3+n}+\frac{1}{24^gg!}.$$
. ∎
For any family $`f:𝒞S`$ of stable curves $`det(f_{}\omega _{𝒞/S})`$ is a line bundle over $`S`$, where $`\omega _{𝒞/S}`$ denotes the relative dualizing sheaf. These determinant sheaves give rise to a $``$-divisor on $`\overline{}_{g,0}`$, which is usually denoted by $`\lambda `$. In a similar way the singular fibers of the above family define devisors, which also give rise to $``$-divisors on $`\overline{}_{g,0}`$, usually denoted by $`\delta _i`$. The irreducible components of the divisor at infinity on $`\overline{}_{g,0}`$ are denoted by $`\mathrm{\Delta }_i`$, $`i=0,\mathrm{},[g/2]`$ with classes
$$[\mathrm{\Delta }_i]=\delta _i\text{f}ori1\text{a}nd[\mathrm{\Delta }_1]=2\delta _1.$$
These are characterized as follows:
* The generic point of $`\mathrm{\Delta }_0`$ corresponds to an irreducible, stable curve of genus $`g1`$ with one ordinary double point. In fact there is a generically 2:1 surjective holomorphic map $`\overline{}_{g1,2}\mathrm{\Delta }_0`$.
* For $`i=1,\mathrm{},[g/2]`$ the generic points of $`\mathrm{\Delta }_i`$ correspond to stable curves with one ordinary double point and two irreducible components of genus $`i`$ and $`gi`$ resp. There exists a surjective holomorphic map
$$\overline{}_{i,1}\times \overline{}_{gi,1}\mathrm{\Delta }_i,$$
which is generically 1:1 for $`ig/2`$ and 2:1 for $`i=g/2`$
(cf. \[H-M\]).
###### Theorem 3.
Let
$$D=p\lambda \underset{j=0}{\overset{[g/2]}{}}q_j\delta _j;p,q_j>0$$
be an effective $``$-divisor on $`\overline{}_{g,0}`$ such that
$$\mu _j=\frac{12q_jp}{p}>0.$$
Then
$$V_{g,0}>\frac{\mu _0}{2}V_{g1,2}+\frac{\mu _1}{48}V_{g1,1}+\underset{j=2}{\overset{[g/2]}{}}\mu _jV_{j,1}V_{gj,1}\mu _{g/2}(V_{g/2,1})^2$$
(where $`\mu _{g/2}=V_{g/2,1}=0`$, if $`g`$ is odd).
###### Proof.
According to \[MU1\]
$$\kappa _1=12\lambda \underset{j=0}{\overset{[g/2]}{}}\delta _j.$$
We want to write the divisor in the form
$$D=\alpha \kappa _1\beta _j\delta _j$$
for $`\alpha ,\beta _j>0`$. This gives
$$\alpha =\frac{p}{12};\beta _j=q_j\frac{p}{12}>0.$$
Now
$$0<\kappa _1^{3g4}D=\alpha \kappa _1^{3g3}\beta _j\kappa _1^{3g4}\delta _j.$$
We use the restriction property of $`\kappa _1`$ with respect to $`\mathrm{\Delta }_j`$.
On $`\overline{}_{g1,2}`$ the class $`\kappa _1`$ is invariant under the action of $`_2`$ exchanging the punctures. So under the natural map $`\overline{}_{g1,2}\mathrm{\Delta }_0`$ it descends to the restriction of the class $`\kappa _1`$ on the ambient space $`\overline{}_{g,0}`$. Now
$$\kappa _1^{3g4}\delta _0=_{\mathrm{\Delta }_0}\kappa _1^{3g4}=\frac{1}{2}V_{g1,2}.$$
In a similar way, we get
$$\kappa _1^{3g4}\delta _1=\frac{1}{2}V_{1,1}V_{g1,1},$$
and for $`i>1`$
$$\kappa _1^{3g4}\delta _i=V_{i,1}V_{gi,1}$$
with an extra factor $`1/2`$ for $`(V_{\frac{g}{2},1})^2`$, if $`g`$ is even. Also we use $`V_{1,1}=1/24`$. ∎
Remark: Combining push-pull formulas and computations of intersections of powers of $`\kappa _1`$ with various effective divisors, one can also estimate the intersection numbers $`V_{g,n}`$ from above in terms of $`V_{g1,n+2}`$, and the numbers $`V_{j,\mathrm{}}`$ with $`j=g`$, $`\mathrm{}<n`$ or $`j<g`$, $`\mathrm{}n+1`$.
The above Theorem 2 follows, since for every rational $`\epsilon >0`$ the $``$-divisor $`D=(11.2+\epsilon )\lambda \delta `$ is ample \[MU1\], where $`\delta =\delta _j`$.
A stronger estimate for $`g23`$ follows from the fact that $`\overline{}_{g,0}`$ has positive Kodaira dimension according to the theorems of Eisenbud, Harris, and Mumford \[E-H, HA1, H-M\]. The equality
$$K_{\overline{}_{g,0}}=13\lambda 2\delta _03\delta _12\underset{j=2}{\overset{[g/2]}{}}\delta _j$$
in $`\mathrm{Pic}(\overline{}_{g,0})`$ was proved in \[H-M\].
This implies for $`g23`$
$$V_{g,0}>\frac{11}{26}V_{g1,2}+\frac{23}{624}V_{g1,1}+\frac{11}{13}\underset{j=2}{\overset{[g/2]}{}}V_{j,1}V_{gj,1}\frac{11}{26}(V_{g/2,1})^2.$$
Further computations of effective divisors in terms of $`\lambda `$ and $`\delta _j`$ were provided in \[H-M, E-H\]. All of these divisors satisfy the hypothesis of Theorem 3 and can be used for new estimates of the Weil-Petersson volumes.
We finally discuss, how to arrive at the asymptotic estimates (3).
A rough estimate following from Theorem 1 is
$$V_{g,n}V_{g,1},$$
which, together with Penner’s lower estimate for $`V_{g,1}`$, already implies the existence of a constant $`c>0`$, independent of $`n1`$, such that for large $`g`$
$$\frac{V_{g,n}}{(3g3+n)!}c^g(2g)!.$$
The corresponding upper estimate is due to Grushevski, so that (3) follows for $`n1`$.
For $`n=0`$ and $`g>1`$ the above Theorem 1 and Theorem 2 give
$$\frac{1}{672}V_{g1,1}V_{g,0}<\frac{2}{(3g2)(7g7)}V_{g,1}.$$
Again, with \[GR, PE\] these inequalities yield the following corollary.
###### Corollary 1.
There exist constants $`0<\stackrel{~}{c}<\stackrel{~}{C}`$ such that for $`g0`$
$$\stackrel{~}{c}^g(2g)!\frac{V_{g,0}}{(3g3)!}\stackrel{~}{C}^g(2g)!.$$
This paper was written, while the second named author was visiting the University of Marburg. He would like to thank the Department of Mathematics for its hospitality.
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# Quantised bulk fields in the Randall-Sundrum compactification model11footnote 1To appear in Physics Letters B.
(May 1, 2000)
## Abstract
The quantisation of a scalar field in the five-dimensional model suggested by Randall and Sundrum is considered. Using the Kaluza-Klein reduction of the scalar field, discussed by Goldberger and Wise, we sum the infinite tower of modes to find the vacuum energy density. Dimensional regularisation is used and we compute the pole term needed for renormalisation, as well as the finite part of the energy density. Some comments are made concerning the possible self-consistent determination of the radius.
For about the past year there has been a resurgence of interest in the old Kaluza-Klein idea that spacetime may have some extra dimensions. One reason for this renewed interest lies in the solution it provides to the hierarchy between the Planck scale and the electroweak scale . In place of the compact extra dimensions used in , Randall and Sundrum have proposed a 5-dimensional model in which the extra dimension has an orbifold compactification with two 3-branes with opposite tensions sitting at two orbifold fixed points. An essential feature of this model is that the 5-dimensional geometry is not a direct product of the 4-dimensional spacetime and the extra dimension. The line element, which is a solution to the 5-dimensional Einstein equations for a particular 3-brane source, is
$$ds^2=e^{2kr|\varphi |}\eta _{\mu \nu }dx^\mu dx^\nu r^2d\varphi ^2.$$
(1)
Here $`x^\mu `$ are the usual coordinates on 4-dimensional spacetime, and $`|\varphi |\pi `$ with the points $`(x^\mu ,\varphi )`$ and $`(x^\mu ,\varphi )`$ identified. The 3-branes sit at $`\varphi =0,\pi `$. $`k`$ is a constant and $`r`$ is also constant and associated with the size of the extra dimensions. The constant $`r`$ is an arbitrary constant of integration found when solving the Einstein equations. Randall and Sundrum showed that for $`kr10`$, essentially because of the exponential appearing in (1), that TeV mass scales could be produced in the 4-dimensional theory from Planck scale quantities in the higher-dimensional theory. Furthermore the Randall-Sundrum model has possible phenomenological consequences that put higher-dimensional models within the experimental range. (See also .)
In the Randall-Sundrum model, the radius $`r`$ of the extra dimension is not fixed. This was also the case for the older Kaluza-Klein theories which were based on the simpler geometry of a direct product of 4-dimensional Minkowski spacetime and a homogeneous space, often chosen to be a sphere. It was realised by Candelas and Weinberg that quantum effects from matter fields, or gravity, could be used to fix the size of the extra dimensions and obtain a self-consistent solution. In view of the interest generated by the Randall-Sundrum scenario, it is reasonable to ask if the radius of the extra dimension in their model can be determined in a similar way. This is the prime motivation for the present paper in which we calculate the vacuum (or Casimir) energy of a quantised scalar field on the Randall-Sundrum background, since this is a first step in all self-consistency calculations. As we will see, the orbifold nature of the Randall-Sundrum spacetime makes this calculation considerably more difficult than was the case in the original Kaluza-Klein models based on direct product spacetimes.
It was shown in that a scalar field on the higher dimensional spacetime could be considered, from the 4-dimensional point of view, in the usual Kaluza-Klein way as an infinite tower of scalar fields whose masses are quantised. The action can be written as
$$S=\frac{1}{2}\underset{n}{}d^4x\left(^\mu \phi _n_\mu \phi _nm_n^2\phi _n^2\right),$$
(2)
where $`m_n^2`$ gives the mass for the $`n^{th}`$ mode. All dependence on the extra fifth dimension has been integrated out here to leave a 4-dimensional theory. For the Randall-Sundrum model (1) Goldberger and Wise showed that the masses $`m_n`$ are given by solutions to the transcendental equation
$$0=y_\nu (ax_n)j_\nu (x_n)j_\nu (ax_n)y_\nu (x_n)$$
(3)
where we have defined $`a=e^{\pi kr}`$ and $`m_n=kax_n`$, with $`x_n`$ the $`n^{th}`$ positive solution to (3). The functions $`j_\nu `$ and $`y_\nu `$ are shorthand for the following combinations of Bessel functions :
$`j_\nu (z)`$ $`=`$ $`2J_\nu (z)+zJ_\nu ^{}(z)`$ (4)
$`y_\nu (z)`$ $`=`$ $`2Y_\nu (z)+zY_\nu ^{}(z).`$ (5)
The order $`\nu `$ of the Bessel functions is given by
$$\nu =\sqrt{4+\frac{m^2}{k^2}},$$
(6)
where $`m`$ is the mass of the 5-dimensional scalar field. The derivation of all of the results (26) is described very clearly in . The fact that $`a=e^{\pi kr}10^{17}`$, which is necessary for a solution to the hierarchy problem , and that as a consequence $`x_11`$ allows the conclusion that the lightest modes have masses of the order of a few TeV. This is in marked contrast to the older Kaluza-Klein picture in which a whole spectrum of light and unobserved particles could appear.
What we wish to do in this paper is calculate the quantum vacuum polarisation effects from the infinite tower of massive particles in (2). As mentioned earlier, similar calculations were performed in the older Kaluza-Klein models based on direct products of Minkowski spacetime with homogeneous spaces, such as spheres (see for example and references therein) with the aim of obtaining self-consistent solutions for the size of the extra dimensions. It seems a natural progression of this earlier work to ask if similar calculations can be done in the new Randall-Sundrum model. A complication which is present in this model is that it is not possible to obtain a closed form expression for the mass modes $`m_n`$, unlike the case of direct products of flat spacetime with homogeneous spaces.<sup>2</sup><sup>2</sup>2Of course it is trivial to obtain approximate solutions to (3) using the asymptotic form of the Bessel functions of large argument, but these are useless for our purposes. Nevertheless we will describe a procedure for evaluating the vacuum energy using only the basic properties of (3) irrespective of the lack of an explicit form for $`m_n`$.
In this paper we will concentrate on the one-loop vacuum energy density for the theory whose Lagrangian is given by (1) which is
$`V^{(1)}`$ $`=`$ $`{\displaystyle \frac{\mathrm{}^ϵ}{2}}{\displaystyle \underset{n}{}}{\displaystyle \frac{d^{3+ϵ}}{(2\pi )^{3+ϵ}}(p^2+m_n^2)^{1/2}}`$ (7)
$`=`$ $`{\displaystyle \frac{1}{16\pi ^2}}\left({\displaystyle \frac{\mathrm{}^2}{4\pi }}\right)^{ϵ/2}\mathrm{\Gamma }(2ϵ/2){\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}(m_n^2)^{2+ϵ/2}.`$
Here we are using dimensional regularisation with $`\mathrm{}`$ the renormalisation length. ($`\zeta `$-function regularisation can be used to the same end.) The $`ϵ0`$ limit is understood here, but before this limit is taken it is necessary to evaluate the sum over modes in (7). Since $`m_n^2=(ka)^2x_n^2`$, we need to evaluate
$$v(ϵ)=\mathrm{\Gamma }(2ϵ/2)\underset{n=1}{\overset{\mathrm{}}{}}x_n^{4+ϵ},$$
(8)
with $`x_n`$ the $`n^{th}`$ positive solution to (3). Then
$$V^{(1)}=\frac{(ka)^4}{16\pi ^2}\left(\frac{k^2a^2\mathrm{}^2}{4\pi }\right)^{ϵ/2}v(ϵ).$$
(9)
will give the quantum vacuum energy (in the limit $`ϵ0`$).
The sum defining $`v(ϵ)`$ can be evaluated in the usual way by first converting it into a contour integral in the region of the complex $`ϵ`$-plane where the sum converges (which is easily seen to be $`\mathrm{}ϵ<5`$ in this case) and then deforming the contour so that we can perform an analytic continuation back to a neighbourhood of $`ϵ=0`$. Similar calculations have been done for the Casimir effect in a spherical shell . It is possible to show that
$$v(ϵ)=\frac{1}{\mathrm{\Gamma }(3+ϵ/2)}_0^{\mathrm{}}𝑑zz^{4+ϵ}\frac{d}{dz}\mathrm{ln}P_\nu (z),$$
(10)
where
$$P_\nu (z)=\frac{2}{\pi }\left[k_\nu (z)i_\nu (az)k_\nu (az)i_\nu (z)\right].$$
(11)
$`k_\nu (z)`$ and $`i_\nu (z)`$ are given by expressions like (4) and (5) but involve the modified Bessel functions $`K_\nu (z)`$ and $`I_\nu (z)`$ in place of $`J_\nu (z)`$ and $`Y_\nu (z)`$.
The remaining task now is to perform the analytic continuation of $`v(ϵ)`$ to a region about $`ϵ=0`$. This can be done by studying the behaviour of the integrand at both small and large $`z`$. The analytic continuation is most easily performed by splitting up the integration range in (10) and defining
$$v(ϵ)=v_1(ϵ)+v_2(ϵ),$$
(12)
where
$$v_1(ϵ)=\frac{1}{\mathrm{\Gamma }(3+ϵ/2)}_0^1𝑑zz^{4+ϵ}\frac{d}{dz}\mathrm{ln}P_\nu (z),$$
(13)
with $`v_2(ϵ)`$ given by a similar expression but with the integration going from 1 to $`\mathrm{}`$. Use of the small $`z`$ behaviour of the Bessel functions shows that $`v_1(ϵ)`$ is analytic at $`ϵ=0`$. This means that we may concentrate on $`v_2(ϵ)`$ whose analytic continuation is less trivial. The basic idea is to study how the integral defining $`v_2(ϵ)`$ diverges as $`ϵ0`$, and by adding and subtracting appropriate terms enable the $`ϵ0`$ limit to be taken everywhere apart from a possible pole term.
The impediment to the convergence of $`v_2(ϵ)`$ at $`ϵ=0`$ is the behaviour of the integrand for large $`z`$. We can use the known asymptotic behaviour of the Bessel functions to calculate the terms we need to add and subtract to let $`ϵ=0`$. We will define
$`i_\nu (z)`$ $`=`$ $`\sqrt{{\displaystyle \frac{z}{2\pi }}}e^z\mathrm{\Sigma }_\nu (z),`$ (14)
$`k_\nu (z)`$ $`=`$ $`\sqrt{{\displaystyle \frac{\pi z}{2}}}e^z\mathrm{\Sigma }_\nu (z),`$ (15)
with $`\mathrm{\Sigma }_\nu (z)1`$ as $`z\mathrm{}`$. This means that the second term in $`P_\nu (z)`$ is the most divergent (see (11)) as $`z\mathrm{}`$. We find
$`v_2(ϵ)`$ $`=`$ $`{\displaystyle \frac{1}{8}}{\displaystyle \frac{1}{10}}(1a)+{\displaystyle \frac{1}{2}}{\displaystyle _1^{\mathrm{}}}𝑑zz^4{\displaystyle \frac{d}{dz}}\mathrm{ln}\left[1{\displaystyle \frac{k_\nu (z)i_\nu (az)}{i_\nu (z)k_\nu (az)}}\right]`$ (16)
$`+{\displaystyle \frac{1}{\mathrm{\Gamma }(3+ϵ/2)}}{\displaystyle _1^{\mathrm{}}}𝑑zz^{4+ϵ}{\displaystyle \frac{d}{dz}}\mathrm{ln}\left[\mathrm{\Sigma }_\nu (z)\mathrm{\Sigma }_\nu (az)\right]+𝒪(ϵ).`$
To obtain this result it is necessary to assume $`\mathrm{}ϵ<5`$, and perform an analytic continuation to $`ϵ=0`$. The first integral in (16) is easily seen to be finite, and we only need to take care with the second integral. Define, for large $`z`$,
$$\mathrm{ln}\mathrm{\Sigma }_\nu (z)\underset{k=1}{\overset{\mathrm{}}{}}\frac{d_k}{z^k}$$
(17)
for some coefficients $`d_k`$. By adding and subtracting the first five terms in the asymptotic expansion we can end up with a result which is finite at $`ϵ=0`$. It is easy to see that
$`v_2(ϵ)`$ $`=`$ $`{\displaystyle \frac{1}{8}}{\displaystyle \frac{1}{10}}(1a)+{\displaystyle \frac{1}{2}}{\displaystyle _1^{\mathrm{}}}𝑑zz^4{\displaystyle \frac{d}{dz}}\mathrm{ln}\left[1{\displaystyle \frac{k_\nu (z)i_\nu (az)}{i_\nu (z)k_\nu (az)}}\right]`$ (18)
$`+{\displaystyle \frac{1}{2}}{\displaystyle _1^{\mathrm{}}}𝑑zz^4{\displaystyle \frac{d}{dz}}\left\{\mathrm{ln}\left[\mathrm{\Sigma }_\nu (z)\mathrm{\Sigma }_\nu (az)\right]{\displaystyle \underset{k=1}{\overset{5}{}}}\left(1+{\displaystyle \frac{(1)^k}{a^k}}\right){\displaystyle \frac{d_k}{z^k}}\right\}`$
$`+{\displaystyle \frac{1}{6}}\left(1{\displaystyle \frac{1}{a}}\right)d_1+{\displaystyle \frac{1}{2}}\left(1+{\displaystyle \frac{1}{a^2}}\right)d_2+{\displaystyle \frac{3}{2}}\left(1{\displaystyle \frac{1}{a^3}}\right)d_3`$
$`+{\displaystyle \frac{4}{ϵ\mathrm{\Gamma }(3+ϵ/2)}}\left(1+{\displaystyle \frac{1}{a^4}}\right)d_4{\displaystyle \frac{5}{2}}\left(1{\displaystyle \frac{1}{a^5}}\right)d_5.`$
The coefficients $`d_k`$ defined in (17) are given in Table 1. The only task remaining is the evaluation of the two integrals, which must be done numerically.
Because $`d_40`$ the vacuum energy density contains a pole term as $`ϵ0`$. The origin of this divergence lies in the orbifold nature of the Randall-Sundrum spacetime. Without the orbifold fixed points the result for the vacuum energy density would have been finite (after regularisation). Similar divergences have been discussed on manifolds with conical singularities and arise in calculations involving black hole entropy. In our case the pole part of $`V^{(1)}`$ is given by
$$\mathrm{P}.\mathrm{P}.\{V^{(1)}\}=\frac{(1+a^4)}{1024\pi ^2ϵ}(16m^4+24m^2k^2+25k^4).$$
(19)
This result has used the expression for $`d_4`$ in Table 1, and the definition of $`\nu `$ in (6). If we had used $`\zeta `$-function regularisation, in place of the above pole term we would have found a non-zero result for $`\zeta (0)`$ given by an expression like (19). Even if we use the fact that $`a<<1`$, there is still a pole term present in the vacuum energy density with a non-zero coefficient. It is also noteworthy that the pole term is present even if $`m=0`$ corresponding to an initially massless scalar field.
The evaluation of the remaining parts of the vacuum energy involves numerical calculations. There is no impediment to doing these calculations for general values of $`m,k,a`$; however, because the main purpose of the present paper is just to illustrate the general method, we will make life simpler for ourselves. If, following , we are interested in $`mk1`$ (in Planck units) we can take $`m/k=3/2`$ leading to (from (7)) $`\nu =5/2`$. This simplifies all of the Bessel functions to powers and exponentials:
$`k_{5/2}(z)`$ $`=`$ $`{\displaystyle \frac{1}{4}}\sqrt{2\pi }z^{5/2}(2z^3+3z^2+3z+3)e^z,`$ (20)
$`i_{5/2}(z)`$ $`=`$ $`(2\pi )^{1/2}z^{5/2}\left[(2z^3+3z)\mathrm{cosh}z3(1+z^2)\mathrm{sinh}z\right].`$ (21)
The other complication present is the extreme smallness of the parameter $`a`$; but because all integrals are now finite it is possible to expand in powers of $`a`$ in a straightforward way and obtain the leading order behaviour (in $`a`$) of the finite part of $`V^{(1)}`$, which we call $`\mathrm{F}.\mathrm{P}.\{V^{(1)}\}`$. After some calculation we find
$`\mathrm{F}.\mathrm{P}.\{V^{(1)}\}`$ $`=`$ $`{\displaystyle \frac{k^4}{16\pi ^2}}\{{\displaystyle \frac{81}{64}}\mathrm{ln}(k^2\mathrm{}^2)+5.84858058+{\displaystyle \frac{15}{16}}a{\displaystyle \frac{3}{8}}a^2`$ (22)
$`a^4[{\displaystyle \frac{81}{64}}\mathrm{ln}(k^2a^2\mathrm{}^2)2.293923]+𝒪(a^5)\}.`$
This completes our evaluation of the vacuum energy density. Before it can be used to examine the problem of self-consistency for the Randall-Sundrum spacetime there are several important things that remain to be done. The first is to perform a careful renormalization analysis. It should be possible to do this following analogous calculations in the case of conical singularities . A second thing that is essential to consider is the effect of the induced gravity term which will arise, as first discussed in , and whose importance in self-consistency was noted by Candelas and Weinberg . Ideally it will be possible to perform the analogous calculations for quantum gravity and study the role of self-consistency in the Randall-Sundrum model. (Some work on the Kaluza-Klein reduction of fields of spin other than zero has been done in .) Whether or not it is possible to fix the value of the arbitrary compactification radius $`r`$ by such a calculation, and provide an alternative to the calculation in , remains to be seen. Further aspects of the calculations will be given elsewhere .
Acknowledgements: I would like to thank A. Flachi for helpful discussions, and J. S. Dowker for supplying some references for heat kernels on cones.
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# Constraints on source and lens parameters from microlensing variability in QSO 0957+561 A,B
## 1 Introduction
The double quasar Q0957+561 A,B was the first discovered multiple image gravitational lens, the first to have a measured time delay and to produce a gravitational lens determination of the Hubble parameter, and the first system in which a microlensing effect was seen in the observational brightness record (Vanderriest et al. 1989). This was not unexpected since such microlensing effects had already been predicted by Chang & Refsdal (1979), see also Kayser et al. (1986), hereafter referred to as KRS, and Schneider & Weiss (1987). With the time delay confirmed by the Vanderriest et al. report, Schild & Smith (1991) noted evidence for fine structure in the microlensing light-curve. In this paper we do not concern ourselves with this reported fine structure, and focus instead on the long-term microlensing trends.
After a period of controversy about the correct time delay value, Kundić et al. (1996) observed an unusually distinct event in the light curve of image A which repeated in image B 417 days later. This time delay was very close to the values obtained earlier by Vanderriest et al. (1989), Schild & Smith (1991) and Pelt et al. (1996), and has later been confirmed with high precision by other workers, see for instance Pelt et al. (1998) who found a time delay of 416.3 days. A precise value of the time delay is necessary in order to subtract out the intrinsic quasar brightness fluctuations on all time scales, and to find the microlensing residual.
The most accurate determination of microlensing variability in Q0957+561 was reported by Pelt et al. (1998). This variability is shown in Fig. 1 (their Fig. 9). Two principal features are seen in the 15 year microlensing light-curve: a rise of 0.25 mag during a period of about 5 years (1982-1986) with a maximum slope of 0.07 mag per year, and a quiet phase of about 8 years (1988-1996) with a variability less than 0.05 mag. The original data in Pelt et al. (1998) show a large scatter which has sometimes been suggested to evidence microlensing on time scales of 100 days and shorter (Schild 1996) but the existence of such low amplitude and rapid fluctuations is complicated by questions of the accuracy of the data; we therefore restrict our analysis to the less controversial long-term microlensing. Although gaps as long as 300 days are present in the long-term microlensing record, segments where the brightness record is more intensively sampled seem to show that amplitudes on sub-year timescales are less than 5% and they average away as shown convincingly in the Pelt et al. (1998) Fig. 9 plot.
This paper is concerned with understanding the two features in the light-curve mentioned above, from models of the gravitational lens and microlensing from stars or other compact objects presumed to lie within the lens galaxy G1. For typical quasar accretion disc sizes it is possible that statistical effects from numerous stars projected in front of the quasar must be considered; a statistical microlensing theory for such a case can be found in Refsdal & Stabell (1991 (RS1), 1993 (RS2) and 1997).
From macro-lens modeling of Q0957+561 one finds the lensing optical depth for the A and B image to be $`\kappa _\mathrm{A}=0.22`$ and $`\kappa _\mathrm{B}=1.24`$, respectively, and with shear terms $`\gamma _\mathrm{A}=\pm 0.17`$ and $`\gamma _\mathrm{B}=\pm 0.9`$ (Lehar, private communication; see also Schmidt & Wambsganss (1998) for similar estimates and a definition of the shear). Our source model is circular with a Gaussian distributed surface luminosity and “radius” $`R`$.
We shall treat three different mass models in our microlensing simulations:
Case 1: All the lensing mass is in identical compact micro-lenses with mass $`M`$.
Case 2: 10% of the lensing mass is in identical compact lenses with mass $`M`$ and 90% in an evenly distributed continuum.
Case 3: 10% of the lensing mass is in compact lenses with $`1M_{}`$ and 90% in compact lenses with mass $`M<M_{}`$.
Our approach is to compare simulated microlensing light-curves with the observed microlensing variability. Because the lens mass $`M`$ and the source radius $`R`$ both affect the amplitude and time scale of variability, we find some interesting constraints for these two parameters in the lens system Q0957+561.
Schmidt & Wambsganss (1998) have examined the available parameter space based upon an observed quiet period of 160 days with a variability less than 0.05 mag.
The results of our study are based on the microlensing variability during 15 years, including a long quiet phase as well as a significant event and naturally gives stronger constraints.
## 2 Observed microlensing variability
We make use of the observed microlensing light-curve covering a time span $`T=15`$ years, see Fig.1. The curve is given by:
$$m(t)=m_\mathrm{A}m_\mathrm{B}^\tau $$
(1)
where $`m_\mathrm{B}^\tau `$ is the magnitude of the B image shifted by the time delay $`\tau =416.3`$ days (see Pelt et al. 1998). Two significant microlensing features are found:
1) A “quiet” period with a variability $`\delta m=m_{\mathrm{max}}m_{\mathrm{min}}`$ less than $`0.05`$ mag lasting $`\delta t=8`$ years.
2) A variation of $`\mathrm{\Delta }m=0.25`$ mag during a time interval $`\mathrm{\Delta }t=5`$ years, with a maximum slope of about 0.07 mag/year (“the event”).
We shall treat these two features independently and investigate by means of simulations how lens and source parameters can be constrained. The results will be demonstrated in “exclusion diagrams”, see Figs. 2 and 3.
Since the optical depth is much larger for the B image than for A, the microlensing variability in the A image is likely to be much smaller than in the B image. We shall therefore in this paper neglect a possible microlensing variability in the A image and assume that the observed microlensing variability comes only from B. From test calculations we find that this approximation has a rather small influence on the constraints obtained (see Discussion).
## 3 Case 1: all mass in identical lenses
We shall first choose all mass to be in lenses with mass $`M`$. This may not be very realistic, but it allows a relatively simple discussion with only two free parameters ($`M`$ and source radius $`R`$), and it gives us an understanding which is quite useful also for more complex mass distributions. We simulated light-curves for the B image with effective transverse motion of the source parallel to the shear and perpendicular to the shear $`(\gamma _\mathrm{B}=\pm 0.9)`$. The results for these two cases are always sufficiently similar so that an intermediate case easily can be interpolated. For practical reasons we shall instead of $`R`$ use the normalized source radius $`r=R/R_\mathrm{E}`$ as a free parameter, where $`R_\mathrm{E}`$ is the Einstein radius (for mass $`M`$), projected into the source plane
$$R_\mathrm{E}=\sqrt{\frac{4GM}{c^2}\frac{D_{\mathrm{ds}}D_\mathrm{s}}{D_\mathrm{d}}}=4.810^{16}\sqrt{\frac{M}{M_{}h_{60}}}\text{cm}$$
(2)
Here the D’s are angular size distances and $`h_{60}`$ is the dimension-less Hubble parameter, $`h_{60}=H_0/(60`$ km s<sup>-1</sup> Mpc$`{}_{}{}^{1})`$. We have used a cosmological model with $`\mathrm{\Omega }=0.5`$ and $`\mathrm{\Lambda }=0`$.
For randomly distributed compact lenses we simulate light-curves for different values of the normalized source radius $`r`$, using well known ray tracing techniques, see KRS and Schneider & Weiss (1987). We obtain for each value chosen for $`r`$ the magnitude of the source as a function of the normalized source position: $`m=m(y)`$, where $`y=\eta /R_\mathrm{E}`$ and $`\eta `$ is the length coordinate along the source track. For each value of $`r`$ we get the light-curve for different (but still all identical) $`M`$-values by a transformation of $`y`$ to the time $`t`$ measured by the observer. For our system we find
$$t=25yv_{600}^1\sqrt{\frac{M}{M_{}h_{60}}}\text{years}$$
(3)
Here $`v_{600}=V/(600`$ km s$`{}_{}{}^{1})`$ where $`V`$ is the effective transverse velocity of the source assumed to be constant, with time measured by the observer, see KRS. When not stated otherwise we shall use $`v_{600}=1`$ and $`h_{60}=1`$ in this paper.
Since $`ty\sqrt{M}`$ the light-curves are simply stretched out in time for larger values of $`M`$. Hence, it is not necessary to calculate light-curves for different $`M`$-values (for a given $`r`$). We have therefore simulated light-curves only for various $`r`$-values ($`r`$=0.02, 0.05, 0.1, 0.3, 1, 3, 5, 10, 20 and 50). The length of the light-curves corresponds in most cases to at least 1000 years for the allowed masses, see Eqs. (3) and (4). We first discuss the constraints on $`M`$ and $`R`$ (and $`r`$) which can be obtained from the quiet phase of 8 years.
i) Quiet phase constraints
For each value of the normalized source radius $`r`$ we calculate the microlensing light-curve $`m(y)`$ for $`\kappa =\kappa _\mathrm{B}=1.24`$ and $`\gamma =\gamma _\mathrm{B}=\pm 0.9`$. Giving equal weight to the two shear directions we can then determine the probability $`P_\mathrm{q}(M,r)`$ to obtain a quiet period lasting at least 8 years during a period of 15 years with a variability $`\delta m=m_{\mathrm{max}}m_{\mathrm{min}}`$ less than $`0.05`$ mag. Since $`t`$ is proportional to $`y\sqrt{M}`$, the probability $`P_\mathrm{q}(M,r)`$ obviously increases monotonically with $`M`$. Hence, we can determine a critical mass $`M_\mathrm{q}(r)`$ by requiring that $`P_\mathrm{q}(M_\mathrm{q},r)`$ is, say 0.01. We then choose to exclude values of $`M(r)`$ less than $`M_\mathrm{q}(r,P_\mathrm{q})`$ since the probability of getting a quiet period of 8 years or more is then less than one percent (compare the paper by Schmidt & Wambsganss, 1998). The dependence of $`M_\mathrm{q}`$ on $`r`$ for $`P_\mathrm{q}=0.01`$ is shown in the exclusion diagram in Fig. 2 as the solid line to the left labeled q.
We can roughly distinguish 3 different parts of the curve $`M_\mathrm{q}(r,0.01)`$:
1) $`r<0.5`$ : Here we find $`M_\mathrm{q}`$ very nearly constant and equal to about $`0.05M_{}`$. This corresponds to $`R_\mathrm{E}10^{16}`$ cm and $`\tau _\mathrm{E}=R_\mathrm{E}/V5`$ years.
The behavior with an almost constant $`M_\mathrm{q}`$ (and $`R_\mathrm{E}`$ and $`\tau _\mathrm{E}`$) for small values of $`r`$ is a consequence of the nearly identical light-curves for different (small) $`r`$ with the same star field. Only during caustic crossings will there be a major difference in the light-curves, but then the variability will be large anyway and therefore violate the quiet phase condition. For small $`r`$ the probability of getting a quiet phase therefore only depends on $`M`$ so that $`M_\mathrm{q}`$ is nearly independent of $`r`$ for $`r<0.5`$.
2) $`0.5<r<10`$ : Here we find that $`M_q`$ is roughly proportional to $`r^2`$. Since $`r=R/R_\mathrm{E}R/\sqrt{M}`$ this means that $`R`$ is roughly constant in this range. We find here $`R410^{15}`$ cm and $`\tau =R/V2`$ years. This result is not unexpected, since the variability time scale for sources with $`r>1`$ is about equal to $`2\tau `$, see RS2.
3) For $`r>10`$ the curve gets flatter and approaches asymptotically a maximum value of $`r_{\mathrm{qm}}40`$. Such an asymptotic behavior was to be expected from the results in RS1, where it was shown that microlensing amplitudes for large sources ($`r>>1`$) are proportional to $`r^1`$.
ii) Event constraints
In the simulated light-curves we consider only events with a variability between $`0.2`$ mag and $`0.3`$ mag ($`\mathrm{\Delta }m\pm 0.05`$), and with a maximum slope during the event steeper than 0.07 mag per year. The reason for choosing also an upper limit for the allowed event magnitudes is that the simulated events for small values of $`r`$ typically show larger amplitudes than the observed one. In our analysis we have found it more practical to use the maximum slope rather than the time scale $`\mathrm{\Delta }t`$. We can then estimate the probability $`P_\mathrm{e}(M,r)`$ of getting such an event during a time span of $`T=15`$ years. Since the slope scales as $`M^{0.5}`$ and the simulated time span as $`M^{0.5}`$, it is clear that $`P_e(M,r)`$ decreases monotonically with increasing $`M`$. We can therefore determine a critical mass $`M_\mathrm{e}(r)`$ by requiring that $`P_\mathrm{e}(M_\mathrm{e},r)`$ is, say $`0.01`$. We then choose to exclude values of $`M(r)`$ larger than $`M_\mathrm{e}(r,P_\mathrm{e})`$.
The dependence of $`M_\mathrm{e}`$ on $`r`$ for $`P_\mathrm{e}=0.01`$ is shown in Fig. 2 as the solid line to the right labeled e. It is seen that $`M_\mathrm{e}(r)`$ is roughly proportional to $`r^2`$, and therefore $`R`$ is fairly constant over a large range of $`r`$-values ($`0.2<r<10`$). We get here $`R10^{16}`$ cm and $`\tau =R/V=5`$ years. This is easy to understand since a typical event time scale is the crossing time $`2R/V`$ for large sources (see RS2) as well as for small sources (caustic crossings).
For small $`r`$-values we see that the $`M_\mathrm{e}`$-curve departs from $`R`$ constant towards smaller values of $`R`$, approaching a constant $`M_\mathrm{e}`$-value of about 5$`M_{}`$ for very small $`r`$. The reason is that for such small $`r`$-values, caustic crossing events as small as $`\mathrm{\Delta }m0.25`$ are quite rare. Of some importance here are, however, some non-caustic events, typically occurring when the track passes just outside a cusp of a caustic. For very small sources($`r<0.1`$) these events are nearly independent of $`r`$, which can explain that $`M_\mathrm{e}(r)`$ stays nearly constant for these small $`r`$-values. We note that the time scale for these events are usually longer than the crossing time $`2R/V`$, see also Wyithe et al. (1999).
For very small $`r`$-values the effect of the so-called ghost caustics which are produced when two lens masses are very close to each other is also of importance, (Grieger et al. 1989, Wambsganss & Kundić 1995). These weak caustics can often “produce” events during caustic crossings as small as about 0.25 mag, even for very small $`r`$-values. An uncertainty arises, however, because relative motion of the stars some times causes large velocities of the ghost caustics so that the “ghost events” will occur more often and on a shorter time scale. Since we have not considered relative motion in our calculations, we may have underestimated the frequency and overestimated the time scale of ghost events. In that case the upper limit for $`M_\mathrm{e}`$ has been underestimated.
For large $`r`$-values the $`M_\mathrm{e}`$-curve also departs from $`R`$ constant towards smaller values of $`R`$ and approaches asymptotically a maximum $`r`$-value of about $`r_{\mathrm{em}}=30`$. This high value of $`r_{\mathrm{em}}`$ were to be expected from the results in RS2, where it was shown that quite significant events can be produced for rather large values of $`r`$.
iii) The allowed region
From the discussion in section i) we found that the quiet phase constraint gives a lower limit, $`M_\mathrm{q}(r)`$ for the lens mass. In section ii) on the other hand, the event constraint gave an upper limit $`M_\mathrm{e}(r)`$. Hence, at a significance level of $`P`$ we are left with an allowed region between the $`M_\mathrm{e}`$ and $`M_\mathrm{q}`$-curves in the $`(M,r)`$-plane which fulfills
$$M_\mathrm{q}(r,P)<M(r)<M_\mathrm{e}(r,P)$$
(4)
For $`P=0.01`$ the allowed region is shown hatched in Fig. 2. For $`r`$ between 0.5 and 20 we find a strip within a relatively narrow range of $`R`$-values ($`410^{15}`$ cm$`<R<10^{16}`$ cm ), covering $`M`$-values between $`10^6M_{}`$ and $`10^1M_{}`$. For $`r<0.5`$, mass values between $`310^2M_{}`$ and $`1M_{}`$ are usually allowed. For the smallest $`r`$-values the upper limit reaches about 5$`M_{}`$. We cannot give a reliable lower limit of $`r`$. However, the lack of any clear high magnification events in all observed quasar light-curves is an argument against too small $`r`$-values (and $`R`$-values) for QSOs in general.
For larger $`P`$-values the allowed region obviously decreases. As an example the allowed region for $`P=0.1`$ is shown more densely hatched in Fig. 2. It is interesting to note that Eq. (4) can now only be fulfilled for $`r<0.6`$ and the possible mass shrinks approximately to the interval between 0.05 $`M_{}`$ and 0.5 $`M_{}`$ (significance level of 10%, and all mass in identical lenses).
## 4 Case 2: 90% of the lensing mass in a continuum
Since a large fraction of the lens mass may be in a smoothly distributed continuum, we have repeated most of the calculations in the previous section, but now with 90% of the mass in a continuum and the remaining 10% in compact lenses with mass $`M`$. The resulting constraint curves are shown as solid lines in Fig. 3 for a significance level of 0.01, and the allowed region is hatched. We find that the asymptotic maximum $`r`$-values $`r_{\mathrm{qm}}`$ and $`r_{\mathrm{em}}`$ decrease with a factor of about 3 relative to Case 1. This was to be expected because the variability amplitude for large $`r`$ is proportional to $`\sqrt{\kappa _\mathrm{s}}/r`$, so that $`r_{\mathrm{qm}}`$ and $`r_{\mathrm{em}}`$ should be proportional to $`\sqrt{\kappa _\mathrm{s}}`$, see RS1. Here $`\kappa _\mathrm{s}`$ is the optical depth for microlensing which is now a factor of 10 smaller than in Case 1.
The possible mass range is now restricted to an interval between $`310^5M_{}`$ and $`10^1M_{}`$.
For $`P=0.1`$ the allowed region is shown more densely hatched. We see that the possible mass now only covers an interval between $`210^3M_{}`$ and $`210^2M_{}`$, and $`r`$ is restricted to be less than about 2.
## 5 Case 3: Two lens populations
From the observed light and spectrum of the lensing galaxy we know that some of the mass must be in stars with typically one solar mass. We have therefore investigated a model with two lens populations; 10% of the mass in solar mass stars and the rest of the mass in objects of smaller mass, $`M`$. Treating as before $`M`$ and $`r=R/R_\mathrm{E}(M)`$ as free parameters we get mostly small changes (relative to Case 1) in the constraint curves in the ($`M,r`$)-plane. The $`M_\mathrm{q}`$-curves are almost unchanged, since the few $`1M_{}`$ stars can only make a minor change in the probability of getting long quiet phases for $`M<1M_{}`$. The $`M_\mathrm{e}`$-curves also follow very closely the corresponding curves for Case 1 (with all mass in identical stars). A separate discussion of this case is therefore not necessary.
## 6 Change of $`V`$ and $`H_0`$
We have until now assumed specific values of the effective transverse velocity $`V`$ and the Hubble parameter $`H_0`$ ($`v_{600}=1`$ and $`h_{60}=1`$). It is easy to see what happens if these values are changed. If for example the transverse velocity is increased by a factor $`F`$, an identical light-curve is obtained if the source size $`R`$ and the Einstein radius $`R_\mathrm{E}`$ are increased by the same factor, such that $`r=R/R_\mathrm{E}`$ is unchanged. Increasing the Einstein radius by a factor F means that each lens mass is increased by a factor $`F^2`$, and to keep the optical depth for lensing unchanged, the linear size of the lens field must of course be scaled up with a factor $`F`$.
Correspondingly, an increase of the Hubble parameter with a factor $`F`$ can be compensated by an increase of each lens mass with the same factor to keep the light-curve unchanged. Since the Einstein radius is unchanged in this case, the optical depth is also unchanged, and no change in the size of the lens field is necessary.
According to the above discussion, the allowed regions in Figs. 2 and 3 will be shifted in mass when $`V`$ and $`H_0`$ are changed.
$$M_\mathrm{q}(r)V^2H_0$$
(5)
$$M_\mathrm{e}(r)V^2H_0$$
(6)
This is most conveniently taken into account by just scaling the mass coordinate appropriately, thereby leaving the curves in Figs. 2 and 3 in the same positions. Obviously the dominating uncertainty in $`M`$ comes from the $`V^2`$ term.
## 7 Discussion
By means of the two features in the microlensing light-curve for the Q0957+561 system and numerical simulations, we have shown that it is possible to set constraints on some parameters for the lens system. With the amount of the continuum mass between 0% (Case 1) and 90% (Case 2) we find the possible lens masses to have values between $`10^6M_{}`$ and $`5M_{}`$ (1% level) and between $`210^3M_{}`$ and $`0.5M_{}`$ (10% level). The maximum source radius is found to be $`10^{16}`$ cm (1% level) and $`610^{15}`$ cm (10% level). By including Case 3 (giving approximately the same reults as Case 1), we find only negligible changes in the above limits.
From the two features in the microlensing light-curve it seems not possible to give a lower limit of $`R`$. The reason for this is that the light-curves for very small (but different) values of $`r`$ are very similar, except during caustic crossings, and then the quiet phase constraint is violated anyhow.
From lens arguments based on the Q0957+561 light-curve alone, we can therefore not rule out very small $`r`$-values (even those much smaller than 0.02), and hence very small $`R`$-values. However, the lack of any clear high magnification events in observed quasar light-curves is an argument against too small $`r`$-values (and $`R`$-values) for QSOs in general.
We have in this paper only considered microlensing effects in the B image. The combined A and B microlensing light-curve will statistically show slightly more variability and less chance of extended quiet phases, and have a slightly higher rate of events than for B alone. Hence, both $`M_\mathrm{q}`$ and $`M_\mathrm{e}`$ will increase slightly. Our test calculations indicate a rather small effect of about 10%.
Since we have neglected the relative motion of the lens masses we may have underestimated the maximum allowed lens mass which occur for small values of $`r`$, particularly in Case 1. This ought to be investigated in the future.
We have in this investigation considered the two features found in the light-curve for Q0957+561 independently and used the probabilities $`P_\mathrm{q}`$ and $`P_\mathrm{e}`$ separately. An alternative approach, which would have given much stronger constraints, might have been to consider the probability for both features to be found within a period of 15 years which must be less than $`P_\mathrm{q}P_\mathrm{e}`$. The largest values of $`P_\mathrm{q}P_\mathrm{e}`$ (between 0.01 and $`0.02)`$ are obviously found within the densely hatched regions in Fig. 2 (Case 1) and Fig. 3 (Case 2). The most probable mass values would then mainly depend on the amount of mass in the continuum ($`0.3M_{}`$ in Case 1 and $`0.01M_{}`$ in Case 2), and the most probable $`r`$-value is less than 1. It is, however, not possible to give levels of significance here since we are obviously doing a posteriori statistics. It seems in fact difficult to avoid a posteriori statistics when making use of a larger part of the information in the light-curve.
Wyithe et al. (1999) have published a series of papers (see astro-ph/9911245 and references therein) where they make use of the distribution of the microlensed light-curve derivatives, and they apply the method on the QSO 2237+0305 light-curve. Their method does not suffer from a posteriori statistics, but more information in the light-curve could perhaps have been taken into account, for instance the correlation time-scale for the derivatives, without running this risk. A comparison with their results is , however, difficult since they are considering a different system. For their system the value of $`R_\mathrm{E}`$ is about a factor of 3 larger, whereas the source radius $`R`$ is roughly a factor 2.5 smaller (estimated from the difference in intrinsic luminosity). Hence, the normalized source radius $`r=R/R_\mathrm{E}`$ will be a factor of about 7 smaller. With this factor in mind we do not see any serious conflict between the rather small $`R`$-values they obtain and our (most probable) values.
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# MEASUREMENTS OF 𝐹₂, 𝑥𝐹₃^𝜈-𝑥𝐹₃^𝜈̄ FROM CCFR 𝜈_𝜇-Fe and 𝜈̄_𝜇-Fe DATA IN A MODEL INDEPENDENT WAY
## 1 Introduction
Deep inelastic lepton-nucleon scattering experiments have been used to determine the quark distributions in the nucleon. However, the quark distributions determined from muon and neutrino experiments were found to be different at small values of $`x`$, because of a disagreement in the extracted structure functions. Here, we report on a measurement of differential cross sections and structure functions from CCFR $`\nu _\mu `$-Fe and $`\overline{\nu }_\mu `$-Fe data. We find that the neutrino-muon difference is resolved by extracting the $`\nu _\mu `$ structure functions in a physics model independent way.
The sum of $`\nu _\mu `$ and $`\overline{\nu }_\mu `$ differential cross sections for charged current interactions on an isoscalar target is related to the structure functions as follows:
| $`F(ϵ)`$ | $`\left[\frac{d^2\sigma ^\nu }{dxdy}+\frac{d^2\sigma ^{\overline{\nu }}}{dxdy}\right]\frac{(1ϵ)\pi }{y^2G_F^2ME_\nu }`$ $`=2xF_1[1+ϵR]+\frac{y(1y/2)}{1+(1y)^2}\mathrm{\Delta }xF_3`$. |
| --- | --- |
Here $`G_F`$ is the Fermi weak coupling constant, $`M`$ is the nucleon mass, $`E_\nu `$ is the incident energy, the scaling variable $`y=E_h/E_\nu `$ is the fractional energy transferred to the hadronic vertex, $`E_h`$ is the final state hadronic energy, and $`ϵ2(1y)/(1+(1y)^2)`$ is the polarization of the virtual $`W`$ boson. The structure function $`2xF_1`$ is expressed in terms of $`F_2`$ by $`2xF_1(x,Q^2)=F_2(x,Q^2)\times \frac{1+4M^2x^2/Q^2}{1+R(x,Q^2)}`$, where $`Q^2`$ is the square of the four-momentum transfer to the nucleon, $`x=Q^2/2ME_h`$ (the Bjorken scaling variable) is the fractional momentum carried by the struck quark, and $`R=\frac{\sigma _L}{\sigma _T}`$ is the ratio of the cross-sections of longitudinally- to transversely-polarized $`W`$-bosons. The $`\mathrm{\Delta }xF_3`$ term, which in leading order $`4x(sc)`$, is not present in the $`\mu `$-scattering case. In addition, in a $`\nu _\mu `$ charged current interaction with $`s`$ (or $`\overline{c}`$) quarks, there is a threshold suppression originating from the production of heavy $`c`$ quarks in the final state. For $`\mu `$-scattering, there is no suppression for scattering from $`s`$ quarks, but more suppression when scattering from $`c`$ quarks since there are two heavy quarks ($`c`$ and $`\overline{c}`$) in the final state.
In previous analyses of $`\nu _\mu `$ data, light-flavor universal physics model dependent (PMD) structure functions were extracted by applying a slow rescaling correction to correct for the charm mass suppression in the final state. In addition, the $`\mathrm{\Delta }xF_3`$ term (used as input in the extraction) was calculated from a leading order charm production model. These resulted in a physics model dependent (PMD) structure functions. In the new analysis reported here, slow rescaling corrections are not applied, and $`\mathrm{\Delta }xF_3`$ and $`F_2`$ are extracted from two parameter fits to the data. We compare the values of $`\mathrm{\Delta }xF_3`$ to various charm production models. The extracted physics model independent (PMI) values for $`F_2^\nu `$ are then compared with $`F_2^\mu `$ within the framework of NLO models for massive charm production.
## 2 Results
The CCFR experiment collected data using the Fermilab Tevatron Quad-Triplet wide-band $`\nu _\mu `$ and $`\overline{\nu }_\mu `$ beam. The raw differential cross sections per nucleon on iron are determined in bins of $`x`$, $`y`$, and $`E_\nu `$ ($`0.01<x<0.65`$, $`0.05<y<0.95`$, and $`30<E_\nu <360`$. GeV). Figure 1 (a) shows typical differential cross sections at $`E_\nu =150`$ GeV. Next, the raw cross sections are corrected for electroweak radiative effects, the $`W`$ boson propagator, and for the 5.67% non-isoscalar excess of neutrons over protons in iron (only important at high $`x`$). Values of $`\mathrm{\Delta }xF_3`$ and $`F_2`$ are extracted from the sums of the corrected $`\nu _\mu `$-Fe and $`\overline{\nu }_\mu `$-Fe differential cross sections at different energy bins according to Eq. (1). It is challenging to fit $`\mathrm{\Delta }xF_3`$, $`R`$, and $`2xF_1`$ using the $`y`$ distribution at a given $`x`$ and $`Q^2`$ because of the strong correlation between the $`\mathrm{\Delta }xF_3`$ and $`R`$ terms, unless the full range of $`y`$ is covered by the data. Covering this range (especially the high $`y`$ region) is hard because of the low acceptance. Therefore, we restrict the analysis to two parameter fits.
Our strategy is to fit $`\mathrm{\Delta }xF_3`$ and $`2xF_1`$ (or equivalently $`F_2`$) for $`x<0.1`$ where the $`\mathrm{\Delta }xF_3`$ contribution is relatively large, while constraining $`R`$ using the $`R_{world}^{\mu /e}`$ QCD inspired empirical fit to all available $`R`$ from electron- and $`\mu `$-scattering data. The $`R_{world}^{\mu /e}`$ fit is also in good agreement with NMC $`R^\mu `$ data at low $`x`$, and with the most recent NNLO QCD calculations (including target mass effects) of $`R`$ by Bodek and Yang
For $`x<0.1`$, $`R`$ in neutrino scattering is expected to be somewhat larger than $`R`$ for muon scattering because of the production of massive charm quarks in the final state. A correction for this difference is applied to $`R_{world}^{\mu /e}`$ using a leading order slow rescaling model to obtain an effective $`R`$ for neutrino scattering, $`R_{eff}^\nu `$. The difference between $`R_{world}^{\mu /e}`$ and $`R_{eff}^\nu `$ is used as a systematic error. Because of the positive correlation between $`R`$ and $`\mathrm{\Delta }xF_3`$, the extracted values of $`F_2`$ are rather insensitive to the input $`R`$. If a large input $`R`$ is used, a larger value of $`xF_3`$ is extracted from the $`y`$ distribution, thus yielding the same value of $`F_2`$. In contrast, the extracted values of $`\mathrm{\Delta }xF_3`$ are sensitive to the assumed value of $`R`$, which is reflected in a larger systematic error. The values of $`\mathrm{\Delta }xF_3`$ are sensitive to the energy dependence of the neutrino flux ($``$ $`y`$ dependence), but are insensitive to the absolute normalization. The uncertainty on the flux shape is estimated by using the constraint that $`F_2`$ and $`xF_3`$ should be flat over $`y`$ (or $`E_\nu `$) for each $`x`$ and $`Q^2`$ bin.
Because of the limited statistics, we use large bins in $`Q^2`$ in the extraction of $`\mathrm{\Delta }xF_3`$ with bin centering corrections from the NLO Thorne & Roberts Variable Flavor Scheme (TR-VFS) calculation with the MRST PDFs. Figure 1 (b) shows the extracted values of $`\mathrm{\Delta }xF_3`$ as a function of $`x`$, including both statistical and systematic errors, compared to various theoretical methods for modeling heavy charm productions within a QCD framework. The three-flavor Fixed Flavor Scheme (FFS) assumes that there is no intrinsic charm in the nucleon, and all scattering from $`c`$ quarks occurs via the gluon-fusion diagram. The concept behind the Variable Flavor Scheme (VFS) proposed by ACOT is that at low scale, $`\mu `$, one uses the three-flavor FFS scheme, and above some scale, one changes to a four-flavor calculation and an intrinsic charm sea (which is evolved from zero) is introduced. The concept in the RT-VFS scheme is that it starts with the three-flavor FFS scheme at a low scale, becomes the four-flavor VFS scheme at high scale, and interpolates smoothly between the two regions. Shown are the predictions from the TR-VFS scheme (as corrected after DIS-2000 and implemented with MRST PDFs), with their suggested scale $`\mu =Q`$, and the predictions of the other two NLO calculations, ACOT-VFS (implemented with CTEQ4HQ and the recent ACOT suggested scale $`\mu =m_c`$ for $`Q<m_c`$, and $`\mu ^2=m_c{}_{}{}^{2}+cQ^2(1m_c{}_{}{}^{2}/Q^2)^n`$ for $`Q<m_c`$ with $`c=0.5`$ and $`n=2`$), and the FFS (implemented with the GRV94 PDFs and GRV94 recommended scale $`\mu =2m_c`$). Also shown are the predictions from $`\mathrm{\Delta }xF_34Ks(x,Q^2)`$ from a leading order model (LO(CCFR)) Buras-Gaemers type fit to the CCFR dimuon data (here K is a slow rescaling correction). Figure 1 (b) (right) also shows the sensitivity to the choice of scale. The data do not favor the ACOT-VFS(CTEQ4HQ) predictions if implemented with an earlier suggested scale of $`\mu =2Pt_{max}`$. With reasonable choices of scale, all the theoretical models yield similar results. However, at low $`Q^2`$ our $`\mathrm{\Delta }xF_3`$ data are higher than all the theortical models. The difference between data and theory may be due to an underestimate of the strange sea (or gluon distribution) at low $`Q^2`$, or from missing NNLO terms.
As discussed above, values of $`F_2`$ (PMI) for $`x<0.1`$ are extracted from two parameter fits to the $`y`$ distributions. In the $`x>0.1`$ region, the contribution from $`\mathrm{\Delta }xF_3`$ is small and the extracted values of $`F_2`$ are insensitive to $`\mathrm{\Delta }xF_3`$. Therefore, we extract values of $`F_2`$ with an input value of $`R`$ and with $`\mathrm{\Delta }xF_3`$ constrained to the TR-VFS(MRST) predictions. As in the case of the two parameter fits for $`x<0.1`$, no corrections for slow rescaling are applied. Fig. 2 (a) shows our $`F_2`$ (PMI) measurements divided by the predictions from the TR-VFS(MRST) theory. Also shown are $`F_2^\mu `$ and $`F_2^e`$ from the NMC divided by the theory predictions. In the calculation of the QCD TR-VFS(MRST) predictions, we have also included corrections for nuclear effects, target mass and higher twist corrections at low values of $`Q^2`$. As seen in Fig. 2, both the CCFR and NMC structure functions are in good agreement with the TR-VFS(MRST) predictions, and therefore in good agreement with each other. A comparison using the ACOT-VFS(CTEQ4HQ) predictions yields similar results.
In the previous analysis of the CCFR data, the extracted values of $`F_2`$ (PMD) at the lowest $`x=0.015`$ and $`Q^2`$ bin were up to $`20\%`$ higher than both the NMC data and the predictions of the light-flavor MRSR2 PDFs. (see figure 2 (b) ). About half of the difference originates from having used a leading order model for $`\mathrm{\Delta }xF_3`$ versus using our new measurement. The other half originates from having used the leading order slow rescaling corrections, instead of using a NLO massive charm production model, and from improved modeling of the low $`Q^2`$ PDFs (which changes the radiative corrections and the overall absolute normalization to the total neutrino cross sections).
## 3 Conclusions
In conclusion, the $`F_2`$ (PMI) values measured in neutrino-iron and muon-deuterium scattering show good agreement with with the predictions of Next to Leading Order PDFs (using massive charm production schemes), thus resolving the long-standing discrepancy between the two sets of data. The first measurements of $`\mathrm{\Delta }xF_3`$ are higher than current theoretical predictions.
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# Role of Coulomb correlations in the optical spectra of semiconductor quantum dots
## Appendix
1. Matrix elements. With $`\varphi ^{e,h}`$ the single-particle states for electrons and holes, the optical matrix elements are of the form $`M_{\nu \mu }^{he}=\mu _o𝑑𝐫\varphi _\nu ^h(𝐫)\varphi _\mu ^e(𝐫)`$, with $`\mu _o`$ the dipole matrix element of the bulk semiconductor. The Coulomb matrix elements read:
$$V_{\mu ^{}\mu ,\nu ^{}\nu }^{ij}=q_iq_j𝑑𝐫𝑑𝐫^{}\frac{\varphi _{\mu ^{}}^{i}{}_{}{}^{}(𝐫)\varphi _{\nu ^{}}^{j}{}_{}{}^{}(𝐫^{})\varphi _\nu ^j(𝐫^{})\varphi _\mu ^i(𝐫)}{\kappa _o|𝐫𝐫^{}|},$$
(3)
with $`\kappa _o`$ the static dielectric constant of the semiconductor, $`i,j=e,h`$ and $`q_{e,h}=1`$ (note that we have neglected electron-hole exchange interactions).
2. Few-particle states. We compute the few-particle electron-hole states within a direct-diagonalization approach. With the creation operators $`c^{}`$ and $`d^{}`$ for electrons and holes, respectively, we define the $`N_e`$-electron and $`N_h`$-hole states $`|\stackrel{}{\mu }_{N_e}=c_{\mu _1}^{}c_{\mu _2}^{}\mathrm{}c_{\mu _{N_e}}^{}|\mathrm{\Phi }_o`$ and $`|\stackrel{}{\nu }_{N_h}=d_{\nu _1}^{}d_{\nu _2}^{}\mathrm{}d_{\mu _{N_h}}^{}|\mathrm{\Phi }_o`$ (vacuum state $`|\mathrm{\Phi }_o`$; spin degrees of freedom have not been indicated explicitly), and we keep the $``$100 few-particle states of lowest single-particle energies. Next, the few-particle Hamiltonian $``$, accounting for all possible electron-electron, electron-hole, and hole-hole Coulomb matrix elements, is expanded within these bases; the few-particle energies $`E_{\mathrm{}}`$ and wavefunctions $`\mathrm{\Psi }_{\stackrel{}{\mu }\stackrel{}{\nu }}^{\mathrm{}}`$ are then obtained through direct diagonalization of the Hamiltonian matrix:
$$E_{\mathrm{}}\mathrm{\Psi }_{\stackrel{}{\mu };\stackrel{}{\nu }}^{\mathrm{}}=\underset{\stackrel{}{\mu }^{},\stackrel{}{\nu }^{}}{}{}_{N_e;N_h}{}^{}\stackrel{}{\mu };\stackrel{}{\nu }||\stackrel{}{\mu }^{};\stackrel{}{\nu }^{}_{N_e;N_h}^{}\mathrm{\Psi }_{\stackrel{}{\mu }^{};\stackrel{}{\nu }^{}}^{\mathrm{}}.$$
(4)
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# POSTSTARBURST MODELS OF LINERS
## 1. INTRODUCTION
For the past two decades, LINERs have been recognized as one of important classes of AGNs (Heckman 1978; Stauffer 1982; Keel 1983a, 1983b Filippenko & Sargent 1985; Véron-Cetty & Véron 1986; Phillips et al. 1986a, 1986b; Maiolino & Rieke 1995; Ho, Filippenko, & Sargent 1997a, 1997c). They are defined by the following criteria based on optical emission-line intensity ratios (Heckman 1978); 1) $`I`$(\[O ii\]$`\lambda `$3727)/$`I`$(\[O iii\]$`\lambda `$5007) $`1`$ where \[O ii\]$`\lambda `$3727 is used to designate the \[O ii\]$`\lambda \lambda `$3726, 3729 doublet, and 2) $`I`$(\[O i\]$`\lambda `$6300)/$`I`$(\[O iii\]$`\lambda `$5007) $`1/3`$. The most recent and comprehensive survey for activity in galactic nuclei conducted by Ho et al. (1997a) has shown that LINERs make up a significant fraction ($``$ 20 – 30 %) of all galaxies and hence they constitute the most populous class of AGNs; i.e., $``$ 50 – 75 % of AGNs are LINERs. Therefore, it is very important to understand the origin of LINERs.
Possible origins of LINERs discussed previously are as follows: a) Low-ionization analog of typical AGNs (Ferland & Netzer 1983; Halpern & Steiner 1983; Filippenko & Halpern 1984; Ho, Filippenko, & Sargent 1993; 1997d), b) powering by shocks driven either by superwinds or by nuclear radio jets (Daltabuit & Cox 1972; Koski & Osterbrock 1976; Fosbury et al. 1978; Heckman 1978; Dopita & Sutherland 1995, 1996; Alonso-Herrero et al. 2000; Sugai & Malkan 2000), c) photoionization by Wolf-Rayet stars (Terlevich & Melnick 1985), d) photoionization by hot O stars (Filippenko & Terlevich 1992; Shields 1992), and e) photoionization by old post-AGB stars in elliptical galaxies (Binette et al. 1994). It is noted that the underlying Balmer absorption due to many A-type stars weakens Balmer emission and thus results in apparently higher \[N ii\]$`\lambda `$6583/H$`\alpha `$ intensity ratios in some cases (e.g., Taniguchi et al. 1996); note that this ratio is often used to distinguish LINERs from H ii regions; i.e., $`I`$(\[N ii\]$`\lambda `$6583)/$`I`$(H$`\alpha `$) $`0.6`$ for LINERs while $`<0.6`$ for H ii nuclei (Veilleux & Osterbrock 1987; Ho et al. 1997a). However, the overall statistics of the incidence of LINERs (e.g., Ho et al. 1997a) are probably not affected by Balmer absorption because careful starlight subtraction has been made to remove the effects of Balmer absorption from the data, by matching the stellar populations of each LINER with a template absorption-line galaxy.
### 1.1. Basic concept
There are several lines of evidence for the presence of pure AGN in some LINERs: i) the presence of broad emission-line region (BLR) (Ho et al. 1997d; Maoz et al. 1998), ii) the presence of hidden BLR in polarized optical spectra (Wilkes et al. 1995; Barth, Filippenko,& Moran 1999), iii) a correlation between line width and critical density of the forbidden lines arising from narrow emission-line region (NLR) (Filippenko & Halpern 1984), iv) the presence of point-like ultraviolet source<sup>1</sup><sup>1</sup>1 Since UV spectra have shown that some of the point-like UV sources are in fact young starburst clusters rather than AGNs (Maoz et al. 1998), the presence of an unresolved UV source may not always indicate the presence of an AGN. (Maoz et al. 1995), v) the presence of radio core and/or jets (e.g., Heckman 1980; Wrobel 1984; Slee et al. 1994), and vi) the detection of hard X-ray continuum as well as Fe K emission (Iyomoto et al. 1996; Ishisaki et al. 1996; Terashima et al. 1998a, 1998b, 2000; Roberts, Warwick, & Ohashi 1999). Despite these firm lines of evidence for the presence of pure AGN in a non-negligibly large number of LINERs (e.g., $``$ 20% of LINERs have the BLR: Ho et al. 1997d), the majority of LINERs have no firm evidence for the presence of pure AGN, suggesting that LINERs are heterogeneous (Heckman 1986).
In this paper, we investigate whether or not there are plausible poststarburst models for the ionization mechanism of LINERs taking account that intermediate-mass stars ($``$ several $`M_{}`$) are also formed in nuclear starbursts (Joseph 1991). Binette et al. (1994) showed that old post-AGB stars provide sufficient ionizing photons to account for the observed H$`\alpha `$ luminosity of LINERs associated with elliptical galaxy nuclei. In their model, it was assumed that post-AGB stars were formed in the early phase of the galaxy evolution (i.e., the initial starburst). On the other hand, in our paper, our interest is addressed to starbursts occurred in recent pasts in nuclei of spiral galaxies. Therefore, our model presented in this paper is different from theirs.
## 2. STARBURST AND POSTSTARBURST EVOLUTION
In order to investigate the evolution of a star cluster formed in a starburst, it is necessary to know both the initial mass function (IMF) and the star formation rate (SFR). The SFR for nuclear starbursts can be estimated using observed H$`\alpha `$ luminosities and so on (Kennicutt 1998). Although it is generally difficult to derive the IMF accurately, it seems reasonable to adopt a power-law form of IMF; $`\varphi (m)m^\mu `$ (e.g., Scalo 1986). In this formulation, there are three free parameters; the power index ($`\mu `$), and the upper and lower mass limits of the IMF ($`m_u`$ and $`m_l`$). The following values are often adopted for the evolution of the solar neighborhood; $`\mu `$ = 1.35, $`m_l=0.1M_{}`$, and $`m_u=60M_{}`$. On the other hand, summarizing various kinds of observational constraints, Joseph (1991) suggests that $`m_u`$ 30 – 60 $`M_{}`$ and $`m_l`$ 3 – 6 $`M_{}`$ for typical nuclear starbursts. Although the top-heavy IMF is suggested for the nuclear starbursts (i.e., $`\mu `$ is smaller than 1.35; e.g., Scalo 1990), we adopt $`\mu `$ =1.35, $`m_u=60M_{}`$, and $`m_l=3M_{}`$ as a modest combination. Another parameter is the duration of the starburst, $`\tau _{\mathrm{SB}}`$; we adopt $`\tau _{\mathrm{SB}}=10^7`$ yr because the negative feedback from supernova explosions to star-forming gas clouds is expected to occur $`10^7`$ yr (i.e., the lifetime of B stars) after the onset of the starburst (e.g., Larson 1987). The lifetime of a starburst can also be estimated as a gas consumption timescale, $`\tau _{\mathrm{SB}}M_{\mathrm{gas}}\eta _{\mathrm{SF}}/SFR`$ where $`M_{\mathrm{gas}}`$ is the gas mass available for the starburst, $`\eta _{\mathrm{SF}}`$ is the star formation efficiency, and $`SFR`$ is the star formation rate. If we adopt $`M_{\mathrm{gas}}10^9M_{}`$, $`\eta _{\mathrm{SF}}0.1`$, and $`SFR10M_{}`$ yr<sup>-1</sup>, we obtain $`\tau _{\mathrm{SB}}10^7`$ y. We will discuss the SFR later.
Now let us consider the evolution of the star cluster with the above parameters. Phase I ($`0t10^7`$ yr): We assume that the starburst lasts for $`10^7`$ yr. The star formation rate is assumed to be constant. In this phase, main photoionization sources are most massive stars. Dynamical effect of supernova explosions may be weak in this phase. Phase II ($`10^7t2\times 10^7`$ yr): Since the star formation ceases at a time $`t=10^7`$ yr in our model, main photoionization sources change from O stars to B stars as time goes. Wolf-Rayet stars also work in the photoionization (Vacca & Conti 1992; Conti 1999). Phase III ($`2\times 10^7t1\times 10^8`$ yr): Continuous supernova explosions develop a superwind and then shock heating also works in this phase (Heckman, Armus, & Miley 1990; Ohyama, Taniguchi, & Terlevich 1997; Alonso-Herrero et al. 2000). Phase IV \[$`t(15)\times 10^8`$ yr\]: Intermediate-mass stars with mass of several $`M_{}`$ play a role as the photoionization sources. The bolometric luminosity of each star in the main sequence phase is $`L_{}1\times 10^3L_{}`$. Each star evolves from the main sequence to the asymptotic giant branch (AGB). Following this AGB phase, each star will lose its gaseous envelope and then a hot stellar core appears and ionizes the surrounding gas, making a so-called planetary nebula (Kaler 1985; Vassiliadis, & Wood 1994). Since this core (i.e., a planetary-nebula nucleus; hereafter PNN) is so hot (i.e., $`T_{\mathrm{eff}}10^5`$ K), they become to be major photoionization sources. Finally, stars with mass of $`3M_{}`$ die at $`t5\times 10^8`$ yr<sup>2</sup><sup>2</sup>2Note that the ages of intermediate-mass stars are estimated as $`5.7\times 10^8`$ yr, $`2.5\times 10^8`$ yr, and $`1.4\times 10^8`$ yr for stellar masses of $`3M_{}`$, $`4M_{}`$, and $`5M_{}`$, respectively using a relation of $`\tau (m)=1.2\times 10^{10}m^{2.78}`$ yr for $`m<10M_{}`$ (Theis, Burkert, & Hensler 1992). and then the effect of the nuclear starburst disappears unless stars less massive than 3 $`M_{}`$ were formed in the starburst. The starburst evolution is summarized schematically in Figure 1.
In many previous studies, Phase IV has not been taken into account seriously. However, the temperature of PNNs is hot (i.e., $`T_{\mathrm{eff}}10^5`$ K) enough to ionize the surrounding nebula. More importantly, this higher temperature can be responsible for the formation of partly-ionized regions in which \[O i\] emission is thought to arise. If the ionization parameter is as low as that for typical LINERs, the optical spectrum of such a nebula is expected to be quite similar to those of LINERs. Note that each PNN has a bolometric luminosity of $`L_{\mathrm{PNN}}10^4L_{}`$ (e.g., Vassiliadis, & Wood 1994), being comparable to that of a hot O star (Filippenko, & Terlevich 1992; Shields 1992).
### 2.1. Poststarburst evolution model
In order to demonstrate the importance of the PNN cluster, we investigate the poststarburst evolution for the star cluster discussed in section 2.1 using luminosity evolution models of Kodama & Arimoto (1997). In Figure 2, we show the time evolution of H$`\alpha `$ luminosity which is normalized by a unit gas mass, $`1M_{}`$. Here we assume that stars are formed from the gas with the solar metallicity for simplicity. The H$`\alpha `$ luminosity is estimated from the number of Lyman continuum photons; $`L(\mathrm{H}\alpha )1.36\times 10^{12}N(\mathrm{Lyc})`$ ergs s<sup>-1</sup> (Leitherer & Heckman 1995). Note that the lower mass cutoff is fixed at 0.1 $`M_{}`$ in Kodama & Arimoto’s models. Therefore, the H$`\alpha `$ luminosity shown in Figure 2 is corrected for the case of $`m_l=3M_{}`$ by us.
According to Kennicutt (1998), the SFR is related to $`L(\mathrm{H}\alpha )`$ as $`SFR=7.9\times 10^{42}L(\mathrm{H}\alpha )M_{}\mathrm{yr}^1`$ for the IMF with $`\mu =1.35`$, $`m_l=0.1M_{}`$, and $`m_u=100M_{}`$. This relation can be replaced by $`SFR=4.0\times 10^{42}L(\mathrm{H}\alpha )M_{}\mathrm{yr}^1`$ for the Salpeter IMF with $`m_l=3M_{}`$ and $`m_u=60M_{}`$. Balzano (1983) found $`L(\mathrm{H}\alpha )10^{40}`$$`6\times 10^{42}`$ ergs s<sup>-1</sup> for her sample of Markarian starburst nuclei with a Hubble constant of $`H_0`$ = 75 km s<sup>-1</sup> Mpc<sup>-1</sup> (see also Kennicutt, Keel, & Blaha 1989; Ho, Filippenko, & Sargent 1997b). Since the duration of H$`\alpha `$-bright phase is very short (e.g., $`10^7`$ yr; see Figure 2) in the starbursts, it seems rarer to detect starbursts with such a bright phase from a statistical point of view. For example, the age of the starburst occurring in Mrk 1259, which is one of the brightest starburst nuclei studied by Balzano (1983), is estimated to be $`5\times 10^6`$ yr (Ohyama et al. 1997). Although the observed H$`\alpha `$ luminosity of this galaxy is $`10^{41}`$ ergs s<sup>-1</sup>, this galaxy would be observed as a less luminous starburst if it will be observed $`10^{78}`$ yr after now. Typical H$`\alpha `$ luminosities listed in Balzano’s catalog are $`10^{40}`$ \- $`10^{41}`$ ergs s<sup>-1</sup>. If the above statistical effect is taken into account, their initial H$`\alpha `$ luminosities may be of the order of $`10^{41}`$ \- $`10^{42}`$ ergs s<sup>-1</sup>. Furthermore, since the mean reddening correction is about a factor of 4-5 for H$`\alpha `$ (e.g., Balzano 1983), the reddening-corrected, initial H$`\alpha `$ luminosities seem to be as luminous as $`10^{42}`$ ergs s<sup>-1</sup> on average and some brightest nuclear starbursts may have H$`\alpha `$ luminosities of $`10^{43}`$ ergs s<sup>-1</sup>. These values lead to $`SFR`$ 10 – 100 $`M_{}`$ yr<sup>-1</sup>. Then, we estimate $`M_{\mathrm{gas}}10^{89}M_{}`$ for the nuclear starbursts. Indeed, the molecular gas mass in circumnuclear regions of typical starburst nuclei is $`M_{\mathrm{gas}}10^9M_{}`$ (Devereux et al. 1994). As shown in Figure 2, the H$`\alpha `$ luminosity density per unit mass at $`t10^8`$ yr is $`L(\mathrm{H}\alpha )/M_{\mathrm{gas}}10^{30}`$ ergs s<sup>-1</sup> $`M_{}^1`$. Therefore, we estimate typical H$`\alpha `$ luminosities of poststarburst nuclei with $`t10^8`$ yr as $`L(\mathrm{H}\alpha )10^{38}`$ ergs s<sup>-1</sup> for a typical LINER powered by PNNs, and $`10^{39}`$ ergs s<sup>-1</sup> in exceptionally bright cases.
It should be mentioned that our poststarburst models cannot be applied to on-going starburst galaxies because most starburst nuclei have $`L(\mathrm{H}\alpha )10^{39}`$ ergs s<sup>-1</sup> (Balzano 1983; Kennicutt et al. 1989; Ho et al. 1997b). Namely, if a nuclear starburst occurs, the photoionization should be dominated by massive stars in the starburst even if there is a cluster of PNNs left from the recent past nuclear starburst.
### 2.2. Optical emission-line properties of the ionized nebula
Next we investigate optical emission-line properties of the gaseous nebula photoionized by the PNN cluster. We use the photoionization code CLOUDY version 90.05 (Ferland 1996), which solves the equations of statistical and thermal equilibrium and produces a self-consistent model of the run of temperature as a function of depth into the nebula. Here we assume that a uniform-density gas cloud with plane-parallel geometry is irradiated by the star cluster formed in the starburst. The parameters for the calculations are 1) the hydrogen density of the cloud ($`n_\mathrm{H}`$), 2) the ionization parameter (Osterbrock 1989); $`U=Q(\mathrm{H}^0)(4\pi r^2n_\mathrm{H}c)^1`$ where $`Q(\mathrm{H}^0)`$ is the number of ionizing photons, $`r`$ is the distance from the ionizing source, and $`c`$ is the light velocity, 3) the spectral energy distribution (SED) of the ionizing radiation, and 4) the chemical compositions.
The SED of the ionizing radiation is calculated during the course of the starburst evolution using Kodama & Arimoto’s models. The SED evolution is shown in Figure 3. Note that the ordinate is the flux in units of $`L_{}`$ Å<sup>-1</sup>, corresponding to the case that the total gas mass is $`1M_{}`$. We perform our photoionization calculations as a function of the ionization parameter between log $`U=4`$ and log $`U=3`$ with a logarithmic interval of 0.2. Following Filippenko & Terlevich (1992), we adopt $`n_\mathrm{H}=10^3`$ cm<sup>-3</sup>. In order to examine the effect of dust grains on emission-line intensity ratios, we perform photoionization calculations for the following two cases; a) the dust-free gas with the solar abundances, and b) the gas with the Orion Nebula abundances where effects of dust grains are taken into account. For both the cases we simply adopted the abundance sets available in CLOUDY. The solar composition is taken from Grevesse & Anders (1989) and Grevesse & Noels (1993). The Orion Nebula abundances are a subjective mean of the abundances determined by Baldwin et al. (1991), Rubin et al. (1991), and Osterbrock, Tran, & Veilleux (1992). The grains are the large-R grains described by Baldwin et al. (1991). The calculations were stopped when the temperature fell to 3000 K, below which little optical emission is expected.
The results are shown in Figures 4 and 5 for the above two cases, respectively. Here we use the same emission-line diagnostic diagrams as those used in Filippenko & Terlevich (1992). It is shown that our poststarburst models with age of $`t(1`$$`5)\times 10^8`$ yr appear consistent with the observed emission-line properties of LINERs; the depletion of metals into dust grains leads to higher relative intensities of low-ionization lines such as \[O i\] and \[S ii\], being more consistent with the observations. However, it should be noted that the \[O i\]/H$`\alpha `$ ratio barely touches the LINER region of the diagnostic diagram in both Figure 4 and 5. In this respect, our models are similar to the previous O-star models for weak \[O i\] LINERs proposed by Filippenko & Terlevich (1992) and Shields (1992); i.e., the weak \[O i\] LINERs are defined as LINERs with $`I`$(\[O i\]$`\lambda `$6300)/$`I`$(H$`\alpha `$) $`1/6`$. However, since PNNs are hot enough to supply such high-energy photons, the other emission-line ratios appear more consistent with the observations (see Figures 4 and 5) than those of the O-star models. The weakness of \[O i\] emission in our models may be attributed to the fact that the production rate of high-energy photons is lower than that of the non-thermal continuum radiation from an AGN.
In order to demonstrate this quantitatively, we show the time evolution of the number of ionizing photons which are capable of ionizing He<sup>0</sup> (i.e., $`h\nu `$ 24.6 eV) in Figure 2 (see the dotted line). Although high-energy photons are supplied by O stars in Phase I, they decrease in number significantly in Phase III. However, they increase in number again in Phase IV because of the supply from high-temperature PNNs. This makes our models more consistent with the observations than the previous O-star photoionization models.
It is known that intermediate-mass stars produce strong Balmer absorption lines in their stellar atmospheres and thus the Balmer emission lines arising from the ionized nebulae are sometimes quenched by the Balmer absorption. In order to investigate this effect quantitatively along the starburst evolution, we show time variations of the equivalent width ($`EW`$) of both Balmer emission and absorption for the case of H$`\beta `$ in Figure 6. The H$`\beta `$ emission dominates in the first a few times $`10^7`$ yr while the H$`\beta `$ absorption does after this age. This comparison suggests that LINERs whose Balmer lines are observed as absorption are more popular by about one order of magnitude than those with the Balmer emission. On the other hand, observations show that the fraction of the former LINERs is significantly smaller than that of the latter ones (Ho et al. 1997a). However, if the Balmer emission has a narrower line width on average than the Balmer absorption, the Balmer emission can be seen even if the $`EW_{\mathrm{absorption}}>EW_{\mathrm{emission}}`$. This trend is actually seen in the observed spectra of many LINERs (Ho et al. 1997a).
## 3. CONFRONTATION WITH OBSERVATION
### 3.1. Frequency of occurrence of poststarburst LINERs
We estimate the frequency of occurrence of poststarburst LINERs. Since LINERs are very common (Ho et al. 1997a), any models are required to explain this property. The lifetime of the poststarburst LINER would be as long as $`10^{10}`$ yr if $`m_l`$ is less than 1 $`M_{}`$ (Binette et al. 1994). However, this case is beyond our scope and should be applied to the nuclear emission-line regions of long-lived elliptical galaxies. Here we adopt that the lifetime of LINER in our poststarburst model is $`\tau (\mathrm{LINER})5\times 10^8`$ yr given that the lowest mass of stars formed in the starburst is 3 $`M_{}`$. On the other hand, the lifetime of the starburst is $`\tau ^{\mathrm{eff}}(\mathrm{SB})2\times 10^7`$ yr where $`\tau ^{\mathrm{eff}}(\mathrm{SB})`$ is the sum total of the durations of both Phase I and II. Therefore, the expected number ratio between starbursts and poststarburst LINERs is $`N(\mathrm{SB})/N(\mathrm{LINER})\tau ^{\mathrm{eff}}(\mathrm{SB})/\tau (\mathrm{LINER})1/25`$. Since typical nuclear-starburst galaxies are observed in $``$ 1% of nearby galaxies (e.g., Balzano 1983), poststarburst LINERs are expected to share $``$ 25% of nearby galaxies, being consistent with the observation (Ho et al. 1997a). However, we should mention why nuclear starbursts with $`\tau ^{\mathrm{eff}}(\mathrm{SB})2\times 10^7`$ yr are observed in $`1`$% of nearby galaxies because the comparison between $`\tau ^{\mathrm{eff}}(\mathrm{SB})`$ and the Hubble time, $`10^{10}`$ yr, would give statistically a frequency of occurrence of starbursts of $``$ 0.2%. The difference between this estimate and the actual frequency may be reconciled if every galaxy experiences starbursts several times in its life.
We mention that the above estimate is based on an assumption that there has been no star formation activity since the recent past starburst. If the galaxy currently experiences nuclear star formation at a level of even $`10^3`$ times the star formation rate of the starburst, then the young stellar population would dominate over PNNs by an order of magnitude, and the emission-line spectrum would be that of an H ii region, not a LINER. If there is any current or very recent star formation, the frequency of LINERs generated by this poststarburst mechanism may be significantly lower than the above estimate. However, it has been argued that the negative feedback from supernova explosions to star-forming gas clouds prevent from new intense star formation for $`10^8`$ yr after the onset of the starburst (e.g., Larson 1987). Therefore, our assumption seems reasonable.
### 3.2. Morphological-type segregation
First of all, we mention that our models cannot be applied to LINERs associated with the nuclei of elliptical galaxies and most of S0 in which nuclear starburst occurs seldom although some S0 galaxies have a lot of molecular gas together with a moderate level of star-formation activity (e.g., Thronson et al. 1989).
It is known that there is a morphological-type segregation between nuclear starbursts and LINERs (Ho et al. 1997a); i.e., starbursts appear preferentially in late-type spirals while LINERs do in early-type ones. As for this issue, we give the following three comments. 1) A possible activity classification bias: If one finds evidence both for an AGN and for a nuclear starburst in a galaxy, one tends to classify it as an AGN. It is noted that infrared evidence for nuclear starbursts has been found even in a significantly large number of early-type spiral galaxies Devereux 1987, 1994). 2) A possible spectroscopic bias: Emission-line classification based on optical spectroscopy with a large aperture may be sometimes unreliable. For example, even if a galaxy has a pure AGN, it may be misclassified as a nuclear-starburst galaxy if it has very luminous circumnuclear star-forming regions (Yoshida et al. 1993). This means that one may misidentify activity types of late-type spirals with an AGN. And, 3) a morphological-evolution effect: It seems likely that a strong nuclear starburst may modify the appearance of a galaxy so that the host tends to be classified as an earlier-type spiral (Alonso-Herrero et al. 2000; Wada, Habe, & Sofue 1995). If starbursts were triggered by minor-mergers of satellite galaxies, the morphology of host galaxies could be altered (Mihos & Hernquist 1994; Hernquist & Mihos 1995; Taniguchi & Wada 1996).
## 4. SUMMARY
We have presented a new poststarburst model of LINERs. In this model, the ionization sources are planetary nebula nuclei (PNNs). Our main point is that the ionization sources are planetary nebula nuclei (PNNs) with temperature of $`10^5`$ K which appear in the late-phase evolution of intermediate-mass stars with mass between $`3M_{}`$ and $`6M_{}`$. Our models are able to reproduce the observed optical narrow emission-line ratios of LINERs although the \[O i\] emission is underpredicted to some extent.
We give a summary of the limitations of our models. 1) Our models cannot be applied to LINERs associated with nuclei of elliptical galaxies and most of S0 galaxies in which nuclear starbursts occur seldom. 2) Our models cannot be applied to LINERs with direct evidence for AGN such as broad-line emission, radio jets, hard X-ray emission, and so on (mostly type 1 LINERs). 3) Our models cannot be applied to on-going starburst nuclei because the photoionization is dominated by massive stars in the starburst rather than a cluster of PNNs left from a recent past starburst. 4) Accordingly, our models are applied to a subset of type 2 LINERs located in nuclear regions of spiral galaxies. Since the H$`\alpha `$ luminosity in the poststarburst LINER phase is less luminous by $``$ 4 orders of magnitude than the initial H$`\alpha `$ luminosity at the onset of the starburst, our models are preferentially applied to low-luminosity LINERs; e.g., $`L`$(H$`\alpha `$) $`10^{38}`$ ergs s<sup>-1</sup>.
Our models are constructed to explain the observed optical narrow emission-line ratios and we have not examined properties at other wavelengths (e.g., radio continuum, hard X rays, and so on). However, since the ionization sources are PNNs, poststarburst LINERs will not show any evidence for the presence of AGNs such as radio jets, hard X-ray emission and so on. As shown in Figure 3, our models predict that the X-ray emission is much weaker than that of typical AGNs, being consistent with some type 2 LINERs (Terashima et al. 2000 and references therein; see also for AGN-like SEDs of LINERs, Ho 1999). One interesting prediction of our models is that the ionization sources may be spatially extended because massive stars in nuclear starbursts are often distributed within central 100-pc regions (Meurer et al. 1995; see also Sugai & Taniguchi 1992). Recently, Pogge et al. (2000) have shown that some LINERs have spatially-extended emission-line regions with sizes of tens to hundreds parsecs. However, since they have either a compact UV source or no bright UV source, they may be not poststarburst LINERs but genuine low-ionization AGNs. In conclusion, in order to examine how many poststarburst LINERs are really present, we need systematic investigations of radio-continuum and X-ray properties of low-luminosity LINERs as well as UV continuum imaging in future.
We would like to thank an anonymous referee for many useful comments and suggestions. YS and TM are supported by JSPS. This work was financially supported in part by Grant-in-Aids for the Scientific Research (Nos. 10044052, and 10304013) of the Japanese Ministry of Education, Culture, Sports, and Science.
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# Macroscopic quantum tunneling and resonances in coupled Bose-Einstein condensates with oscillating atomic scattering length
## I Introduction
Recently much attention has been paid to the investigation of macroscopic tunneling phenomena between coupled Bose-Einstein condensates (BEC) . The theory has been developed for weakly as well as for strongly overlapped condensates -. As a result, it was shown that, in oscillations of the relative atomic population, periodic processes of two types can exist. One which is characterized by a zero mean atomic inbalance function $`(<z(t)>=0)`$, corresponding to the periodic flux of atoms between two condensates due to the overlaps of the wavefunctions. The second type has a net atomic population balance $`<z(t)>0`$, corresponding to the so-called macroscopic self-trapping , with localization of atoms in one of the condensates, and a periodically varying population around a constant value. The different regimes are connected with the value of the initial phase differences $`\theta =\varphi _1\varphi _2`$ between the condensates and the effective nonlinearity $`\mathrm{\Lambda }(\alpha _1+\alpha _2)N`$, where $`\alpha _i`$ are the nonlinearity parameters depending on the atomic scattering length $`a_s`$ (which in turn depends on the two-atoms interactions), and $`N`$ is the total number of atoms.
If a periodic variation of the parameters of the trap and of the BEC is introduced, new phenomena show up. Examples of this are the resonant phenomena with a periodically varying trap potential considered recently for the case of weakly coupled BEC in and for the strongly overlaped BEC in . The optical analog of this process is the electromagnetic wave transmission in a nonlinear coupler with periodic variation of tunnel coupling . Further interesting cases occur when a time-dependent atomic scattering length in BEC is considered. This is made possible due to the fact that $`a_s`$ depends resonantly on the external magnetic field or on the laser field -. This means that the scattering length can be varied in time, in particular, periodically. The main interest now is to investigate the influence of this periodic variation of atomic scattering length on the macroscopic interference phenomena. The central point here consists in that, as the interference phenomena are time-periodic processes in a nonlinear system, we can expect a nonlinear resonant response of the system under time-periodic perturbations, leading to a variety of phenomena, such as hysteresis or chaos.
## II Formulation of problem. Slowly time-varying atomic scattering length
The problem of weakly coupled BEC’s in a double-well trap potential with time-dependent atomic scattering length $`a_s`$ can be described by the two-modes model
$`ih{\displaystyle \frac{\psi _1}{t}}=(E_1+\alpha _1(t)|\psi _1|^2)\psi _1K\psi _2,`$ (1)
$`ih{\displaystyle \frac{\psi _2}{t}}=(E_2+\alpha _2(t)\psi _2|^2)\psi _2K\psi _1,`$ (2)
where the parameters $`E_i,\alpha _i,K`$ are defined by the overlaps integrals of the time dependent Gross-Pitaevsky eigenfunctions and derived in . In particular, $`\alpha _i(t)`$ are the parameters characterizing the nonlinear interactions between atoms and $`\alpha _ig_0=4\pi h^2a_s(t)/m`$, where $`a_s(t)`$ is the time dependent atomic scattering length. This system of equations is valid in the approximation of a weak link between condensates. As comparison with the numerical solution of GPE shows that it is a good aproximation for $`z0.6.`$ (see the definition of $`z(t)`$ below). Anther limitation is connected with the time dependence of $`a_s`$. For a harmonic modulation in time we should require that the perturbation does not introduce transitions between the ground state and excited states in the trap, whose energy gap is of the order of $`\mathrm{}\omega _0`$, where $`\omega _0`$ is the harmonic frequency of trap. Then we have the restriction $`\mathrm{\Omega }\omega _0`$. Another characteristic frequency is given by the coupling between condensates and is defined as $`\omega _L=2K`$. For the weak coupling case when $`\omega _L\omega _0`$ the regimes of resonant $`\mathrm{\Omega }\omega _L`$ and rapid modulations of $`a_s`$ can be realized. The periodic modulation of nonlinearity in the Gross-Pitaevsky equation can also lead to the parametric instability of the collective excitations of the condensate , that, in turn, leads to the break down of the two-modes approximation. Because these frequencies are far from the ones considered below we wiil use the model (1) in our analysis.
Introducing new variables $`\psi _i=\sqrt{N_i}\mathrm{exp}(i\theta _i),z=(N_1N_2)/N_T,N_T=N_1+N_2,\mathrm{\Phi }=\theta _1\theta _2`$, where $`N_i,\theta _i`$ are the number of atoms and phases in the i-th trap, we get the system
$`z_t=\sqrt{1z^2}\mathrm{sin}\mathrm{\Phi },`$ (3)
$`\mathrm{\Phi }_t=\nu \mathrm{\Lambda }(t)z+{\displaystyle \frac{z}{\sqrt{1z^2}}}\mathrm{cos}\mathrm{\Phi }+\mathrm{\Delta }E_0,`$ (4)
where $`\mathrm{\Lambda }=(\alpha _1(t)+\alpha _2(t))N_T/2K`$, $`t=\omega _Lt`$, $`\omega _L=2K`$, $`\omega _L`$ being the frequency of the linear Rabi oscillations, $`\mathrm{\Delta }E_0=\mathrm{\Delta }E/2K`$, and $`\nu =\pm 1`$ for the positive and negative atomic scattering length, respectively.
The Hamiltonian of the unperturbed system ($`a_s(t)=\mathrm{const}`$, $`\mathrm{\Delta }E_0=\mathrm{const}`$, $`\eta =0`$) is
$$H=\frac{\nu \mathrm{\Lambda }z^2}{2}\sqrt{1z^2}\mathrm{cos}(\mathrm{\Phi })+\mathrm{\Delta }E_0z.$$
(5)
Below, if not specified otherwise, we will consider the case of positive scattering length $`\nu =1`$,and $`\mathrm{\Delta }E_0=0`$. Consider the initial phase difference to be zero, $`\mathrm{\Phi }(0)=0`$ – the zero-phase mode. Depending on the number of atoms, $`N`$, and two-atoms interaction, $`\alpha `$, i.e. $`\mathrm{\Lambda }`$, there may exist a macroscopic quantum tunneling regime with $`<z(t)>=0`$,corresponding to $`\mathrm{\Lambda }<\mathrm{\Lambda }_c`$ and a self-trapped regime with $`<z(t)>0`$, corresponding to $`\mathrm{\Lambda }>\mathrm{\Lambda }_c`$. The same is valid for the initial phase difference, $`\mathrm{\Phi }(0)=\pi `$ – the $`\pi `$-phase mode.
Let us consider now the case of periodic modulations of the atomic scattering length, when $`a_s(t)=a_0(1+ϵ\mathrm{cos}\mathrm{\Omega }t)`$, inducing a variation of $`\mathrm{\Lambda }`$, i.e, $`\mathrm{\Lambda }=\mathrm{\Lambda }_0+\mathrm{\Lambda }_1\mathrm{cos}(\mathrm{\Omega }t)`$. Such a variation can be achieved by the variation of the external magnetic field \- a magnetic field induced Feshbach resonances for example. Near the resonance the scattering length is varied dispersively and can have positive and negative values. We will consider situations when we are not close to resonance and variations of scattering length are small, that is, $`\mathrm{\Lambda }_1/\mathrm{\Lambda }_0=\epsilon 1`$.
### A Resonances in zero phase modes
Let us first consider the case where $`\mathrm{\Phi }(0)=0`$. Take the case of $`z^21`$ and $`\mathrm{\Lambda }_1>z`$. Taking into account that $`z_t\mathrm{\Phi }z^2`$, assumed to be small, we obtain the equation for the relative atomic population
$$z_{tt}+(\mathrm{\Lambda }(t)+1)z\frac{\mathrm{\Lambda }(t)+1}{2}z^3=0.$$
(6)
where, $`\mathrm{\Lambda }(t)=\mathrm{\Lambda }_0(1+ϵ\mathrm{sin}(\mathrm{\Omega }t))`$. For $`ϵ1`$ we can reduce the equation to
$$z_{tt}+\omega _0^2z\beta z^3+ϵ\mathrm{\Lambda }_0\mathrm{cos}(\mathrm{\Omega }t)z=0,$$
(7)
where $`\omega _0=\sqrt{1+\mathrm{\Lambda }_0}`$ is the frequency of linear oscillations, and $`\beta =(\mathrm{\Lambda }_0+1)/2`$. It is useful to introduce a new variable $`y=\sqrt{(\mathrm{\Lambda }_0+1)}z1`$ and a scaled time $`\tau =t\omega _0`$. Then we have the equation
$$y_{\tau \tau }+y\beta _1y^3+\epsilon _1\mathrm{cos}(\mathrm{\Omega }_1\tau )y=0,$$
(8)
where $`\beta _1=1/(2\sqrt{\mathrm{\Lambda }_0+1}),\epsilon _1=\epsilon \mathrm{\Lambda }_0/(\mathrm{\Lambda }_0+1),\mathrm{\Omega }_1=\mathrm{\Omega }/\omega _0`$.
Let us consider the dynamics in the parametric resonance region when $`\mathrm{\Omega }_1=2+\mathrm{\Delta }`$. Using the results of perturbation theory , we find that the parametric resonance occurs when the condition for the square of increment $`s^2>0`$ is satisfied, where
$$s^2=\frac{1}{4}[(\frac{\epsilon _1}{2}(1\frac{3\beta _1a^2}{8}))^2\mathrm{\Delta }^2].$$
(9)
We have included here a nonlinear correction to the frequency $`3\beta _1a^2/8`$. Then we obtain the limit for the amplitude for parametric resonance to occur,
$$a_p^2\frac{8(\epsilon _12\mathrm{\Delta })}{3\beta _1\epsilon _1},$$
(10)
and $`z_p=a_p/\sqrt{\mathrm{\Lambda }_0+1}`$.
In Fig.(1) we present the oscillations of $`z(t)`$ when the parameters are in the region of the parametric resonance, by integrating numerically Eq.(3).
### B Resonances in $`\pi `$ \- phase modes
We now come to the case $`\varphi (0)=\pi `$, with for $`\mathrm{\Lambda }<1`$ and $`z^2<<1`$. The system (3) is simplified and reduces to one equation for $`\mathrm{\Phi }`$, having the form of the equation of a double pendulum. Inspection of this equation shows that there exists a valley in the effective energy $`V(\mathrm{\Phi })`$ around $`\mathrm{\Phi }=\pi `$. The maximum of the valley depth is achieved when $`\mathrm{\Lambda }_01`$. Thus we can search the solution of the system in this region of parameters assuming $`z1,\mathrm{\Phi }=\pi +\delta (t)`$. It results that we obtain the following equation for $`z(t)`$
$$z_{tt}+(1\mathrm{\Lambda }_0)z\mathrm{\Lambda }_1\mathrm{cos}(\mathrm{\Omega }t)z=0.$$
(11)
Here $`\mathrm{\Lambda }_0<1`$ .The parametric resonance in the oscillations of the atomic population appears when $`\mathrm{\Omega }=2\sqrt{1\mathrm{\Lambda }_0}+\mathrm{\Delta }`$. The parametric instability occurs when $`\mathrm{\Lambda }_1/(2\sqrt{1\mathrm{\Lambda }_0})\mathrm{\Delta }`$. The increment $`s`$ is, in this case,
$$s=\frac{1}{2}\sqrt{(\frac{\mathrm{\Lambda }_1}{2\sqrt{1\mathrm{\Lambda }_0}})^2\mathrm{\Delta }^2}.$$
(12)
For example, when $`\mathrm{\Lambda }_0=0.36,\mathrm{\Delta }=0.1`$ we have $`\mathrm{\Omega }=1.7`$. In Fig.(2) we plot the oscillations of $`z(t)`$ for this choice of parameters, and with $`\mathrm{\Lambda }_1=0.162`$, by a nummerical integration of Eq.(3).
Note that the frequency of linear Rabi oscillations is $`\omega _L=2K`$, or one in the units used. Thus the resonance in $`\pi `$ -mode occurs at frequencies lower than the Rabi frequency. It is worth mentioning that the metastable $`\pi `$-phase mode has been recently observed in the weak link separating two reservoirs of superfluid He<sup>3</sup> .
## III Chaotic dynamics in oscillations of the relative atomic population
The oscillations induced by the periodic perturbations of the scattering length may be comples and become chaotic, for certain regions of the parameters. Note that the unperturbed system is equivalent to the Duffing oscillator. Thus, as is known, the resonance overlaps under periodic perturbations can lead to chaotic variations of $`z(t)`$ . In such cases, to study the chaotic motion, it is useful to look for the motion near the separatrix of the unperturbed system. The total Hamiltonian is $`H=H_0+V`$, where $`H_0`$ is unperturbed Hamiltonian (4) and the perturbation is $`V=\mathrm{\Lambda }_1\mathrm{cos}(\mathrm{\Omega }t)z^2/2`$. The value of the Hamiltonian on separatrix is $`H_s=1`$. The separatrix divides the regions with macroscopic quantum tunnelling, with $`<z>=0`$, from the self-trapped regions where $`<z>0`$. The expressions for the separatrix solution are:
$`z_s(t)=\sqrt{{\displaystyle \frac{a}{b}}}\text{sech}(\sqrt{2a}t),`$ (13)
$`\mathrm{sin}^2(\mathrm{\Phi }_s)={\displaystyle \frac{a^2\text{sech}^2(\sqrt{2a}t)\mathrm{tanh}^2(\sqrt{2a}t)}{2K_0^2(ba\text{sech}^2(\sqrt{2a}t))}},`$ (14)
where $`a=|1\mathrm{\Lambda }|/2`$, $`b=\mathrm{\Lambda }^2/8`$. For example, for the zero-phase mode the separatrix is given by the hyperbolic fixed points $`\mathrm{\Phi }=\pm \pi ,z=0`$.
The existence of chaos can be proved by calculating the Melnikov function. The calculations are similar to the ones performed in and shows that the Melnikov function has an the infinite number of zeros, so the existence of chaotic regimes in the atomic population oscillations is proved. In Fig.(3) we plot the Melnikov distance versus the frequency $`\mathrm{\Omega }`$ for $`\mathrm{\Lambda }=10`$. We see that the maximum value is achieved for $`\mathrm{\Omega }3.89`$. In Fig.4 we plot the difference of atomic population $`z(t)`$ as function of time for $`\mathrm{\Lambda }_0=10,z(0)=0.4,\mathrm{\Lambda }_1=0.2,\mathrm{\Omega }=3.89`$.
It is interesting to obtain an analytical estimate for the critical amplitude of modulation $`\mathrm{\Lambda }_1`$ leading to the chaotic behavior in the relative atomic population. Using the expression for the separatrix (13) we can calculate the Melnikov-Arnold integral and find that the energy change for a perturbation $`\mathrm{\Lambda }_1\mathrm{cos}(\mathrm{\Omega }t+\psi )`$(where $`\psi `$ is the phase) is
$$\mathrm{\Delta }H=_{\mathrm{}}^{\mathrm{}}(\frac{H_0}{t}+[H_0,V])𝑑t=\alpha \mathrm{sin}(\psi ),\alpha =\frac{\pi \epsilon |1\mathrm{\Lambda }|\mathrm{\Omega }^2}{\mathrm{\Lambda }\mathrm{sinh}(\pi \mathrm{\Omega }/2\sqrt{\mathrm{\Lambda }1})},$$
(15)
where $`[\mathrm{}]`$ denotes the Poisson bracket.
The period of the unperturbed motion is $`T=4\kappa K(\kappa )/(C\mathrm{\Lambda })`$, where :
$$\kappa ^2=\frac{1}{2}[1+\frac{H_0\mathrm{\Lambda }1}{\sqrt{\mathrm{\Lambda }^2+12H_0\mathrm{\Lambda }}}],$$
(16)
$$C^2=\frac{2}{\mathrm{\Lambda }^2}[(H_0\mathrm{\Lambda }1)+\sqrt{\mathrm{\Lambda }^2+12H_0\mathrm{\Lambda }}],$$
(17)
where $`K(\kappa )`$ is the complete elliptic integral of the first kind. Near the separatrix, when $`H_01,`$ we have $`\kappa ^21`$. Introducing the parameter $`\delta =1H_0`$, $`H_0<1`$, $`\delta <<1`$, and taking into account that $`K(\kappa )\mathrm{ln}(4/\sqrt{1\kappa ^2})`$, when $`\kappa ^21`$, we obtain the estimate for the period of oscillations near the separatrix
$$T\frac{2}{\sqrt{\mathrm{\Lambda }1}}\mathrm{ln}\left(\frac{4\sqrt{2(\mathrm{\Lambda }1)}}{\sqrt{\mathrm{\Lambda }\delta }}\right).$$
(18)
Using the expressions (15) and (18) we find the whisker-type map for our problem
$`\delta _{n+1}=\delta _n+\alpha \mathrm{sin}(\psi _n),`$ (19)
$`\psi _{n+1}=\psi _n+\gamma \mathrm{ln}\left({\displaystyle \frac{\mu }{\sqrt{\delta _{n+1}}}}\right).`$ (20)
where $`\gamma =2\mathrm{\Omega }/\sqrt{\mathrm{\Lambda }1}`$ and $`\mu =4\sqrt{2(\mathrm{\Lambda }1)}/\sqrt{\mathrm{\Lambda }}`$.
This map can be simplified, using the linearization around the fixed points . The fixed points are defined by
$$2\pi l=\gamma \mathrm{ln}\left(\frac{\mu }{\sqrt{\delta _f^{(l)}}}\right),l=1,2\mathrm{}$$
(21)
Let us introduce the dimensionless energy $`I_n`$, defined by
$$I_n=\frac{\gamma }{\delta _f^{(l)}}(\delta _n\delta _f^{(l)}).$$
(22)
Substituting this expression into Eq.(19), and redefining $`I_nI_n/2`$, we obtain the standard map
$`I_{n+1}=I_nK\mathrm{sin}(\psi _n),`$ (23)
$`\psi _{n+1}=\psi _n+I_{n+1},`$ (24)
where the parameter $`K`$ is
$$K=\frac{\alpha \gamma }{2\delta _f^{(l)}}.$$
(25)
As is well known, chaos appears when $`K1`$ (exact value is $`K^{}=0.9716`$). Then we have the estimate for the critical amplitude of modulation
$$\mathrm{\Lambda }_1\frac{\delta _f^{(l)}\mathrm{\Lambda }^3\mathrm{sinh}(\pi \mathrm{\Omega }/2\sqrt{\mathrm{\Lambda }1})}{\pi \sqrt{\mathrm{\Lambda }1}\mathrm{\Omega }^3}.$$
(26)
From the above, we may calculate the diffusion coefficient for the motion in the stochastic layer
$$D=\frac{<(\mathrm{\Delta }E)^2>}{t}.$$
(27)
It results that we have $`D=K^2/2`$. For frequencies $`\mathrm{\Omega }2\sqrt{\mathrm{\Lambda }1}/\pi `$ we have the estimate
$$D\frac{2^{10}\pi ^2\epsilon ^2(\mathrm{\Lambda }1)^3\mathrm{\Omega }^4}{\mathrm{\Lambda }^2}\mathrm{exp}(\frac{4\pi l\sqrt{\mathrm{\Lambda }1}}{\mathrm{\Omega }}).$$
(28)
For $`\epsilon =0.1`$, $`\mathrm{\Omega }=2.5`$, $`\mathrm{\Lambda }=9`$, $`l=1`$, the time of exit from the stochastic layer is $`10`$ i.e. $`4`$ periods of the oscillations of the atomic scattering length.
## IV Rapidly varying scattering length
In the case when the frequency of oscillations of the atomic scattering length is much larger than the tunneling frequency $`\omega _L=2K`$. i.e. $`\mathrm{\Lambda }=\mathrm{\Lambda }(t/ϵ)`$, where $`ϵ=\omega _L/\mathrm{\Omega }<<1`$, it is useful to derive an averaged set of equations. Using a multiscale approach and expanding $`z=\overline{z}+ϵz_1+\mathrm{},\mathrm{\Phi }=\overline{\mathrm{\Phi }}+ϵ\mathrm{\Phi }_1\mathrm{}.`$ we can derive the system of equations for the slowly varying functions $`\overline{z},\overline{\mathrm{\Phi }}`$. The limits of validity is $`\overline{z}^2<<1`$.
The averaged system turns out to be:
$`\overline{z}_t=\sqrt{1\overline{z}^2}\mathrm{sin}(\overline{\mathrm{\Phi }})(1\delta \overline{z}^2),`$ (29)
$`\overline{\mathrm{\Phi }}_t=\nu <\mathrm{\Lambda }>\overline{z}+{\displaystyle \frac{\overline{z}\mathrm{cos}(\overline{\mathrm{\Phi }})}{\sqrt{1\overline{z}^2}}}(1+2\delta 3\delta \overline{z}^2).`$ (30)
where $`\delta `$ is proporttional to $`ϵ^2`$, and for the harmonic modulation is $`\mathrm{\Lambda }_1^2/(4\mathrm{\Omega }^2)`$. The corresponding Hamiltonian is
$$H=\frac{\nu <\mathrm{\Lambda }>}{2}\overline{z}^2\sqrt{1\overline{z}^2}\mathrm{cos}(\overline{\mathrm{\Phi }})(1\delta \overline{z}^2).$$
(31)
With respect to the unperturbed case, a new stable point can occur. A bifurcation is observed when $`\delta `$ is increased. The unstable point (in the undisturbed case) at $`z=0,\mathrm{\Phi }=\pi `$ may become stable under rapid perturbations. This phenomenon is the analog of the Kapitza stabilization of the unstable fixed point of the pendulum by rapidly oscillating the suspension point. By passing though the bifurcation, a change in the topology of the phase portrait happens. The value of $`\delta `$ for this to occur depends on $`<\mathrm{\Lambda }>`$. To have an estimate of the critical value of $`\delta `$, consider $`z^2<<1`$ and $`\mathrm{\Phi }=\pi +\psi ,\psi <<\pi `$. The equation for $`z`$ reads:
$$z_{tt}(<\mathrm{\Lambda }>12\delta )z=0.$$
(32)
Thus, when $`\delta >\delta _c=(<\mathrm{\Lambda }>1)/2`$ we have oscillatory motion around $`(0,\pi )`$ in the phase plane. When $`\delta <\delta _c`$ we have an hyperbolic point, and $`(0,\pi )`$ is no longer stable. For example, with $`\mathrm{\Lambda }=2`$, the critical $`\delta `$ is $`0.50`$. For $`<\mathrm{\Lambda }>=1`$, the point is always stable, as $`\delta _c=0`$. To support these results, we have numerically integrated the original systems given by Eq.(3) with an explicit harmonic fast-time variation for $`\mathrm{\Lambda }(t)`$. The above described features immediatly show up.
Note also that we can cross the sign of $`a_s`$ from positive to negative. Experiments with very a large monotonic change of sign show that the BEC shrinks very fast. It would be interesting to check if the condensate with attractive interaction of atoms stable under the rapid variation of scattering length. This problem requires, however, separate investigations. In any case in quasi $`1D`$ geometries collapse is suppressed and this analysis is relevant.
In order to better understand the physical consequences of the averaging, we note that for small $`\overline{z}^2<<1`$ we can rewrite the Hamiltonian in the form
$$H=\frac{\nu \mathrm{\Lambda }_rz_r^2}{2}\sqrt{1z_r^2}\mathrm{cos}(\mathrm{\Phi }),$$
where $`z_r=(1+\delta )\overline{z}`$, $`\mathrm{\Lambda }_r=(12\delta )<\mathrm{\Lambda }>`$. Thus we can conclude that the result of the averaging consists, for a fixed initial phase difference between condensates, in the lowering of the critical relative population for self-trapping regime to occured and in the increase of the critical value of the nonlinear parameter $`\mathrm{\Lambda }`$. We come to the picture of a more rigid pendulum in comparison with the constant case. The threshold for the switching from the macroscopic quantum tunneling regime to the self-trapping regime is shifted to lower initial values of the atomic imbalance.
## V Conclusion
In this paper we have studied the new effects occuring due to time-periodic variations of the atomic scattering length. The resonances in oscillations of the atomic imbalance for the different regimes have been analyzed.
The interaction of resonances occuring under oscillating scattering length can lead to complicated behavior of the atomic oscillations and gives rise to the phenomenon of chaotic macroscopic quantum tunneling. In this work we prove the existence of such regime, using the Melnikov function approach and calculating the Melnikov distance, characterizing the width of the stochastic layer near the separatrix of the unperturbed system. We derive a whisker type map for the problem and obtain the estimate for the critical value of amplitude $`\mathrm{\Lambda }_1`$ leading to the chaotic motion. The diffusion coefficient for the motion in the stochastic layer is calculated and shown that the time to cross the stochastic layer is the order of few periods of modulations. We also consider the evolution of the system under rapidly varying oscillations of the scattering length and derive the averaged system for the slow-time variations of the relative atomic population $`\overline{z}(t)`$ and the phase $`\overline{\mathrm{\Phi }}(t)`$. The analysis of the fixed points shows that the stabilization of the system under rapid perturbation, when the unpertubed is unstable, is possible in the $`\pi `$ -phase mode.
## VI Acknowledgments
F.Kh.A. is grateful to FAPESP (Brazil) and INTAS (Grant 96-339) for partial financial support, and to Instituto de Física Teórica – UNESP, São Paulo, for hospitality. R.A.K. thanks FAPESP and CNPq (Brazil) for partial financial support.
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# Anti-self-dual Yang-Mills equations on noncommutative spacetime
## 1 Introduction
The deformed ADHM construction of Nekrasov and Schwarz suggests that the anti-self-dual Yang-Mills (ASDYM) equations will be “integrable” on noncommutative spacetimes as well. This is also advocated by the work of Kapustin et al. that extends the ordinary twistorial interpretation of the ADHM construction to the noncommutative $`^4`$. Since twistor theory is a clue to the integrability of the ordinary ASDYM equations , it is natural to expect that the ASDYM equations on the noncommutative $`^4`$, too, will be integrable.
This issue is also interesting from the point of view of integrable systems of two-dimensional field theories, such as the principal chiral field (PCF) model and Hitchin’s Higgs pair equations . It is well known that these integrable systems can be derived from the ASDYM equations by dimensional reduction. If a similar reduction procedure works on the noncommutative $`^4`$, it seems likely that the integrability of the four-dimensional system will be inherited by the two-dimensional systems. This four-dimensional point of view can be an alternative approach to recent studies on the PCF and Wess-Zumino-Witten (WZW) models on noncommutative spacetimes .
This paper aims to answer these questions. The gauge group is assumed to be $`U(N)`$ throughout the paper. The ASDYM equations on the noncommutative $`^4`$ are then obtained from the ordinary ASDYM equations by replacing the product of fields in the field equations with the so called “$``$-product” (the commutator of which is the Moyal bracket ). We shall show that almost all part of the twistorial and integrable structures of the ordinary ASDYM equations can be extended to the noncommutative $`^4`$ by the same substitution rule. What breaks down is the part where tools of finite dimensional linear algebra (determinants, Camer’s formula, characteristic polynomials, etc.) are used.
This paper is organized as follows. Section 2 presents the formulation of the $``$-product ASDYM equations. Section 3 deals with the twistorial and integrable structures of the $``$-product ASDYM equations. Section 4 is concerned with some implications of the deformed ADHM construction. Section 5 is devoted to two-dimensional reductions. Section 6 is for conclusion.
## 2 ASDYM equations on noncommutative $`^4`$
### 2.1 Spacetime coordinates
The noncommutative $`^4`$ is characterized by the commutation relations
$`[x_j,x_k]=i\theta _{jk}`$ (2.1)
of the spacetime coordinates, where $`\theta _{jk}`$ are real constants. These commutation relations can be extended to the associative $``$-product
$`fg(x)=\mathrm{exp}\left({\displaystyle \underset{j,k=1}{\overset{4}{}}}{\displaystyle \frac{i}{2}}\theta _{jk}_{x_j}_{\stackrel{~}{x}_k}\right)f(x)g(\stackrel{~}{x})|_{\stackrel{~}{x}=x}`$ (2.2)
of functions $`f`$ and $`g`$ on the spacetime.
We now introduce complex coordinates $`(z_1,z_2)`$ that satisfy commutation relations of the form
$`[z_1,z_2]=\zeta _{},[\overline{z}_1,\overline{z}_2]=\overline{\zeta }_{},[z_1,\overline{z}_1]+[z_2,\overline{z}_2]=\zeta _{}.`$ (2.3)
For instance, $`z_1=x_3+ix_4`$ and $`z_2=x_1+ix_2`$ give such complex coordinates after a suitable orthogonal transformation of the real coordinates. The complex constant $`\zeta _{}`$ and the real constant $`\zeta _{}`$ form a three-vector $`(\zeta _{},\zeta _{})`$ in $`\times ^3`$, and can be rotated to any direction by the $`SU(2)`$ action
$`\left(\begin{array}{cc}z_1& z_2\\ \overline{z}_2& \overline{z}_1\end{array}\right)\left(\begin{array}{cc}\alpha & \beta \\ \overline{\beta }& \overline{\alpha }\end{array}\right)\left(\begin{array}{cc}z_1& z_2\\ \overline{z}_2& \overline{z}_1\end{array}\right)`$ (2.10)
($`|\alpha |^2+|\beta |^2=1`$) on the spacetime coordinates.
### 2.2 ASDYM equations
Let $`A`$ be an $`N\times N`$ matrix-valued one-form representing a $`U(N)`$-connection. Let $`A_1,A_2,A_3,A_4`$ denote the components in the real coordinate frame, and $`A_{z_1},A_{z_2},A_{\overline{z}_1},A_{\overline{z}_2}`$ the components in the complex coordinate frame:
$`A={\displaystyle \underset{j=1}{\overset{4}{}}}A_jdx_j={\displaystyle \underset{a=1,2}{}}A_{z_a}dz_a+{\displaystyle \underset{a=1,2}{}}A_{\overline{z}_a}d\overline{z}_a.`$ (2.11)
The covariant derivatives can be accordingly written
$`_{x_j}=_{x_j}+A_j,_{z_a}=_{z_a}+A_{z_a},_{\overline{z}_a}=_{\overline{z}_a}+A_{\overline{z}_a}.`$ (2.12)
On the noncommutative $`^4`$, the components of the curvature two-form $`F=_{j,k}F_{jk}dx_jdx_k`$ are defined as
$`F_{jk}=_{x_j}A_k_{x_k}A_j+[A_j,A_k]_{}.`$ (2.13)
Note that the usual matrix commutators $`[A_j,A_k]=A_jA_kA_kA_j`$ are now replaced by the $``$-product commutators
$`[A_j,A_k]_{}=A_jA_kA_kA_j.`$ (2.14)
The components $`F_{z_1z_2}`$, $`F_{\overline{z}_1\overline{z}_2}`$ and $`F_{z_a\overline{z}_b}`$ in the complex coordinate frame are similarly written in terms of $`A_{z_a}`$ and $`A_{\overline{z}_a}`$. The ASDYM equations in the complex coordinate frame take the neat form
$`F_{z_1z_2}=0,F_{\overline{z}_1\overline{z}_2}=0,F_{z_1\overline{z}_1}+F_{z_2\overline{z}_2}=0.`$ (2.15)
### 2.3 Reduced form of ASDYM equations
It is well known that the ASDYM equations can be converted to a (classical) field theory with a Lagrangian formalism. Actually, two types of such expressions are known. One is Yang’s equation (also called the four-dimensional Donaldson-Nair-Schiff equation ). Another expression is due to Leznov and Parkes . Both can be extended to the noncommutative spacetime as follows.
Let us consider the first equation $`F_{z_1z_2}=0`$ of the ASDYM equations. This is a partial (two dimensional) flatness condition. In the ordinary (complexified) spacetime, this implies that $`A_{z_1}`$ and and $`A_{z_2}`$ can be expressed as
$`A_{z_1}=h^1_{z_1}h,A_{z_2}=h^1_{z_2}h`$ (2.16)
with an $`N\times N`$ matrix-valued function $`h`$ of the spacetime coordinates. (Of course, this is, in general, a local expression.) This expression persists to be true on the noncommutative spacetime if the ordinary product in the matrix multiplication are replaced by the $``$-product:
$`A_{z_1}=(h)_{}^1_{z_1}h,A_{z_2}=(h)_{}^1_{z_2}h.`$ (2.17)
Here $`(h)_{}^1`$ stands for an inverse with respect to the $``$-product, namely, $`h(h)_{}^1=(h)_{}^1h=1`$. One can prove this $``$-product version of Frobenius’ theorem in much the same way as a proof in the ordinary spacetime.
In order to derive a $``$-product analogue of Yang’s equation, we solve another flatness condition $`F_{\overline{z}_1,\overline{z}_2}=0`$ in the ASDYM equation as
$`A_{\overline{z}_1}=(k)_{}^1_{\overline{z}_1}k,A_{\overline{z}_2}=(k)_{}^1_{\overline{z}_2}k,`$ (2.18)
and the “matrix-ratio”
$`g=k(h)_{}^1`$ (2.19)
of $`h`$ and $`k`$. As we shall show below, this matrix-valued field turns out to obey the field equation
$`_{z_1}\left((g)_{}^1_{\overline{z}_1}g\right)+_{z_2}\left((g)_{}^1_{\overline{z}_2}g\right)=0`$ (2.20)
This gives an analogue of Yang’s equation on the noncommutative spacetime. The Lagrangian formalism in the commutative case can be readily extended to the noncommutative case. Note that the field equation
$`_z\left((g)_{}^1_{\overline{z}}g\right)+_{\overline{z}}\left((g)_{}^1_zg\right)=0`$ (2.21)
of the noncommutative PCF model can be derived by dimensional reduction.
The foregoing noncommutative analogue of Yang’s equation can be derived by the following trick. Let us consider the finite gauge transformation by $`h`$. Two of the four gauge potentials, $`A_{z_a}`$ ($`a=1,2`$), are thereby gauged away as
$`_{z_a}h_{z_a}(h)_{}^1=_{z_a},`$ (2.22)
and the other two are transformed as
$`_{\overline{z}_a}`$ $``$ $`h_{\overline{z}_a}(h)_{}^1`$ (2.23)
$`=`$ $`_{\overline{z}_a}_{\overline{z}_a}h(h)_{}^1+h(k)_{}^1_{\overline{z}_a}k(h)_{}^1`$
$`=`$ $`_{\overline{z}_a}+(g)_{}^1_{\overline{z}_a}g.`$
The gauge potentials are now in a half-flat gauge in which two of the gauge potentials vanish,
$`A_{z_1}=0,A_{z_2}=0,`$ (2.24)
and the other two gauge potentials are written as
$`A_{\overline{z}_a}=(g)_{}^1_{\overline{z}_a}g.`$ (2.25)
The remaining equation $`F_{z_1\overline{z}_1}+F_{z_2\overline{z}_2}=0`$ of the $``$-product ASDYM equations, which now takes the simplified form
$`_{z_1}A_{\overline{z}_1}+_{z_2}A_{\overline{z}_2}=0,`$ (2.26)
gives the $``$-product analogue of Yang’s equation.
One can see, from this derivation of Yang’s equation, the existence of a field theoretical “dual” of Yang’s equation as well. Note that the $``$-product ASDYM equations in the foregoing half-flat gauge with $`A_{z_1}=A_{z_2}=0`$ consist of the two equations
$`_{\overline{z}_1}A_{\overline{z}_2}_{\overline{z}_2}A_{\overline{z}_1}+[A_{\overline{z}_1},A_{\overline{z}_2}]_{}=0,_{z_1}A_{\overline{z}_1}+_{z_2}A_{\overline{z}_2}=0.`$ (2.27)
If one solves the first equation as a partial flatness condition, as we have seen above, Yang’s equation emerges from the second equation. Meanwhile, one can also solves the second equation as
$`A_{\overline{z}_1}=_{z_2}\varphi ,A_{\overline{z}_2}=_{z_1}\varphi `$ (2.28)
for a matrix-valued potential $`\varphi `$. The first equation then takes the form
$`(_{z_1}_{\overline{z}_1}+_{z_2}_{\overline{z}_2})\varphi +[_{z_1}\varphi ,_{z_2}\varphi ]_{}=0.`$ (2.29)
This is a $``$-product version of the field equation of Leznov and Parkes.
## 3 Twistor theory and integrability
### 3.1 Twistor geometry
Twistor theory encodes various fields on spacetime to a geometric structure on another (complex) manifold called the “twistor space” . In the case of four dimensional flat spacetime, the twistor space is the three dimensional complex projective space $`_{}^3`$. Roughtly speaking, twistor theory is a kind of “tomography”, namely, to “scan” the spacetime by a three-parameter family of two-dimensional surfaces (“twistor surfaces”) $`S(\xi )`$ labelled by the point $`\xi `$ of the twistor space. We review the essence of twistor geometry in the following.
To define the twistor surfaces, however, the real (Euclidean) spacetime $`^4`$ has to be extended to the complexified spacetime $`^4`$, in which $`(z_1,z_2,\overline{z}_1,\overline{z}_2)`$ are independent complex coordinates. The twistor surfaces in the complexified spacetime $`^4`$ are labelled by three parameters $`(\lambda ,u_1,u_2)`$, and defined by the equations
$`z_1\lambda \overline{z}_2=u_1,z_2+\lambda \overline{z}_1=u_2.`$ (3.1)
The parameters $`(\lambda ,u_1,u_2)`$ are local coordinates on a coordinate patch of the whole twistor space $`_{}^3`$. Furthermore, $`\lambda `$ turns out to play the role of the “spectral parameter” in the theory of integrable systems.
Various real spacetimes, such as the Minkowski spacetime and the spacetime with $`(2,2)`$ signature, are embedded in the complexified spacetime $`^4`$ as “real slices”. Although the twistor surface $`S(\lambda ,u_1,u_2)`$ intersects with the Euclidean spacetime at most at a point, the intersection with the Minkowski spacetime is a null line, and the intersection with the $`2+2`$ spacetime is a totally null surface (i.e., the inner product of any two tangent vectors vanish).
The twistor space $`_{}^3`$ itself appears in the description of a compacitified spacetime, such as the one-point compactification $`S^4=^4\{\mathrm{}\}`$ of the Euclidean spacetime. Let us introduce the complex Grassmann variety
$`\mathrm{Gr}_{}(2,4)=\{V_2V_2^4,dimV_2=2\}`$ (3.2)
of vector subspaces of $`^4`$ and the flag variety
$`\mathrm{Fl}_{}(1,2,4)=\{(V_1,V_2)V_1V_2^4,dimV_1=1,dimV_2=2\}`$ (3.3)
of pairs of nested vector subspaces of $`^4`$. The Grassmann variety is a natural complexification of $`S^4`$. Twistor geometry connects these compact (and complexified) spacetimes with the twistor space $`_{}^3`$ by the “Klein correspondence”
$`\mathrm{Gr}_{}(2,4)\stackrel{p_2}{}\mathrm{Fl}_{}(1,2,4)\stackrel{p_1}{}_{}^3,`$ (3.4)
where the projections $`p_1`$ and $`p_2`$ send the flag $`(V_1,V_2)`$ to $`V_1_{}^3`$ and $`V_2\mathrm{Gr}_{}(2,4)`$, respectively. The subset $`S(\xi )=p_2(p_1^1(\xi ))`$ is isomorphic to $`_C^2`$ and gives a compactification of the foregoing twistor surface $`S(\lambda ,u_1,u_2)`$ in $`^4`$. Similarly, the subset $`L(x)=p_1(p_2^1(x))`$ is isomorphic to $`_{}^1`$ and plays a key role in decoding the twistorial data.
The uncompactified spacetime $`^4`$ (or, rather, its complexification $`^4`$) can be described by the open twistor space
$`𝒯=_{}^3_{}^1`$ (3.5)
with a line $`_{}^1`$ deleted. It is rather this twistor space that we mostly consider in the following. This open twistor space has the projection
$`\begin{array}{ccc}\hfill \pi :𝒯& & _{}^1\hfill \\ \hfill \xi =[\xi _0:\xi _1:\xi _2:\xi _3]& & [\xi _0:\xi _1]\hfill \end{array}`$ (3.8)
and is covered by the two standard coordinate patches $`𝒰=\{\xi _00\}`$ and $`\widehat{𝒰}=\{\xi _10\}`$. The deleted line $`_{}^1`$ is the locus where $`\xi _0=\xi _1=0`$. The three parameters $`(\lambda ,u_1,u_2)`$ can be identified with the standard local coordinates on $`𝒰`$:
$`\lambda =\xi _1/\xi _0,u_1=\xi _2/\xi _0,u_2=\xi _3/\xi _0.`$ (3.9)
Thus, in particular, $`\lambda `$ is an affine coordinate of the base, and $`u_1`$ and $`u_2`$ are coordinates along the fibers.
### 3.2 Flatness on twistor surfaces
The three members of the ASDYM equations can be combined to a single equation of the form
$`F(_{\overline{z}_1}\lambda _{z_2},_{\overline{z}_2}+\lambda _{z_1})=F_{z_1z_2}\lambda (F_{z_1\overline{z}_1}+F_{z_2\overline{z}_2})+\lambda ^2F_{\overline{z}_1\overline{z}_2}=0.`$ (3.10)
Here $`F(v,v^{})`$ stands for the contraction of $`F`$ by two vector fields $`v,v^{}`$. Since the two vector fields $`_{\overline{z}_1}\lambda _{z_2},_{\overline{z}_2}+\lambda _{z_1}`$ on the left hand side span the tangent planes of the twistor surface $`S(\lambda ,u_1,u_2)`$, the foregoing equation means the flatness
$`F|_{S(\lambda ,u_1,u_2)}=0`$ (3.11)
of the connection on all twistor surfaces.
Frobenius’ theorem connects this flatness (or “zero-curvature”) condition with the integrability of the linear system
$`(_{\overline{z}_1}\lambda _{z_2})\mathrm{\Psi }(\lambda )=0,(_{\overline{z}_2}+\lambda _{z_1})\mathrm{\Psi }(\lambda )=0,`$ (3.12)
where $`\mathrm{\Psi }(\lambda )`$ is a vector- or matrix-valued unknown function (which, of course, depends on the spacetime coordinates as well). Having this linear system, one can now apply a number of techniques for integrable systems to the ASDYM equations .
Note that the first two equations of the ASDYM equations (from which $`h`$ and $`k`$ were derived) correspond to the flatness on the twistor surfaces with $`\lambda =0`$ and $`\lambda =\mathrm{}`$. Accordingly, one can choose two matrix-valued solutions $`\mathrm{\Psi }(\lambda )`$ and $`\widehat{\mathrm{\Psi }}(\lambda )`$ of (3.12) to be such that
$`\mathrm{\Psi }(0)=h,\widehat{\mathrm{\Psi }}(\mathrm{})=k.`$ (3.13)
In other words, $`\mathrm{\Psi }(\lambda )`$ and $`\widehat{\mathrm{\Psi }}(\lambda )`$ are one-parameter deformations of $`h`$ and $`k`$. Moreover, the Laurent expansion
$`\mathrm{\Psi }(\lambda )=h+{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}w_n\lambda ^n,\widehat{\mathrm{\Psi }}(\lambda )=k+{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\widehat{w}_n\lambda ^n`$ (3.14)
of these solutions of (3.12) are related to two infinite series of nonlocal conservation laws . There is no reason that these two solutions coincide. They rather give a pair that arise in the so called Riemann-Hilbert problem. We now turn to this issue.
### 3.3 Vector bundle and Riemann-Hilbert problem
The twistor transformation encodes a solution of the ASDYM equations to a holomorphic vector bundle $``$ over the twistor space. Given a solution of the ASDYM equations, one can consider an associated rank $`N`$ vector bundle $`E`$ over the spacetime with an induced connection. This connection is flat on each twistor surface $`S(\xi )`$. The fiber $`_\xi `$ of the bundle $``$ at a point $`\xi `$ of the twistor space is, by definition, the vector space of flat sections of $`E|_{S(\xi )}`$. In a down-to-earth language, the fiber $`_\xi `$ is the vector space of $`E`$-valued solutions of linear system (3.12) restricted to $`S(\lambda ,u_1,u_2)`$.
This bundle $``$ need not be defined over the whole twistor space. If the solution of the ASDYM equations is defined in a small neighborhood of a spacetime point $`x`$, the bundle $``$ is accordingly defined only in a neighborhood of the line $`L(x)=\{\xi _{}^3xS(\xi )\}`$. Instanton solutions are global solutions that give rise to a globally defined vector bundle on the whole twistor space.
The holomorphic vector bundle $``$ has the special property that the restriction $`|_{L(x)}`$ to the line $`L(x)_{}^1`$ is holomorphically trivial for any spacetime point $`x`$ in the domain where the gauge potentials are defined. This property of $``$ plays a key role in the inverse transformation, namely, to reproduce the solution of the ASDYM equation from the vector bundle $``$.
It is here that the notion of Riemann-Hilbert problem emerges. Let us recall that any holomorphic vector bundle over $`_{}^1`$ can be represented by the “patching function” $`p(\lambda )`$ on the intersection $`D\widehat{D}`$ of two affine coordinate patches $`\{D,\widehat{D}\}`$ of $`_{}^1`$. The patching function $`p(\lambda )`$ is an $`GL(N,)`$-valued holomorpihc function. If the vector bundle is holomorphically trivial, the patching function can be expressed as
$`p(\lambda )=\widehat{\mathrm{\Psi }}(\lambda )^1\mathrm{\Psi }(\lambda ),`$ (3.15)
where $`\mathrm{\Psi }(\lambda )`$ and $`\widehat{\mathrm{\Psi }}(\lambda )`$ are $`GL(N,)`$-valued holomorphic functions of $`\lambda `$ on $`D`$ and $`\widehat{D}`$, respectively. Finding such a pair of matrix-valued functions to the given data $`p(\lambda )`$ is a kind of Riemann-Hilbert problem (also called the “splitting” problem in the terminology of Ward).
The patching function $`p(\lambda )`$ is determined by a patching function of the vector bundle $``$ itself. As already remarked, the twistor space $`𝒯=_{}^3_{}^1`$ is covered by the two coordinate patches $`𝒰`$ and $`\widehat{𝒰}`$. The vector bundle $``$ is described by a $`GL(N,)`$-valued function $`P(\lambda ,u_1,u_2)`$ that glues together the rank-$`N`$ trivial bundles over $`𝒰`$ and $`\widehat{𝒰}`$. Its restriction on the line $`L(x)`$ is nothing but the patching function $`p(\lambda )`$ of $`|_{L(x)}`$:
$`p(\lambda )=P(\lambda ,z_1\lambda \overline{z}_2,z_2+\lambda \overline{z}_1).`$ (3.16)
In particular, the patching function $`p(\lambda )`$ turns out to obey the linear differential equations
$`(_{\overline{z}_1}\lambda _{z_2})p(\lambda )=0,(_{\overline{z}_2}+\lambda _{z_1})p(\lambda )=0.`$ (3.17)
(Note that this is a rather simplified setup. If the solution is defined in a general domain of spacetime, we need a more refined cohomological language — see Ivanova’s review and references cited therein.)
Given such a patching function, one can prove that the Riemann-Hilbert problem indeed solves the ASDYM equations. We shall review this proof later on in the framework of the noncommutative spacetime. The converse is also true. Namely, if $`\mathrm{\Psi }(\lambda )`$ and $`\widehat{\mathrm{\Psi }}(\lambda )`$ are a pair of arbitrary solutions of (3.12), its matrix ratio $`\widehat{\mathrm{\Psi }}(\lambda )^1\mathrm{\Psi }(\lambda )`$ satisfies equation (3.17). This can be confirmed by direct calculations.
Solving the Riemann-Hilbert problem explicitly is usually very difficult. Explicit solutions are known for special cases only. The so called “Ward Ansatz” (or “Atiyah-Ward Ansatz”) solutions provide such an example. The corresponding Riemann-Hilbert problem can be solved by linear algebra.
The existence of a large set of hidden symmetries of the ASDYM equations can be explained by the Riemann-Hilbert problem . Those symmetries are generated by the left and right action
$`p(\lambda )g_L(\lambda )p(\lambda )g_R(\lambda )^1`$ (3.18)
of $`GL(N,)`$-valued functions $`g_L(\lambda )`$ and $`g_R(\lambda )`$ of $`(\lambda ,z_1\lambda \overline{z}_2,z_2+\lambda \overline{z}_1)`$. The infinitesimal form of these symmetries can be determined explicitly and reproduces the previously known results . For subsequent progress on finite transformations, see Popov’s paper .
### 3.4 Integrability of $``$-product ASDYM equations
Having reviewed the twistorial and integrable structures of the ordinary ASDYM equations, we now turn to the $``$-product ASDYM equations.
The geometric setup of twistor theory can be extended to the noncommutative spacetime rather straightforward. To see this, let us notice that the commutation relations of the complex coordinates $`(z_1,z_2,\overline{z}_1,\overline{z}_2)`$ can be rewritten
$`[z_1\lambda \overline{z}_2,z_2+\lambda \overline{z}_1]=\zeta _{}\lambda \zeta _{}+\lambda ^2\overline{\zeta }_{}.`$ (3.19)
The linear combinations of the spacetime coordinates on the left hand side are exactly those in the definition of the twistor surface $`S(\lambda ,u_1,u_2)`$. Accordingly, whereas $`\lambda `$ persists to be a commutative coordinate, the coordinates $`u_1`$ and $`u_2`$ of the fibers of $`\pi :𝒯_{}^1`$ turn out to have to obey the commutation relation
$`[u_1,u_2]=\zeta _{}\lambda \zeta _{}+\lambda ^2\overline{\zeta }_{}.`$ (3.20)
Thus the twistor space, like the spacetime, becomes a noncommutative manifold. This will be an alternative interpretation of the results of of Kapustin et al. .
The linear system for $`\mathrm{\Psi }(\lambda )`$ is now replaced by the $``$-product version
$`(_{\overline{z}_1}\lambda _{z_2})\mathrm{\Psi }(\lambda )=(_{\overline{z}_1}\lambda _{z_2})\mathrm{\Psi }(\lambda )+(A_{\overline{z}_1}\lambda A_{z_2})\mathrm{\Psi }(\lambda )=0,`$
$`(_{\overline{z}_2}+\lambda _{z_1})\mathrm{\Psi }(\lambda )=(_{\overline{z}_2}+\lambda _{z_1})\mathrm{\Psi }(\lambda )+(A_{\overline{z}_2}+\lambda A_{z_1})\mathrm{\Psi }(\lambda )=0.`$ (3.21)
Although the notion of vector bundles on the noncommutative twistor space is complicated , the Riemann-Hilbert problem itself remains intact except that the product $`\widehat{\mathrm{\Psi }}(\lambda )^1\mathrm{\Psi }(\lambda )`$ is replaced by the $``$-product:
$`p(\lambda )=(\widehat{\mathrm{\Psi }}(\lambda ))_{}^1\mathrm{\Psi }(\lambda ).`$ (3.22)
The patching function $`p(\lambda )`$ is required to satisfy the same linear differential equations as (3.17), or, equivalently, to be of the form $`P(\lambda ,z_1\lambda \overline{z}_2,z_2+\lambda \overline{z}_1)`$.
Let us confirm that this Riemann-Hilbert problem indeed solves the $``$-product ASDYM equations. The reasoning is fully parallel to the ordinary ASDYM equations. We first note that (3.17) implies the equations
$`(_{\overline{z}_1}\lambda _{z_2})\mathrm{\Psi }(\lambda )(\mathrm{\Psi }(\lambda ))_{}^1`$ $`=`$ $`(_{\overline{z}_1}\lambda _{z_2})\widehat{\mathrm{\Psi }}(\lambda )(\widehat{\mathrm{\Psi }}(\lambda ))_{}^1,`$
$`(_{\overline{z}_2}+\lambda _{z_1})\mathrm{\Psi }(\lambda )(\mathrm{\Psi }(\lambda ))_{}^1`$ $`=`$ $`(_{\overline{z}_2}+\lambda _{z_1})\widehat{\mathrm{\Psi }}(\lambda )(\widehat{\mathrm{\Psi }}(\lambda ))_{}^1.`$ (3.23)
Since $`D`$ and $`\widehat{D}`$ cover the whole Riemann sphere, both hand sides of these equations define a matrix-valued meromorphic function with the only possible poles being at $`\lambda =\mathrm{}`$ and of the first order. By Liouville’s theorem, they are a linear function of $`\lambda `$ with matrix coefficients. Let us express these linear functions as $`A_{\overline{z}_1}+\lambda A_{z_2}`$ and $`A_{\overline{z}_2}\lambda A_{z_1}`$. The coefficients $`A_{z_1},A_{z_2},A_{\overline{z}_1},A_{\overline{z}_2}`$ are to be identified with the gauge potentials. Thus $`\mathrm{\Psi }(\lambda )`$ and $`\widehat{\mathrm{\Psi }}(\lambda )`$ turn out to satisfy (3.4), from which the $``$-product ASDYM equations are derived.
The other part of the forgoing discussion, too, can be mostly extended to the noncommutative spacetime. For instance, hidden symmetries are again generated by the action
$`p(\lambda )g_L(\lambda )p(\lambda )g_R(\lambda )^1`$ (3.24)
of $`GL(N,)`$-valued functions $`g_L(\lambda )`$ and $`g_R(\lambda )`$ of $`(\lambda ,z_1\lambda \overline{z}_2,z_2+\lambda \overline{z}_1)`$. The associated infinitesimal symmetries take the same form as those for the ordinary ASDYM equations (with, of course, the product of spacetime functions being replaced by the $``$-product).
An essential difference can be seen in the places where finite dimensional linear algebra is used. A typical example is the Ward Ansatz. In the noncommutative framework, such a linear algebraic structure has to be replaced by an infinite dimensional counterpart. As for the Ward Ansatz, for instance, we do not know how to extend it to the $``$-product ASDYM equations.
## 4 Deformed ADHM construction
### 4.1 How to deform ADHM construction
The ordinary ADHM construction of a $`U(N)`$-instanton solution with instanton number $`k`$ is based on the $`2k\times (2k+N)`$ matrix-valued function
$`\mathrm{\Delta }(z)=\left(\begin{array}{ccc}B_1+z_11& B_2+z_21& I\\ B_2^{}\overline{z}_21& B_1^{}+\overline{z}_11& J^{}\end{array}\right)`$ (4.3)
of $`z=(z_1,z_2,\overline{z}_1,\overline{z}_2)`$. Here $`B_1`$ and $`B_2`$ are $`k\times k`$ matrices, $`I`$ a $`k\times N`$ matrix, $`J`$ an $`N\times k`$ matrix, and $`B_1^{},B_2^{},I^{},J^{}`$ their Hermitian conjugate. Assuming a nondegeneracy condition, one can construct a $`(2k+N)\times N`$ matrix $`v(z)`$ that satisfies the equations
$`\mathrm{\Delta }(z)v(z)=0,v(z)^{}v(z)=1.`$ (4.4)
If the so called ADHM equations
$`[B_1,B_2]+IJ=0,`$
$`[B_1,B_1^{}]+[B_2,B_2^{}]+II^{}J^{}J=0`$ (4.5)
are satisfied, the gauge potentials defined by
$`A=v(z)^{}dv(z)`$ (4.6)
give a solution (instanton solution) of the ASDYM equations.
As Nekrasov and Schwarz pointed out, the instanton solutions of the ASDYM equations on the noncommutative $`^4`$ can be obtained by deforming the ADHM equations as
$`[B_1,B_2]+IJ=\zeta _{}1,`$
$`[B_1,B_1^{}]+[B_2,B_2^{}]+II^{}J^{}J=\zeta _{}1.`$ (4.7)
The connection form is now given by the $``$-product
$`A=v(z)^{}dv(z).`$ (4.8)
### 4.2 Solution of Riemann-Hilbert problem
We here present, as an application of the ADHM construction, an explicit construction of the solution of the Riemann-Hilbert problem for the instanton solutions. This is based on the work of Corrigan et al. on the Dirac equation with the instanton gauge potentials
According to one of their results, the parallel translation (i.e., the “Wilson line operator”)
$`w(z,z^{})=\text{P-exp}\left({\displaystyle _z^{}^z}A\right)`$ (4.9)
between two points $`z,z^{}`$ on the same twistor surface $`S(\lambda ,u_1,u_2)`$ is given by the simple formula
$`w(z,z^{})=v(z)^{}v(z^{}).`$ (4.10)
Consequently, this matrix obeys the group law
$`w(z,z^{})w(z^{},z^{\prime \prime })=w(z,z^{\prime \prime })`$ (4.11)
for any triple $`z,z^{},z^{\prime \prime }`$ of points on $`S(\lambda ,u_1,u_2)`$.
Let us apply this group law to the special points
$`z^{\mathrm{}}(\lambda )`$ $`=`$ $`(z_1\lambda \overline{z}_2,z_2+\lambda \overline{z}_1,0,0),`$
$`z^0(\lambda )`$ $`=`$ $`(0,0,\overline{z}_1+\lambda ^1z_2,\overline{z}_2\lambda ^1z_1)`$ (4.12)
that are on the same twistor surface as $`z=(z_1,z_2,\overline{z}_1,\overline{z}_2)`$. Acccordingly, we have the relation
$`w(z^{\mathrm{}}(\lambda ),z^0(\lambda ))=w(z,z^{\mathrm{}}(\lambda ))^1w(z,z^0(\lambda )).`$ (4.13)
This relation is exactly the Riemann-Hilbert problem with the patching function
$`p(\lambda )=w(z^{\mathrm{}}(\lambda ),z^0(\lambda )),`$ (4.14)
for which we thus obtain the explicit solution
$`\mathrm{\Psi }(\lambda )=w(z,z^0(\lambda )),\widehat{\mathrm{\Psi }}(\lambda )=w(z,z^{\mathrm{}}(\lambda )).`$ (4.15)
This construction carries over to the noncommutative case if the ordinary matrix products therein are replaced by the $``$-product. The parallel translation along the twistor surface is given by the $``$-product
$`w(z,z^{})=v(z)^{}v(z^{}),`$ (4.16)
and the foregoing expressions of $`g(\lambda )`$, $`\mathrm{\Psi }(\lambda )`$ and $`\widehat{\mathrm{\Psi }}(\lambda )`$ remain valid.
### 4.3 Remarks on ADHM equations
It is well known that the left hand side of the ADHM equations, i.e.,
$`\mu _{}`$ $`=`$ $`[B_1,B_2]+IJ,`$
$`\mu _{}`$ $`=`$ $`[B_1,B_1^{}]+[B_2,B_2^{}]+II^{}J^{}J,`$ (4.17)
are a pair of moment maps for the hyper-Kähler quotient construction of the moduli space of (both undeformed and deformed) ADHM instantons. In particular, the pair $`(\mu _{},\mu _{})`$ transforms just like the three-vector $`(\zeta _{},\zeta _{})`$ under the $`SU(2)`$ rotation (2.10) of spacetime coordinates. Therefore it is natural to combine the three moment maps $`\mu _{},\mu _{},\mu _{}^{}`$ to the one-parameter family
$`\mu (\lambda )=\mu _{}+\lambda \mu _{}\lambda ^2\mu _{}^{}=[B_1\lambda B_2^{},B_2+\lambda B_1^{}]+(I\lambda J^{})(J+\lambda I^{})`$ (4.18)
of moment maps. The $`SU(2)`$ action is now represented by fractional transformations of $`\lambda `$:
$`\lambda {\displaystyle \frac{\beta +\alpha \lambda }{\overline{\alpha }+\overline{\beta }\lambda }}.`$ (4.19)
A “pencil” of moment maps of this type generally appears in the quotient construction of the twistor space associated with a hyper-Kähler quotient . The twistor space $`𝒵`$ of a hyper-Kähler manifold is fibered over $`_{}^1`$ by a map $`\pi :𝒵_{}^1`$, and each fiber $`\pi ^1(\lambda )`$ is a complex symplectic manifold. The moment map $`\mu (\lambda )`$ is used to make the symplectic quotient of $`\pi ^1(\lambda )`$. Roughly speaking, this fiberwise symplectic quotient of $`𝒵`$ gives the twistor space for the hyper-Kähler quotient.
This pencil of moment maps is also interesting in the context of finite dimensional integrable systems. Following Gorsky, Nekrasov and Rubtsov , let us introduce the symplectic form
$`\mathrm{\Omega }=Tr(dB_1dB_2+dIdJ)`$ (4.20)
on the space of the quadruples $`(B_1,B_2,I,J)`$. As they pointed out, $`\mu _{}`$ may be interpreted as the moment map of the action
$`(B_1,B_2,I,J)(gB_1g^1,gB_2g^1,gI,Jg^1)`$ (4.21)
of $`G=GL(N,)`$, and the reduced phase space (actually, with $`I`$ and $`J`$ being further constrained to a special $`G`$-orbit) has the structure of an integrable system with the Poisson-commutative Hamiltonians $`TrB_2^{\mathrm{}}`$, $`\mathrm{}=1,\mathrm{},N`$. If $`k`$ is equal to $`1`$ and $`B_1`$ and $`B_2`$ are restricted to Hermitian matrices, this integrable system reduces to the rational Calogero-Moser system; the case for $`k>1`$ is relate to a generalized Calogero-Moser system . Now, what occurs if one repeats the same construction for the pencil $`\mu (\lambda )`$ of moment maps? Note that the symplectic form, too, has to be deformed as
$`\mathrm{\Omega }(\lambda )=Tr\left(d(B_1\lambda B_2^{})d(B_2+\lambda B_1^{})+d(I\lambda J^{})d(J+\lambda I^{})\right).`$ (4.22)
Upon taking the symplectic quotient, a one-parameter family of integrable systems will emerge. In fact, $`\mathrm{\Omega }(\lambda )`$ is exactly the symplectic form of the fiber $`\pi ^1(\lambda )`$ of the twistor space before taking the quotient. Thus the phase space of the aforementioned one-parameter family of integrable systems turns out to be nothing but the fibers $`\pi ^1(\lambda )`$ of the twistor space of the instanton moduli space.
## 5 Two-dimensional reductions
### 5.1 PCF model and Hitchin’s equations
We here examine the PCF model and Hitchin’s Higgs pair equations as dimensional reductions of the ASDYM equations.
The PCF model can be derived by letting
$`_{\overline{z}_1}_z+A_z,`$ $`_{z_2}_z,`$
$`_{\overline{z}_2}_{\overline{z}}+A_{\overline{z}},`$ $`_{z_1}_{\overline{z}}`$ (5.1)
under the gauge $`A_{z_1}=A_{z_2}=0`$. The associated linear system reads
$`((1\lambda )_zA_z)\mathrm{\Psi }(\lambda )=0,((1+\lambda )_{\overline{z}}+A_{\overline{z}})\mathrm{\Psi }(\lambda )=0.`$ (5.2)
On the noncommutative spacetime, $`z`$ and $`\overline{z}`$ are assumed to obey the commutation relation
$`[z,\overline{z}]=\zeta 1`$ (5.3)
for a real constant $`\zeta `$, and the linear system is replaced by the $``$-product analogue
$`((1\lambda )_zA_z)\mathrm{\Psi }(\lambda )=0,((1+\lambda )_{\overline{z}}+A_{\overline{z}})\mathrm{\Psi }(\lambda )=0.`$ (5.4)
Conservation lows, infinitesimal symmetries, the Riemann-Hilbert problem, etc. can be extended to the $``$-product PCF model straightforward.
Hitchin’s Higgs pair equations
$`F_{z\overline{z}}=[\mathrm{\Phi },\mathrm{\Phi }^{}],_{\overline{z}}\mathrm{\Phi }=0,_z\mathrm{\Phi }^{}=0`$ (5.5)
can be derived from the ASDYM equations by first exchanging $`z_2\overline{z}_2`$ (which interchanges anti-self-duality and self-duality), then reducing
$`_{\overline{z}_2}_z,`$ $`_{z_1}\mathrm{\Phi },`$
$`_{z_2}_{\overline{z}},`$ $`_{\overline{z}_1}\mathrm{\Phi }^{}`$ (5.6)
while letting $`_{z_1}0`$ and $`_{\overline{z}_1}0`$. The associated linear system can be written
$`(_z+\lambda \mathrm{\Phi })\mathrm{\Psi }(\lambda )=0,(_{\overline{z}}\lambda ^1\mathrm{\Phi }^{})\mathrm{\Psi }(\lambda )=0.`$ (5.7)
A natural $``$-product analogue of these equations are, of course,
$`F_{z\overline{z}}=[\mathrm{\Phi },\mathrm{\Phi }^{}]_{},_{\overline{z}}\mathrm{\Phi }=0,_z\mathrm{\Phi }^{}=0`$ (5.8)
and
$`(_z+\lambda \mathrm{\Phi })\mathrm{\Psi }(\lambda )=0,(_{\overline{z}}\lambda ^1\mathrm{\Phi }^{})\mathrm{\Psi }(\lambda )=0.`$ (5.9)
### 5.2 Some more remarks on Hitchin’s equations
Hitchin’s equations are formulated on any compact Riemann surface . If the genus of the Riemann surface is greater than one, the moduli space of “stable” Higgs pairs is a smooth (but noncompact) symplectic manifold with the structure of an “algebraically integrable Hamiltonian system” . (This fact is further extended to punctured Riemann surfaces including tori.) The “spectral curve” $`det(\mathrm{\Phi }\zeta 1)=0`$ plays a central role therein.
What about the $``$-product analogue of Hitchin’s equations? Unfortunately, we do not know if the moduli space of solutions has any structure of an integrable system, because, first of all, the notion of determinant (hence, of spectral curve) ceases to exist. This is a place where a linear algebraic structure breaks down again. One can nevertheless expect that some yet unknown mechanism might give rise to an integrable structure in the moduli space solutions. This issue will be closely related to the notion of “noncommutative Riemann surfaces” that has been pursued by Bertoldi et al. .
Let us finally mention that Hitchin’s equations are also related to a class of conformal field theories — e.g., the (non-affine) Toda field theories and $`W`$-gravity . The associated $``$-product analogues will be interesting from the point of view of the Chern-Simons and WZW models on noncommutative spaces . Note, however, that the naive substitution prescription $`e^{\alpha \varphi }(e^{\alpha \varphi })_{}`$ in the Toda field theories does not lead to an integrable system. A correct integrable deformation is the so called “nonabelian Toda field theory”, which does not take such an exponential form.
## 6 Conclusion
We have shown that many properties of the ASDYM equations are inherited by its analogue on the noncommutative $`^4`$. After all, the rule of game is quite simple — just to replace the ordinary product by the $``$-product. This rather naive prescription has turned out to fit surprisingly well into the twistorial and integrable structures of the ASDYM equations. Moreover, these structures are preserved under dimensional reduction to the PCF model and Hitchin’s Higgs pair equations. However, linear algebraic structures, such as the Ward Ansatz solutions, mostly loose its meaning in the noncommutative spacetimes.
We have also pointed out a few interesting structures in the ADHM construction. These structures deserve to be studied in more detail.
Another important issue, which we have not addressed in this paper, is that of the Nahm equations. The Nahm construction of BPS monopoles has been extended to a noncommutative spacetime , in which a $``$-product analogue of the Nahm equations is used. The $``$-product Nahm equations have been independently studied in the context of the M-theory as well .
### Acknowledgements
I am grateful to David Fairlie for useful comments.
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# Kinematic Age Estimates for 4 Compact Symmetric Objects from the Pearson-Readhead Survey
## 1 Introduction
While most radio galaxies exhibit an asymmetric, core-jet morphology on the parsec scale, approximately 7% of radio galaxies in complete, flux-limited samples display parsec-scale jets and lobe emission on both sides of a central engine (Taylor et al. 1996). This emission is thought to be free of Doppler boosting effects (Wilkinson et al. 1994) and the sources are believed to be physically small as opposed to appearing small due to projection.
The measurement of a kinematic age for the CSO 0710+439 by Owsianik & Conway (1998) of just 1100 $`\pm `$ 100 years lent strong support to the theory that CSOs are small by virtue of their youth and not because of confinement. This was the favored interpretation by Phillips & Mutel (1980, 1982) who first drew attention to a group of 4 compact double sources with steep spectra and slow motions compared to the majority of core-jet sources. Although further multi-frequency VLBI observations have revealed that 2 of these (CTD93 - Shaffer, Kellermann & Cornwell 1999; and 3C 395 - Taylor 2000) are actually asymmetric core-jet sources, many of the speculations of Phillips & Mutel have been borne out. Their misidentification of two sources emphasizes the fact that sensitive multi-frequency VLBI observations are required to demonstrate symmetric structure on the parsec scale. In particular, not all (or even most) GHz Peaked Spectrum (GPS) radio galaxies belong to the CSO class. Failure to find any evidence of the extremely dense environment required to confine CSOs, along with some indirect age measurements, also led Readhead et al. (1996a) to favor the idea that CSOs are young. In fact, angular separation rates for the CSOs measured to date indicate typical hotspot velocities of $``$0.1 $`h^1`$ c (e.g., Owsianik et al. 1998, Owsianik & Conway 1998, Polatidis et al. 1999). Such speeds are about an order of magnitude larger than predicted by Readhead et al. (1996a) and so indicate even younger ages.
The evolution of CSOs is of considerable current interest. To understand the evolution of radio galaxies we need to be able to recognize the evolutionary state of any given galaxy. Readhead et al. (1996b) have proposed an evolutionary model for powerful radio sources in which CSO’s evolve first into Compact Steep Spectrum doubles and then into large Fanaroff-Riley (1974) Type II objects. Based on a different sample of somewhat larger objects, Fanti et al. (1995) reached the same conclusions. Measurements of source ages are crucially important to the understanding of the evolution of this intriguing class of objects.
In 1994 Taylor, Readhead, and Pearson (1996) performed VLBA observations of 4 CSOs from the Pearson-Readhead survey (Pearson & Readhead 1988) at 15 GHz in order to pinpoint their centers of activity. The results were surprising in that the components inferred to be the cores, because of strongly inverted spectra and compactness, were found not to be associated with the strongest central components seen at 5 GHz and lower frequencies, but instead were unresolved components very close to the midway point between the hotspots (e.g., Fig. 1). The bright components seen at 5 GHz (in some cases previously mis-identified as cores) turn out to be one-sided jets. To further characterize the properties of the core components we carried out 43 GHz VLBA observations of several PR CSOs in 1996 and 1999.
In §3 we present multi-epoch VLBA observations at 15 and 43 GHz for four CSOs from the PR survey. These results are used in §4 to measure the advance speeds of the hot spots, and to thereby determine kinematic age estimates. Velocity measurements are expressed in terms of $`h`$ = H<sub>0</sub> / 100 km s<sup>-1</sup> Mpc<sup>-1</sup>, and where physical scales are quoted (or drawn) we assume $`h=0.65`$.
## 2 Observations and Data Reduction
Observations were made at 15 and 43 GHz using the Very Long Baseline Array (VLBA)<sup>1</sup><sup>1</sup>1The National Radio Astronomy Observatory is a facility of the National Science Foundation operated under cooperative agreement by Associated Universities, Inc. telescope at multiple epochs between 1994 and 1999. The details of the observations are provided in Table 1. All sources were observed with 5–8 scans spread across a wide range in hour angle in order to obtain good ($`u`$, $`v`$) coverage. The VLBA correlator produced typically 16 frequency channels across every 8 MHz of observing bandwidth in each 2 second integration period.
Calibration procedures were followed for the 15 GHz data in a manner similar to that used for the first epoch observations as described in Taylor, Readhead & Pearson (1996). At 43 GHz the delays and IF phase offsets were determined from the injected pulse-calibration and a short observation of the strong calibrator 3C 84. The data were then averaged in frequency to a single channel of 64 MHz. No global fringe-fitting was attempted as the scatter in the solutions so derived was judged to be worse than the small residual delay error. After phase self-calibration with a 10 s solution interval and a point-source model, the data were coherently averaged to 10 s integrations. All editing, imaging, deconvolution, and self-calibration were performed using Difmap (Shepherd, Pearson & Taylor 1994, 1995). Several iterations of phase self-calibration and imaging were performed with each data set before any attempt at amplitude self-calibration, and no amplitude self-calibration was attempted at 43 GHz. No reliable 43 GHz image could be obtained for the weakest source, 0108+388. At each iteration, windows for clean components were added, if necessary, to provide support and reject sidelobes.
Once the data were completely self-calibrated, Gaussian model-fitting was performed using Difmap. To determine relative motions the component shapes and sizes were frozen in the model-fitting to equal the fitted values in the first epoch, while the flux densities and positions of components were allowed to vary. The AIPS task JMFIT was also used to perform Gaussian model-fitting to the hotspots in the image plane. In all cases the AIPS JMFIT produced very similar results to fitting the visibility data.
## 3 Results
A history of the VLBI observations of each source studied herein can be found in Taylor, Readhead & Pearson (1996; hereafter TRP96). For convenience source identifications and redshifts are given in Table 2, and measurements of core properties are summarized in Table 3. Here we describe our velocity measurements and compare them with recent results in the literature. We compute the kinematic age of the source in its rest frame, $`\tau _k`$, as $`\tau _k=\theta _{\mathrm{hs}}/\mu _{\mathrm{hs}}(1+z)`$ where $`\theta _{\mathrm{hs}}`$ is the angular separation of the hot spots, and $`\mu _{\mathrm{hs}}`$ is the angular separation rate of the hot spots.
### 3.1 0108+388
In Fig. 1 we show a naturally weighted 15 GHz image of 0108+388. This image shows continuous emission connecting the two main components (labeled C1 through C7). Component C3 was identified as the core based on its inverted spectrum and compactness (TRP96; and see Table 3). A detailed spectral analysis between 1.6 and 15 GHz has been carried out by Marr, Taylor, and Crawford (2000a,b) who find that while C3 is likely to have a spectral turnover at a high frequency due to synchrotron self-absorption, the spectra of other components turn over around 5 GHz because of free-free absorption by an ionized disk centered on the nucleus. More evidence for a relatively dense circumnuclear environment comes from Hi absorption measurements by Carilli et al. (1998) who find an optical depth of 0.44 $`\pm `$ 0.04 and implied column density of 80.7 $`\times `$ 10<sup>18</sup> T$`{}_{s}{}^{}/f`$ cm<sup>-2</sup>, where T<sub>s</sub> is the spin temperature and $`f`$ is the Hi covering factor. Owsianik, Conway & Polatidis (1998) measured an angular separation rate for the outer components (C1 and C7) of 9.3 $`\pm `$ 1.2 $`\mu `$as yr<sup>-1</sup> from 3 epochs of global VLBI observations at 5 GHz spread over 12 years.
Given the identification of C3 as the core in 0108+388, the relative velocity of 0.57 $`\pm `$ 0.35 $`h^1`$c (see Fig. 1 and Table 4) is at first glance somewhat alarming since one traditionally assumes the core component to be stationary. The core component, however, is likely to be just the optically thick base of the jet (Blandford & Königl 1979), and could well appear to move in either direction along the jet axis as new jet components emerge. In one of the few cases where absolute motions have been obtained there is evidence that the “core” of 1928+738 moves in this way (Ros et al. 1999). For this reason, the core component makes a poor choice as a reference feature. Instead we adopt the westernmost component (C1) as the reference feature. As we will argue, this is likely to be a subrelativistic hotspot. This choice is somewhat arbitrary, and we could just as well have taken the easternmost hot spot (C7) as a reference.
From our 3 VLBA epochs at 15 GHz we find a separation rate for C1 and C7 of 11 $`\pm `$ 2 $`\mu `$as yr<sup>-1</sup>, which is in agreement with the measurements of Owsianik et al. (1998) and also with recent measurements at 8.4 GHz (Polatidis et al. 1999). Our kinematic age estimate for 0108+388 is 310 $`\pm `$ 70 yrs. Assuming that C1 and C7 are moving apart at equal speeds, this gives each an advance speed of 0.12 $`h^1`$c. The best-fit models from which these velocities were derived are listed in Table 4. A significant velocity of 0.79 $`\pm `$ 0.04 $`h^1`$c is found for component C5, although since this is relative to C1 the true velocity of C5 likely to be smaller by 0.12 $`h^1`$c.
Due to the low flux density of 0108+388 at 43 GHz, it was not possible to self-calibrate the data and make an image. Even so we can still estimate that the core flux density must be less than $``$40 mJy at 43 GHz, otherwise it would have been readily detected. This indicates that the spectral index is not as steeply inverted as at the lower frequencies (see Table 3).
### 3.2 0710+439
Owsianik & Conway (1998) reported the first significant detection of a hot spot advance speed based on 5 epochs on 0710+439 at 5 GHz between 1980 and 1993. The hot spots (components A2 and C2) were found to have a separation rate of 14 $`\pm `$ 1.6 $`\mu `$as yr<sup>-1</sup>.
In Fig. 2 we show our 2nd epoch 15 GHz observations with motions derived from the two epochs indicated by arrows. Components A and C are both leading-edge brightened with emission fading gradually towards the center of the source. Component B is more compact at the southern end and becomes wider to the north with an opening angle of $``$20 degrees. At the base of the jet a compact inverted spectrum component (B5 in Fig. 2) was identified and inferred to be the core by TRP96. Gaussian model-fits and component motions are given in Table 5. In general our 15 GHz model components correspond to the 5 GHz model of Owsianik & Conway (1998), but at 15 GHz the extended components A1, B1 and C1 are resolved out. The other difference in the 5 and 15 GHz models is the core component (B5) which was too weak to be included in the 5 GHz model.
We find a separation rate between the hotspots (A2 and C2) of 29 $`\pm `$ 8.7 $`\mu `$as yr<sup>-1</sup>. The implied velocity of advance, assuming equal speeds, is 0.26 $`h^1`$c. This velocity is nearly twice that found by Owsianik & Conway (1998). The kinematic age we derive for 0710+439 is 550 $`\pm `$ 160 yrs. Since the epochs do not overlap, one explanation could be that the hot spot advance speed has recently doubled. A more likely explanation for the discrepancy is that the 15 GHz observations are more sensitive to the motion of a bright, compact working surface, while the 5 GHz observations measure the more stable overall expansion of the lobe. In support of this idea we note that the size of component A2 is 0.51 $`\times `$ 0.33 mas at 15 GHz, and 0.84 $`\times `$ 0.57 mas at 5 GHz.
We also find a substantial velocity of 1.36 $`\pm `$ 0.16 $`h^1`$c for component B3 relative to C2. This velocity, however, is accompanied by a change in the flux density of B3 from 116 mJy in 1994.971 to only 63 mJy in 1999.587. The large drop in flux could indicate a region or subcomponent that faded significantly, causing a large shift in the centroid.
In Fig. 3 we show an image of 0710+439 at 43 GHz. Only the core, inner jet complex, and brightest region of the northern hot spot are detected. The hot spot is well resolved in these observations and shows a similar twist to the southeast as indicated in the 15 GHz image. The core component is prominent and still unresolved. The core has a 15-43 GHz spectral index of $``$0.2 $`\pm `$ 0.3 (where $`S_\nu \nu ^\alpha `$). Despite the higher resolution afforded by the 43 GHz observations, no significant component motions were detected due to the shorter time baseline and lower signal-to-noise compared to the 15 GHz observations.
### 3.3 1031+567
A map made from our 15 GHz observations of 1031+567 is shown in Fig. 4 and the modelfit components are listed in Table 6. No central core component is apparent. It is likely that the two outermost components are the working surfaces and lobes of two oppositely directed jets, given the edge brightened appearance and steep spectra (TRP96) and that no core emission is detected. This source is also remarkable for the very small (1 – 9%) change in the flux density of its components.
We report here a first tentative detection of the expansion of 1031+567. We find a separation rate between the hot spots W1 and E1 of 14.6 $`\pm `$ 4.8 $`\mu `$as yr<sup>-1</sup>, although not along the axis of the source (see Fig. 4). A larger expansion rate is found between W1 and E2 of 37.6 $`\pm `$ 8.4 $`\mu `$as yr<sup>-1</sup> roughly along the position angle of the axis. The kinematic age derived from the W1–E2 expansion rate is 620 $`\pm `$ 140 yrs. Assuming equal advance speeds this corresponds to a velocity of 0.31 $`h^1`$c.
### 3.4 2352+495
The hot spots (labeled A and C after Conway et al. 1992) are well resolved in our 15 GHz image (Fig. 5). We also see a faint narrow jet to the south of the bright B1-B5 complex with a compact component (D) embedded in it and identified by TRP96 as the core. Based on global VLBI observations at 5 GHz over 14 years, Owsianik, Conway & Polatidis (1999) measured a relative angular separation rate for the hot spots of 21.1 $`\pm `$ 2.7 $`\mu `$as yr<sup>-1</sup>.
The model-fit results at 15 GHz for 2352+495 are listed in Table 7. From our two epoch VLBA observations over 4.6 years we find an angular separation velocity for the hot spots of 33 $`\pm `$ 11 $`\mu `$as yr<sup>-1</sup>. Given our large uncertainty, these observations are consistent with the findings of Owsianik et al. (1999). Our measured component separation rate yields a kinematic age for 2352+495 of 1200 $`\pm `$ 400 years. Assuming equal advance speeds for the hot spots yields a velocity of 0.16 $`h^1`$c. The jet components in the component B complex have velocities ranging from 0.27 – 0.76 $`h^1`$c.
In Fig. 6. we show a 43 GHz image of 2352+495. Only the core and B1-B5 complex is detected. The core component appears well isolated and still unresolved, and has a 15-43 GHz spectral index of $``$0.6 $`\pm `$ 0.4.
## 4 Discussion
### 4.1 Baby Radio Galaxies?
The kinematic ages derived for the 4 bright CSOs studied here range from 300 to 1200 yrs. This is part of a growing body of evidence conclusively demonstrating that CSOs are indeed young objects. These ages are derived under the assumption that the velocity has been constant over the entire period of activity. The kpc-scale bridge of emission in 0108+388 (Baum et al. 1990) suggests prior active phases in at least this CSO. Owsianik et al. (1998) have suggested that the one-sided appearance of the large scale structure is due to this periodic nature and light travel time effects.
In 3 of the 4 CSOs observed here the hot spots appear to be advancing along the jet axis. The one exception is 1031+567 for which the separation of the hot spot heads is not along the axis of the source, although the lobe complex is moving out along the source axis. In the dentist drill model (Scheuer 1974) for radio galaxy evolution, the hotspots wander around the leading edge of the lobe. At lower frequencies (see for example Owsianik et al. 1999) the northern lobe of 2352+495 appears to extend further from the core than the hotspot. There is also a sharp bend at the hotspot so it is possible that the jet is deflected at the “primary” hot spot toward the east.
### 4.2 Core and Jet Properties
Three of the four CSO’s studied here contain central components inferred to be the cores. All three of these apparent core components display spectra that peak around 15 GHz (see Table 3) and are unresolved. If the turnovers are due to synchrotron self-absorption, following Marscher (1983) one finds that the magnetic fields in the cores must be less than 7 $`\times 10^4`$, 80, and 600 Gauss in 0108+388, 0710+439, and 2352+495, respectively. Free-free absorption is unlikely to be significant at 15 GHz, although it has been demonstrated to be significant at frequencies below 5 GHz in the central regions of some CSOs (e.g., 1946+708 – Peck, Taylor, & Conway 1999; 0108+388 – Marr, Taylor, & Crawford 2000a,b). One might also expect the cores to turn over more sharply than is indicated in Table 3 if free-free absorption is the dominant mechanism.
In 0108+388 we find marginal evidence for motion of the core component. Less significant motions are detected in 0710+439 and 2352+495. These motions are expected in the Blandford & Königl (1979) model where the core is the base of an optically thick jet as new components emerge from the center of activity.
Jet components moving relativistically from the cores towards the hot spots are seen in all sources studied here except 1031+567. In 0710+439 and 2352+495 the jet components appear stronger on one side of the core, consistent with their fast motions and probably indicating that these sources do not lie in the plane of the sky. Without some additional constraint on jet velocities, however, the inclination angle is not well determined.
## 5 Conclusions
We confirm hot spot advance speeds of $``$0.2 $`h^1`$c in three CSOs and present a new detection for the CSO 1031+567. These growth rates correspond to ages between 300 and 1200 yrs for the current phase of activity in the sources studied here. Further observations should allow for a more precise determination of the hot spot velocities and will test the dentist drill model for radio sources.
We thank the referee, Ken Kellermann, for insightful comments on the manuscript. This research has made use of the NASA/IPAC Extragalactic Database (NED) which is operated by the Jet Propulsion Laboratory, Caltech, under contract with NASA. This research has made use of data from the University of Michigan Radio Astronomy Observatory which is supported by the NSF and by funds from the University of Michigan.
TABLE 1
Observational Parameters
| Source | Date | Frequency | Bandwidth<sup>a</sup> | Scan Length | Total Time |
| --- | --- | --- | --- | --- | --- |
| | | (GHz) | (MHz) | (min) | (hours) |
| 0108+388 | 1994.971 | 15.4 | 16 | 26 | 2.2 |
| | 1996.594 | 43.2 | 64 | 5 | 1.3 |
| | 1997.092 | 15.4 | 64 | 5 | 4.3 |
| | 1999.587 | 15.4 | 64 | 5 | 1.2 |
| | 1999.587 | 43.2 | 64 | 7 | 1.6 |
| 0710+439 | 1994.971 | 15.4 | 16 | 13 | 1.3 |
| | 1996.594 | 43.2 | 64 | 5 | 1.3 |
| | 1999.587 | 15.4 | 64 | 5 | 1.2 |
| | 1999.587 | 43.2 | 64 | 7 | 1.5 |
| 1031+567 | 1995.411 | 15.4 | 16 | 13 | 1.4 |
| | 1999.587 | 15.4 | 64 | 5 | 1.0 |
| 2352+495 | 1994.971 | 15.4 | 16 | 26 | 3.5 |
| | 1994.971 | 43.2 | 64 | 6 | 1.6 |
| | 1999.587 | 15.4 | 64 | 5 | 1.2 |
<sup>a</sup> Total bandwidth refers to the sum of RCP and LCP bandwidths when both were observed in 1994.971.
TABLE 2
Source Identifications
| Source | R.A. | Declination | $`V`$ | ID | $`z`$ |
| --- | --- | --- | --- | --- | --- |
| (1) | (2) | (3) | (4) | (5) | (6) |
| 0108+388 | 01 11 37.3118 | 39 06 28.110 | 22 | G | 0.6703 |
| 0710+439 | 07 13 38.1766 | 43 49 17.005 | 21 | G | 0.518 |
| 1031+567 | 10 35 07.0451 | 56 28 46.697 | 20 | G | 0.4597 |
| 2352+495 | 23 55 09.4539 | 49 50 08.362 | 20 | G | 0.237 |
Notes – Col.(1): B1950 Source name according to IAU convention. Cols.(2)-(3):J2000.0 right ascension and declination (Patnaik et al. 1992). Col.(4): Visual magnitude (Pearson & Readhead 1988). Col.(5): Optical identification: G, galaxy. Col.(6): Redshift (Lawrence et al. 1996).
TABLE 3
Core Properties
| Source | $`S_\mathrm{c}`$(8 GHz) | $`S_\mathrm{c}`$(15 GHz) | $`S_\mathrm{c}`$(43 GHz) | $`\alpha _\mathrm{c}`$(8-15) | $`\alpha _\mathrm{c}`$(15-43) |
| --- | --- | --- | --- | --- | --- |
| (1) | (2) | (3) | (4) | (5) | (6) |
| 0108+388 | 6.3 $`\pm `$ 1 | 14 $`\pm `$ 2 | $`<`$40 | 1.3 $`\pm `$ 0.4 | $`<`$1.0 |
| 0710+439 | 42 $`\pm `$ 4 | 45 $`\pm `$ 5 | 37 $`\pm `$ 7 | 0.1 $`\pm `$ 0.3 | $`0.2`$ $`\pm `$ 0.3 |
| 1031+567 | $`<`$3 | $`<2`$ | | | |
| 2352+495 | 13 $`\pm `$ 3 | 12 $`\pm `$ 3 | 6.5 $`\pm `$ 1 | $``$0.1 $`\pm `$ 0.7 | $`0.6`$ $`\pm `$ 0.4 |
Notes – Col.(1): Source name. Col.(2): Flux density of the core in mJy at 8.4 GHz from Xu (1994). Col.(3): Flux density of core in mJy at 15 GHz. Col.(4): Flux density of core in mJy at 43 GHz. Col.(5): Spectral index of core at between 8.4 and 15 GHz. Col.(6): Spectral index of core at between 15 and 43 GHz.
TABLE 4
Gaussian Model and Relative Proper Motions for 0108+388
| Component | Epoch | $`S`$ | $`r`$ | $`\theta `$ | $`a`$ | $`b/a`$ | $`\mathrm{\Phi }`$ | $`\mu `$ | $`v`$ | p.a. |
| --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- |
| | | (Jy) | (mas) | () | (mas) | | () | ($`\mu `$as/yr) | ($`h^1`$ c) | () |
| C1… | 1994.971 | 0.119 | 0.0 | 0.0 | 0.53 | 0.74 | $``$30.5 | | | |
| | 1997.092 | 0.115 | 0.0 | 0.0 | 0.53 | 0.74 | $``$30.5 | | | |
| | 1999.587 | 0.105 | 0.0 | 0.0 | 0.53 | 0.74 | $``$30.5 | reference | | |
| C2… | 1994.971 | 0.033 | 0.894 | 91.85 | 0.82 | 0.72 | 66.4 | | | |
| | 1997.092 | 0.027 | 0.931 | 92.42 | 0.82 | 0.72 | 66.4 | | | |
| | 1999.587 | 0.023 | 0.963 | 92.90 | 0.82 | 0.72 | 66.4 | 16 $`\pm `$ 20 | 0.35 $`\pm `$ 0.44 | 106 |
| C3… | 1994.971 | 0.010 | 2.510 | 55.98 | 0.29 | 1.0 | 0.0 | | | |
| | 1997.092 | 0.014 | 2.507 | 55.87 | 0.29 | 1.0 | 0.0 | | | |
| | 1999.587 | 0.012 | 2.615 | 55.59 | 0.29 | 1.0 | 0.0 | 26 $`\pm `$ 16 | 0.57 $`\pm `$ 0.35 | 48 |
| C4… | 1994.971 | 0.037 | 3.664 | 54.76 | 0.89 | 0.20 | 74.7 | | | |
| | 1997.092 | 0.030 | 3.750 | 54.58 | 0.89 | 0.20 | 74.7 | | | |
| | 1999.587 | 0.024 | 3.664 | 54.26 | 0.89 | 0.20 | 74.7 | 23 $`\pm `$ 20 | 0.51 $`\pm `$ 0.44 | 44 |
| C5… | 1994.971 | 0.054 | 4.022 | 53.04 | 0.17 | 0.28 | $``$76.6 | | | |
| | 1997.092 | 0.061 | 4.097 | 53.60 | 0.17 | 0.28 | $``$76.6 | | | |
| | 1999.587 | 0.069 | 4.170 | 53.99 | 0.17 | 0.28 | $``$76.6 | 36 $`\pm `$ 2 | 0.79 $`\pm `$ 0.04 | 79 |
| C6… | 1994.971 | 0.104 | 4.856 | 57.68 | 0.93 | 0.37 | $``$84.8 | | | |
| | 1997.092 | 0.089 | 4.859 | 57.68 | 0.93 | 0.37 | $``$84.8 | | | |
| | 1997.092 | 0.083 | 4.844 | 57.44 | 0.93 | 0.37 | $``$84.8 | 4.7 $`\pm `$ 6 | 0.10 $`\pm `$ 0.12 | $``$67 |
| C7… | 1994.971 | 0.172 | 5.830 | 59.84 | 0.47 | 0.48 | $``$79.3 | | | |
| | 1997.092 | 0.152 | 5.841 | 59.88 | 0.47 | 0.48 | $``$79.3 | | | |
| | 1999.587 | 0.130 | 5.878 | 59.98 | 0.47 | 0.48 | $``$79.3 | 11 $`\pm `$ 2 | 0.24 $`\pm `$ 0.04 | 76 |
Notes to Table 4
NOTE – Parameters of each Gaussian component of the model brightness distribution: $`S`$, flux density; $`r,\theta `$, polar coordinates of the center of the component relative to an arbitrary origin, with polar angle measured from north through east; $`a,b`$, major and minor axes of the FWHM contour; $`\mathrm{\Phi }`$, position angle of the major axis measured from north through east; $`\mu `$, relative proper motion of the component; $`v`$, relative projected velocity in units of $`h^1`$ c ($`h`$ = H<sub>0</sub>/100 km s<sup>-1</sup> Mpc<sup>-1</sup>), along the given position angle (p.a.).
TABLE 5
Gaussian Model and Relative Proper Motions for 0710+439
| Component | Epoch | $`S`$ | $`r`$ | $`\theta `$ | $`a`$ | $`b/a`$ | $`\mathrm{\Phi }`$ | $`\mu `$ | $`v`$ | p.a. |
| --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- |
| | | (Jy) | (mas) | () | (mas) | | () | ($`\mu `$as/yr) | ($`h^1`$ c) | () |
| C2… | 1994.971 | 0.051 | 0.0 | 0.0 | 1.30 | 0.62 | $``$7.3 | | | |
| | 1999.587 | 0.043 | 0.0 | 0.0 | 1.30 | 0.62 | $``$7.3 | reference | | |
| B5… | 1994.971 | 0.042 | 12.941 | 2.17 | 0.56 | 0.33 | $``$13.1 | | | |
| | 1999.587 | 0.041 | 12.974 | 2.29 | 0.56 | 0.33 | $``$13.1 | 9.3 $`\pm `$ 9.7 | 0.17 $`\pm `$ 0.17 | 42 |
| B4… | 1994.971 | 0.121 | 14.205 | 2.03 | 0.46 | 0.25 | 15.8 | | | |
| | 1999.587 | 0.107 | 14.184 | 2.02 | 0.46 | 0.25 | 15.8 | 4.5 $`\pm `$ 8.7 | 0.08 $`\pm `$ 0.16 | 171 |
| B3… | 1994.971 | 0.116 | 15.355 | 2.13 | 1.15 | 0.31 | $``$25.5 | | | |
| | 1999.587 | 0.063 | 15.670 | 1.56 | 1.15 | 0.31 | $``$25.5 | 76.0 $`\pm `$ 8.7 | 1.36 $`\pm `$ 0.16 | $``$24 |
| B2… | 1994.971 | 0.115 | 15.854 | 2.66 | 1.05 | 0.29 | 13.1 | | | |
| | 1999.587 | 0.126 | 15.826 | 2.80 | 1.05 | 0.29 | 13.1 | 10.4 $`\pm `$ 8.7 | 0.19 $`\pm `$ 0.16 | 129 |
| A3… | 1994.971 | 0.126 | 23.743 | 0.43 | 2.11 | 0.60 | $``$68.7 | | | |
| | 1997.092 | 0.095 | 23.679 | 0.60 | 2.11 | 0.60 | $``$68.7 | 20.6 $`\pm `$ 11 | 0.37 $`\pm `$ 0.20 | 133 |
| A2… | 1994.971 | 0.185 | 24.333 | 1.16 | 0.51 | 0.65 | 51.4 | | | |
| | 1999.587 | 0.189 | 24.468 | 1.17 | 0.51 | 0.65 | 51.4 | 29.2 $`\pm `$ 8.7 | 0.52 $`\pm `$ 0.16 | 3 |
Notes to Table 5
See notes to Table 4.
TABLE 6
Gaussian Model and Relative Proper Motions for 1031+567
| Component | Epoch | $`S`$ | $`r`$ | $`\theta `$ | $`a`$ | $`b/a`$ | $`\mathrm{\Phi }`$ | $`\mu `$ | $`v`$ | p.a. |
| --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- |
| | | (Jy) | (mas) | () | (mas) | | () | ($`\mu `$as/yr) | ($`h^1`$ c) | () |
| W1… | 1995.411 | 0.078 | 0.0 | 0.0 | 0.69 | 0.71 | 13.2 | | | |
| | 1999.587 | 0.080 | 0.0 | 0.0 | 0.69 | 0.71 | 13.2 | reference | | |
| W2… | 1995.411 | 0.097 | 1.075 | 62.24 | 1.68 | 0.54 | 71.5 | | | |
| | 1999.587 | 0.099 | 1.113 | 59.31 | 1.68 | 0.54 | 71.5 | 16.3 $`\pm `$ 12 | 0.27 $`\pm `$ 0.20 | 5 |
| E2… | 1995.411 | 0.085 | 32.057 | 49.64 | 2.42 | 0.56 | 44.0 | | | |
| | 1999.587 | 0.092 | 32.211 | 49.74 | 2.42 | 0.56 | 44.0 | 37.6 $`\pm `$ 8.4 | 0.61 $`\pm `$ 0.14 | 70 |
| E1… | 1995.411 | 0.065 | 33.699 | 50.34 | 0.72 | 1.00 | 0.00 | | | |
| | 1999.587 | 0.064 | 33.748 | 50.26 | 0.72 | 1.00 | 0.00 | 14.6 $`\pm `$ 4.8 | 0.24 $`\pm `$ 0.08 | 6 |
Notes to Table 6
See notes to Table 4.
TABLE 7
Gaussian Model and Relative Proper Motions for 2352+495
| Component | Epoch | $`S`$ | $`r`$ | $`\theta `$ | $`a`$ | $`b/a`$ | $`\mathrm{\Phi }`$ | $`\mu `$ | $`v`$ | p.a. |
| --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- |
| | | (Jy) | (mas) | () | (mas) | | () | ($`\mu `$as/yr) | ($`h^1`$ c) | () |
| C… | 1994.971 | 0.024 | 0.0 | 0.0 | 2.06 | 0.61 | $``$81.0 | | | |
| | 1999.587 | 0.019 | 0.0 | 0.0 | 2.06 | 0.61 | $``$81.0 | reference | | |
| D… | 1994.971 | 0.013 | 25.237 | $``$20.29 | 0.64 | 0.17 | 4.1 | | | |
| | 1999.587 | 0.012 | 25.310 | $``$20.18 | 0.64 | 0.17 | 4.1 | 19.1 $`\pm `$ 32 | 0.18 $`\pm `$ 0.31 | 13 |
| B5… | 1994.971 | 0.084 | 30.104 | $``$15.72 | 0.53 | 0.53 | 8.4 | | | |
| | 1999.587 | 0.048 | 30.433 | $``$15.42 | 0.53 | 0.53 | 8.4 | 79.1 $`\pm `$ 11 | 0.76 $`\pm `$ 0.11 | 10 |
| B4… | 1994.971 | 0.156 | 30.913 | $``$16.30 | 0.60 | 0.80 | $``$54.3 | | | |
| | 1999.587 | 0.117 | 31.025 | $``$16.43 | 0.60 | 0.80 | $``$54.3 | 28.6 $`\pm `$ 11 | 0.27 $`\pm `$ 0.11 | $``$48 |
| B3… | 1994.971 | 0.075 | 31.390 | $``$15.66 | 1.56 | 0.39 | 0.4 | | | |
| | 1999.587 | 0.101 | 31.227 | $``$15.74 | 1.56 | 0.39 | 0.4 | 36.6 $`\pm `$ 11 | 0.35 $`\pm `$ 0.11 | 179 |
| B2… | 1994.971 | 0.120 | 33.047 | $``$15.38 | 0.94 | 0.63 | $``$25.4 | | | |
| | 1997.092 | 0.064 | 33.229 | $``$15.36 | 0.94 | 0.63 | $``$25.4 | 39.4 $`\pm `$ 11 | 0.38 $`\pm `$ 0.11 | $``$12 |
| B1… | 1994.971 | 0.085 | 33.360 | $``$17.03 | 1.20 | 0.70 | $``$5.8 | | | |
| | 1999.587 | 0.058 | 33.592 | $``$17.00 | 1.20 | 0.70 | $``$5.8 | 50.5 $`\pm `$ 11 | 0.48 $`\pm `$ 0.11 | $``$13 |
| A… | 1994.971 | 0.055 | 49.243 | $``$17.47 | 1.20 | 0.69 | $``$68.2 | | | |
| | 1999.587 | 0.042 | 49.392 | $``$17.44 | 1.20 | 0.69 | $``$68.2 | 32.7 $`\pm `$ 11 | 0.31 $`\pm `$ 0.11 | $``$8 |
Notes to Table 7
See notes to Table 4.
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# Collapse-revivals and population trapping in the 𝑚-photon mazer
## I Introduction
Laser cooling of atoms is a rapidly developing field in quantum optics. Cold and ultracold atoms introduce new regimes in atomic physics often not considered in the past. Recently, Scully et al. have shown that a new kind of induced emission occurs when a micromaser is pumped by ultracold atoms, requiring a quantum-mechanical treatment of the center-of-mass motion. They called this particular process mazer action to insist on the quantized $`z`$-motion feature of the induced emission.
The detailed quantum theory of the mazer has been presented in a series of three papers by Scully and co-workers . They showed that the induced emission probability is strongly dependent on the cavity mode profile. Analytical calculations were presented for the mesa and the sech<sup>2</sup> mode profiles. For sinusoidal modes, WKB solutions were detailed.
Retamal et al. showed that we must go beyond the WKB solutions for the sinusoidal mode case when we consider strictly the ultracold regime. Remarkably, they showed that the resonances in the emission probability are not completely smeared out for actual interaction and cavity parameters. In a recent work , we proposed a numerical method for calculating efficiently the induced emission probability for arbitrary cavity field modes. In particular, the gaussian potential was considered, thinking in open cavities in the microwave or optical field regime. Differences with respect to the sech<sup>2</sup> mode case were found. Calculations for sinusoidal potentials were also performed and divergences with WKB results were reported, confirming results given in .
Zhang et al. extended the concept of the mazer to the two-photon process by proposing the idea of the two-photon mazer. Their work was focused on the study of its induced emission probability in the special case of the mesa mode function. Under the condition of an initial coherent field state, they showed that this probability exhibits with respect to the interaction length the collapse-revivals phenomenon, which have different features in different regimes. They are similar to those in the two-photon Jaynes-Cummings model only in the thermal-atom regime.
The collapse-revivals of the atomic excitation in the framework of the Jaynes-Cummings model was predicted in the early 1980s by Eberly and co-workers . Fleischhauer and Schleich showed later that the shape of each revival is a direct reflection of the shape of the initial photon-number distribution $`P_n`$, assuming that the atom is prepared completely in the upper state or in the lower state and that the distribution $`P_n`$ is sufficiently smooth. It was also noticed that, under some special conditions of the initial atom-field state, the revivals can be largely and even completely suppressed . This phenomenon was denominated “population trapping” to refer, as noted by Yoo and Eberly , to a persistent probability of finding the atom in a given level in spite of the existence of both the radiation field and allowed transitions to other levels. The initial atom-field states giving rise to this phenomenon were called “trapping states” in . Let us mention that this denomination is actually used in various physical contexts whenever a degree of freedom is found unaltered in spite of the existence of an interaction able to change its value. For instance trapping states in the context of the micromaser theory have been predicted and very recently measured by Filipowicz et al. and Weidinger et al. respectively. Nevertheless, these trapping states do not relate with the suppression of the collapse-revivals we are dealing here.
An elegant explanation of the population trapping phenomenon has been proposed just very recently by Jonathan et al., who noticed that the key to understand the collapse-revival patterns under very general conditions is to consider the joint initial properties of the atom-field system, even if this one is completely disentangled before the interaction. By defining an appropriate coordinate system, the dressed-state coordinates, they were able to yield simple analytical expressions for the atomic populations which exhibit the conditions needed for population trapping.
At the present time, no work has been devoted to know whether the revivals predicted by Zhang et al. for the two-photon mazer may also be suppressed by use of an appropriate initial state of the atom-field system. An answer to this question is given at the end of this paper. To be not restricted to the case of the two-photon mazer, our analysis is generalized to the arbitrary $`m`$-photon mazer system, although the construction of real multiphoton cavities results in a formidable experimental task.
In Sec. II, we write the quantum theory of the $`m`$-photon mazer by use of the dressed-state coordinate formalism as it was very efficient in the description of the population trapping phenomenon in the Jaynes-Cummings model . General expressions are derived for the atomic populations and the cavity photon distribution after the interaction of the atom with the cavity. The theory is written for any initial pure state of the atom-field system (entangled or not). We consider zero temperature and no dissipation in the high-$`Q`$ cavity. In Sec. III, results of Zhang et al. are extended to the 1, 2 and 3-photon mazer systems and to various cavity mode profiles (mesa, sech<sup>2</sup>, and gaussian ones). Collapse-revival patterns are described for atoms prepared completely in the upper state or in the lower state. Sec. IV is devoted to the study of the $`m`$-photon mazer trapping states which suppress the collapse-revivals. A brief summary of our results is given in Sec. V.
## II The model
### A The Hamiltonian
We consider a two-level atom moving along the $`z`$-direction in the way to a cavity of length $`L`$. The atom is coupled resonantly with an $`m`$-photon transition to a single mode of the quantized field present in the cavity. The atom-field interaction is modulated by the cavity field mode function. The atomic center-of-mass motion is described quantum mechanically and the rotating-wave approximation is made. In the interaction picture, the Hamiltonian describing the system is
$$H=\frac{p^2}{2M}+\mathrm{}gu(z)(a^m\sigma +a^m\sigma ^{}),$$
(1)
where $`p`$ is the atomic center-of-mass momentum along the $`z`$-axis, $`M`$ is the atomic mass, $`\sigma =|ba|`$ ($`|a`$ and $`|b`$ are respectively the upper and lower levels of the $`m`$-photon transition), $`a`$ and $`a^{}`$ are respectively the annihilation and creation operators of the cavity radiation field, $`g`$ is the atom-field coupling strength (half the Rabi frequency) and $`u(z)`$ is the cavity field mode.
### B The wavefunctions
In the $`z`$-representation and in the dressed-state basis
$$\{\begin{array}{c}|b,0,\mathrm{},|b,m1,\hfill \\ |\pm ,n=\frac{1}{\sqrt{2}}\left(|a,n\pm |b,n+m\right),\hfill \end{array}$$
(2)
$`|n`$ being the photon-number states, the problem reduces to the scattering of the atom upon the potentials $`V_n^\pm (z)=\pm \mathrm{}g\sqrt{(n+1)\mathrm{}(n+m)}u(z)`$. Indeed, the set of wavefunction components
$$\psi _n^\pm (z,t)=z,\pm ,n|\psi (t),$$
(3)
where $`|\psi (t)`$ is the atom-field state satisfy the Schrödinger equation
$$i\mathrm{}\frac{}{t}\psi _n^\pm (z,t)=\left(\frac{\mathrm{}^2}{2M}\frac{^2}{z^2}+V_n^\pm (z)\right)\psi _n^\pm (z,t).$$
(4)
The general solution of (4) is
$$\psi _n^\pm (z,t)=𝑑k\varphi _n^\pm (k)e^{i\frac{\mathrm{}k^2}{2M}t}\phi _n^\pm (k,z),$$
(5)
where $`\phi _n^\pm (k,z)`$ is solution of the time-independent Schrödinger equation
$$\left(\frac{^2}{z^2}+k^2\kappa _n^2u(z)\right)\phi _n^\pm (k,z)=0,$$
(6)
with
$$\kappa _n=\kappa \sqrt[4]{(n+1)\mathrm{}(n+m)}$$
(7)
and
$$\kappa =\sqrt{2Mg/\mathrm{}}.$$
(8)
The wavefunction components ($`n=1,\mathrm{},m`$)
$$\psi _n(z,t)=z,b,mn|\psi (t)$$
(9)
satisfy a Schrödinger equation characterized with a null potential and are therefore not affected by the interaction of the atom with the cavity. The atom in the lower state cannot obviously interact with the cavity field that contains less than $`m`$ photons. The components (9) describe a free particle problem.
We assume that, initially, the atomic center-of-mass motion is not correlated to the other degrees of freedom. We describe it by the wave packet
$$\chi (z)z|\chi =𝑑kA(k)e^{ikz}\theta (z),$$
(10)
where $`\theta (z)`$ is the Heaviside step function (indicating that the atoms are incident from the left of the cavity). No restrictions are made for the initial conditions of the atomic internal state and the cavity field state, except that pure states are only considered. By use of an expansion over the dressed-state basis (2), we may write
$$|\psi (0)=|\chi \left(\underset{n=1}{\overset{m}{}}w_ne^{i\chi _n}|b,mn+\underset{n=0}{\overset{\mathrm{}}{}}w_ne^{i\chi _n}|\beta _n\right),$$
(11)
with
$$|\beta _n=\mathrm{cos}\left(\frac{\theta _n}{2}\right)|+,n+e^{i\varphi _n}\mathrm{sin}\left(\frac{\theta _n}{2}\right)|,n.$$
(12)
The parameters $`w_n[0,1]`$, $`\theta _n[0,\pi ]`$ and $`\chi _n`$, $`\varphi _n[0,2\pi ]`$ are called dressed-state coordinates . The normalisation condition is
$$\underset{n=m}{\overset{\mathrm{}}{}}w_n^2=1$$
(13)
and the phase factor $`\chi _m`$ may be set to $`0`$ without loss of generality.
We consider therefore
$$\{\begin{array}{c}\psi _n(z,0)=c_n\chi (z),\hfill \\ \psi _n^\pm (z,0)=c_n^\pm \chi (z),\hfill \end{array}$$
(14)
with
$$\{\begin{array}{c}c_n=w_ne^{i\chi _n},\hfill \\ c_n^+=w_ne^{i\chi _n}\mathrm{cos}\left(\theta _n/2\right),\hfill \\ c_n^{}=w_ne^{i(\chi _n\varphi _n)}\mathrm{sin}\left(\theta _n/2\right).\hfill \end{array}$$
(15)
Inserting Eqs. (2) and (10) into Eq. (11), we get
$`|\psi (0)`$ $`=`$ $`{\displaystyle }dz{\displaystyle }dkA(k)\times `$ (19)
$`({\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}[S_{a,n}e^{ikz}\theta (z)|z,a,n`$
$`+S_{b,n+m}e^{ikz}\theta (z)|z,b,n+m]`$
$`+{\displaystyle \underset{n=1}{\overset{m}{}}}w_ne^{i\chi _n}e^{ikz}\theta (z)|z,b,mn),`$
with
$$\left(\begin{array}{c}S_{a,n}\\ S_{b,n+m}\end{array}\right)=\stackrel{~}{A}_n\left(\begin{array}{c}1\\ 1\end{array}\right)$$
(20)
and
$$\stackrel{~}{A}_n=\frac{w_ne^{i\chi _n}}{\sqrt{2}}\left(\begin{array}{cc}\mathrm{cos}\left(\theta _n/2\right)& e^{i\varphi _n}\mathrm{sin}\left(\theta _n/2\right),\\ \mathrm{cos}\left(\theta _n/2\right)& e^{i\varphi _n}\mathrm{sin}\left(\theta _n/2\right),\end{array}\right)$$
(21)
After the atom has left the interaction region, the wavefunctions $`\phi _n^\pm (k,z)`$ can be written as
$$\phi _n^\pm (k,z)=\{\begin{array}{cc}r_n^\pm (k)e^{ikz}\hfill & (z<0)\hfill \\ t_n^\pm (k)e^{ik(zL)}\hfill & (z>L)\hfill \end{array},$$
(22)
where $`r_n^\pm (k)`$ and $`t_n^\pm (k)`$ are respectively the reflection and transmission coefficient associated with the scattering of the particle of momentum $`\mathrm{}k`$ upon the potential $`V_n^\pm (z)`$ (Eq. 6). The initial state components $`\psi _n^\pm (z,0)`$ have evolved into
$`\psi _n^\pm (z,t)`$ $`=`$ $`c_n^\pm {\displaystyle }dkA(k)e^{i\frac{\mathrm{}k^2}{2M}t}[r_n^\pm (k)e^{ikz}\theta (z)`$ (24)
$`+t_n^\pm (k)e^{ik(zL)}\theta (zL)]`$
whereas the free particle wavefunction components $`\psi _n(z,0)`$ become
$$\psi _n(z,t)=c_n𝑑kA(k)e^{i\frac{\mathrm{}k^2}{2M}t}e^{ik(zL)}\theta (zL).$$
(25)
We thus obtain
$`|\psi (t)`$ $`=`$ $`{\displaystyle }dz{\displaystyle }dkA(k)e^{i\frac{\mathrm{}k^2}{2M}t}\times `$ (31)
$`({\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}[R_{a,n}(k)e^{ikz}\theta (z)|z,a,n`$
$`+T_{a,n}(k)e^{ik(zL)}\theta (zL)|z,a,n`$
$`+R_{b,n+m}(k)e^{ikz}\theta (z)|z,b,n+m`$
$`+T_{b,n+m}(k)e^{ik(zL)}\theta (zL)|z,b,n+m]`$
$`+{\displaystyle \underset{n=1}{\overset{m}{}}}w_ne^{i\chi _n}e^{ik(zL)}\theta (zL)|z,b,mn),`$
in which
$$\left(\begin{array}{c}R_{a,n}(k)\\ R_{b,n+m}(k)\end{array}\right)=\stackrel{~}{A}_n\left(\begin{array}{c}r_n^+(k)\\ r_n^{}(k)\end{array}\right),$$
(33)
$$\left(\begin{array}{c}T_{a,n}(k)\\ T_{b,n+m}(k)\end{array}\right)=\stackrel{~}{A}_n\left(\begin{array}{c}t_n^+(k)\\ t_n^{}(k)\end{array}\right).$$
(34)
If initially the electromagnetic field is in the state $`|n`$ and the atom is in the excited state $`|a`$, the only non-zero dressed-state coordinates are $`w_n=1`$ and $`\theta _n=\pi /2`$. We get therefore
$$\stackrel{~}{A}_n=\frac{1}{2}\left(\begin{array}{cc}1& 1\\ 1& 1\end{array}\right)$$
(35)
and Eqs. (II B) lead to the same results given by Meyer et al. who considered in detail this case for the one-photon mazer.
### C Atomic populations
The reduced density matrix $`\sigma (t)`$ for the atomic internal degree of freedom is given by the trace over the radiation and the atomic external variables of the atom-field density matrix, that is its elements $`i,j=a,b`$ are
$$\sigma _{ij}(t)=\underset{n}{}𝑑zz,i,n|\psi (t)\psi (t)|z,j,n.$$
(36)
The atomic populations $`\sigma _{ii}`$ follows immediately from Eq. (36) :
$$\sigma _{ii}(t)=\underset{n}{}𝑑z|z,i,n|\psi (t)|^2.$$
(37)
Inserting Eqs. (19) and (31) into Eq. (37) and using Eqs. (20) and (II B), we get for an incident atom of momentum $`\mathrm{}k`$ :
$`\sigma _{aa}(0)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[1{\displaystyle \underset{n=1}{\overset{m}{}}}w_n^2+{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}w_n^2\mathrm{sin}(\theta _n)\mathrm{cos}(\varphi _n)\right],`$ (38)
$`\sigma _{aa}(t)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[1{\displaystyle \underset{n=1}{\overset{m}{}}}w_n^2+{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}w_n^2\mathrm{sin}(\theta _n)\text{Re}(e^{i\varphi _n}K_n)\right],`$ (39)
where
$$K_n=r_n^+r_n^{}+t_n^+t_n^{}.$$
(41)
The change of the atomic population $`\sigma _{aa}`$ induced by the interaction of the incident atom with the cavity radiation field is then given by
$$\delta \sigma _{aa}=\sigma _{aa}(t)\sigma _{aa}(0),$$
(42)
with the time $`t`$ chosen long after the interaction.
Thus we have
$$\delta \sigma _{aa}=\underset{n=0}{\overset{\mathrm{}}{}}\mathrm{\Delta }_n,$$
(43)
with
$$\mathrm{\Delta }_n=\frac{w_n^2}{2}\mathrm{sin}(\theta _n)\left[\text{Re}\left(e^{i\varphi _n}K_n\right)\mathrm{cos}(\varphi _n)\right].$$
(44)
As expected, the components $`w_ne^{i\chi _n}`$ of the initial state $`|\psi (0)`$ over the states $`|b,n`$ $`(n<m)`$ do not play any role in the dynamics of the system.
We have to emphasize that in Eq. (43) $`\mathrm{\Delta }_n`$ *cannot* be interpreted strictly as the change in the $`\sigma _{aa}`$ population induced by the interaction of the two-level atom with the cavity radiation field containing $`n`$ photons. This is only true when the incident atom is prepared in the excited state. Indeed, if initially the internal atomic state is $`c_a|a+c_b|b`$ and the field state is $`|n`$ $`(nm)`$, then the only non-zero dressed-state coordinates are $`w_n=|c_a|`$, $`\chi _n=\mathrm{arg}(c_a)`$, $`\theta _n=\pi /2`$, $`w_{nm}=|c_b|`$, $`\chi _{nm}=\mathrm{arg}(c_b)`$, $`\theta _{nm}=\pi /2`$ and $`\varphi _{nm}=\pi `$. We thus have in that case
$`\delta \sigma _{aa}`$ $`=`$ $`\mathrm{\Delta }_n+\mathrm{\Delta }_{nm},`$ (45)
$`=`$ $`\mathrm{\Delta }_n\text{iff}c_b=0.`$ (46)
### D Photon statistics
The reduced density matrix $`\rho (t)`$ for the cavity radiation field is given by the trace over the internal and external atomic degrees of freedom of the atom-field density matrix, that is its elements $`n,n^{}`$ are
$$\rho _{nn^{}}(t)=\underset{i=a,b}{}𝑑zz,i,n|\psi (t)\psi (t)|z,i,n^{}.$$
(47)
The photon distribution $`P_n=\rho _{nn}`$ follows immediately from Eq. (47) :
$$P_n(t)=\underset{i=a,b}{}𝑑z|z,i,n|\psi (t)|^2.$$
(48)
The change $`\delta P_n`$ in the cavity photon distribution induced by the interaction of the cavity electromagnetic field with the incident atom is then given by
$$\delta P_n=P_n(t)P_n(0).$$
(49)
Inserting Eqs. (19) and (31) into Eq. (48) and using Eqs. (20) and (II B), we get for an incident atom of momentum $`\mathrm{}k`$ :
$$\delta P_n=\{\begin{array}{cc}\mathrm{\Delta }_n\mathrm{\Delta }_{nm}\hfill & (nm),\hfill \\ \mathrm{\Delta }_n\hfill & (n<m).\hfill \end{array}$$
(50)
We see that if the initial state is $`|a,n`$ we have
$$\delta \sigma _{aa}+\delta P_{n+m}=0,$$
(51)
which gives an intuitive population conservation condition.
### E Reflection and transmission probabilities
The reflection and transmission probabilities of the incident atom upon the cavity are respectively given by
$`R`$ $`=`$ $`{\displaystyle \underset{i=a,b}{}}{\displaystyle \underset{n}{}}{\displaystyle _{\mathrm{}}^0}𝑑z|z,i,n|\psi (t)|^2,`$ (53)
$`T`$ $`=`$ $`{\displaystyle \underset{i=a,b}{}}{\displaystyle \underset{n}{}}{\displaystyle _L^{\mathrm{}}}𝑑z|z,i,n|\psi (t)|^2.`$ (54)
Inserting Eq. (31) into Eqs. (II E), we get for an incident atom of momentum $`\mathrm{}k`$ :
$`R`$ $`=`$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}w_n^2\left(\mathrm{cos}^2(\theta _n/2)|r_n^+|^2+\mathrm{sin}^2(\theta _n/2)|r_n^{}|^2\right),`$ (56)
$`T`$ $`=`$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}w_n^2\left(\mathrm{cos}^2(\theta _n/2)|t_n^+|^2+\mathrm{sin}^2(\theta _n/2)|t_n^{}|^2\right)`$ (58)
$`+{\displaystyle \underset{n=1}{\overset{m}{}}}w_n^2.`$
One verifies immediately that the results of Meyer et al. about the reflection and transmission probabilities are well recovered by Eqs. (II E) when their initial conditions are considered. Indeed, when the atom-field system is initially in the state $`|a,n`$, Eqs. (II E) become
$`R`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(|r_n^+|^2+|r_n^{}|^2\right),`$ (60)
$`T`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(|t_n^+|^2+|t_n^{}|^2\right).`$ (61)
We get the same results if the atom-field system is initially in the state $`|b,n`$ with $`nm`$, except that $`n`$ must be replaced by $`nm`$ in Eqs. (II E). In the case $`n<m`$, we have obviously $`T=1`$.
### F Final remarks
All the results here above (about the atomic populations, the photon statistics, and the reflection and transmission probabilities) may be very easily generalized for any momentum wavefunction $`A(k)`$ of the initial wave packet. The various expressions must simply be weighted by $`|A(k)|^2`$ and integrated over $`k`$. For instance, Eq. (43) becomes
$$\delta \sigma _{aa}=𝑑k|A(k)|^2\underset{n=0}{\overset{\mathrm{}}{}}\mathrm{\Delta }_n,$$
(62)
where $`\mathrm{\Delta }_n`$ depends on $`k`$ through the reflection and transmission coefficients, $`r_n^\pm (k)`$ and $`t_n^\pm (k)`$ respectively, in $`K_n`$ (see Eq. (44)).
The expressions obtained for all these various physical quantities are very simple in the framework of the dressed-state formalism, even though they are very general. They take a form much more complicated when the usual coordinates of the atom-field system are used (the complex coefficients $`c_a`$, $`c_b`$ and $`c(n)`$ of the atom-field states written as $`(c_a|a+c_b|b)_nc(n)|n`$). Also entangled initial states may be considered by this formalism. The great advantage of the dressed-state coordinates was already pointed out by Jonathan et al. who used them to express various physical quantities in the Jaynes-Cummings model.
## III Collapse-revivals
Expressions (43) and (50) show that the features of the changes in the atomic populations and in the photon distribution $`P_n`$ with respect to the interaction length $`\kappa L`$ are directly related to the characteristics of $`\mathrm{\Delta }_n`$. This in turn depends on the atom-field initial state (through the dressed-state coordinates) and on the cavity field mode profile $`u(z)`$ which affect the reflection and transmission coefficients, $`r_n^\pm `$ and $`t_n^\pm `$ respectively, and thus $`K_n`$.
If $`K_n`$ have a strong oscillatory behaviour with respect to $`\kappa L`$, we may expect collapse-revivals in the population changes when several modes of the field are initially filled. In the following we will restrict ourselves to the description of the collapse-revivals when the atom is prepared completely in the upper or in the lower state, and the field is in the state $`_nc(n)|n`$.
When the atom is initially in the upper state $`|a`$, $`\sigma _{bb}(t)=1\sigma _{aa}(t)`$ represents the probability that a photon be emitted by the atom due to its interaction with the cavity. From Eq. (39), one gets that this induced emission probability is given by
$$P_{em}=\underset{n=0}{\overset{\mathrm{}}{}}p(n)P_{em}(n),$$
(63)
with $`p(n)=|c(n)|^2`$ and
$$P_{em}(n)=\frac{1\text{Re}(K_n)}{2}.$$
(64)
When the atom is initially in the lower state $`|b`$, $`\sigma _{aa}(t)`$ represents the probability that a cavity photon be absorbed by the atom. This absorption probability is identical to the induced emission probability $`P_{em}`$, except that $`p(n)`$ must be replaced by $`p(n+m)`$ in Eq. (63).
The induced emission probability $`P_{em}`$ is studied hereafter for different cavity mode profiles: mesa, sech<sup>2</sup> and gaussian modes. Our description is restricted to the ultracold regime (incident atoms with a momentum $`\mathrm{}k`$ such that $`k/\kappa 1`$).
### A Mesa mode
In the special case where the cavity field mode profile is given by the mesa function
$$u(z)=\{\begin{array}{cc}1\hfill & \text{for}\mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}0}<z<L\hfill \\ 0\hfill & \text{elsewhere}\hfill \end{array}$$
(65)
the reflection and transmission coefficients $`r_n^\pm (k)`$ and $`t_n^\pm (k)`$ respectively may be calculated analytically. Their expression has been given for the one-photon mazer by Löffler et al.. Same results are obtained for the $`m`$-photon mazer, except that the value of the parameter $`\kappa _n`$ must be changed accordingly (see Eq. (7)). Inserting these results into Eq. (64), one gets when $`\mathrm{exp}(\kappa _nL)1`$ and $`(\kappa _n/2k)^2\mathrm{exp}(\kappa _nL)\mathrm{sin}(\kappa _nL)1`$ that
$$P_{em}(n)=\frac{\frac{1}{2}\left[1+\frac{1}{2}\mathrm{sin}(2\kappa _nL)\right]}{1+(\kappa _n/2k)^2\mathrm{sin}^2(\kappa _nL)}.$$
(66)
As pointed out by Löffler et al., Eq. (66) is similar to the Airy function of the classical optics which gives the transmitted intensity in a Fabry-Perot interferometer. This by no means exhibits a strong oscillatory behaviour. Thus, we have no chance to obtain similar collapse-revivals as those in the Jaynes-Cummings model when several mode of the field are initially filled. This is illustrated in Fig. 1 for the 1, 2 and 3-photon mazers where the cavity field is taken initially in a coherent state ($`p(n)=e^{\overline{n}}\frac{\overline{n}^n}{n!}`$) with a mean photon number $`\overline{n}=10`$. A chaotic behaviour in the curves $`P_{em}(\kappa L)`$ is clearly obtained for $`k/\kappa =0.1`$.
### B Sech<sup>2</sup> mode
When the cavity field mode profile is given by the sech<sup>2</sup> function
$$u(z)=\text{sech}^2(z/L)$$
(67)
the reflection and transmission coefficients may also be calculated analytically. Their expression has been given for the one-photon mazer by Löffler et al.. Again, same results are obtained for the $`m`$-photon mazer, except that the value of the parameter $`\kappa _n`$ must be changed accordingly. Hence, the curves $`P_{em}(n)`$ with respect to the interaction length $`\kappa _nL`$ are identical for the one-photon and the $`m`$-photon mazers. Two such curves have been presented by Löffler et al. for $`k/\kappa _n=0.1`$ and $`k/\kappa _n=0.01`$. These curves present well resolved resonances that get smeared for large values of the interaction length. We have calculated for the 1, 2 and 3-photon mazers the induced emission probability in the case of a cavity field initially in a coherent state (with $`\overline{n}=10`$). These results are presented on Fig. 2 for $`k/\kappa =0.1`$. Evidence for collapse-revivals is shown on these figures. They are stronger in the case of the 3-photon mazer.
### C Gaussian mode
For a cavity field mode profile described by the gaussian function
$$u(z)=e^{\frac{z^2}{2\sigma ^2}}$$
(68)
the reflection and transmission coefficients can no more be calculated analytically.
We proposed recently a numerical method for computing efficiently these coefficients and the induced emission probability $`P_{em}(n)`$. We compare on Fig. 3 the results obtained for this probability to those calculated in the case of the sech<sup>2</sup> mode profile. We have considered $`k/\kappa _n=0.1`$ and interaction lengths $`\kappa _nL`$ varying between 0 and 20. The parameter $`\sigma `$ in Eq. (68) was fixed to $`\sqrt{2/\pi }L`$ in order to adopt the same normalization factor for the two profiles (identical area under the modes). As we pointed out in , the resonances in the curves get smeared with increasing values of $`\kappa _nL`$ for both profiles. But this fact is not so marked in the case of the gaussian profile where the resonances still exist for longer interaction lengths. It is a result to be expected as the gaussian profile is growing more abruptly than the $`\text{sech}^2`$ one. Thus it is in some sense “closer” to the mesa mode, which exhibits resonances at infinity.
We have then calculated the induced emission probability $`P_{em}`$ with respect to $`\kappa L`$ for a field initially in a coherent state (with $`\overline{n}=10`$). The result is presented on Fig. 4 for the 1, 2 and 3-photon mazers. As the curves for the probability $`P_{em}(n)`$ are qualitatively similar in the cases of the gaussian and the sech<sup>2</sup> modes, it is not a surprising result that the collapse-revivals are also similar in both cases. Nevertheless, they are stronger for the gaussian potential because the oscillations in $`P_{em}(n)`$ are stronger too.
Squeezing the field inside the cavity has an effect on the collapse-revival patterns. We have considered initial photon distributions $`p(n)`$ inside the cavity corresponding to various squeezed coherent states $`|\alpha ,re^{i\theta }`$, namely
$`p(n)`$ $`=`$ $`{\displaystyle \frac{(\mathrm{tanh}r)^n}{2^nn!\mathrm{cosh}r}}\left|H_n\left({\displaystyle \frac{\alpha e^{i\theta /2}}{\sqrt{2\mathrm{cosh}r\mathrm{sinh}r}}}\right)\right|^2`$ (70)
$`\times \mathrm{exp}\left[|\alpha |^2+{\displaystyle \frac{1}{2}}(e^{i\theta }\alpha ^2+e^{i\theta }(\alpha ^{})^2)\mathrm{tanh}r\right],`$
where $`H_n(z)`$ designates the $`n^{th}`$ order Hermite polynomial. We noticed that squeeze parameters $`r`$ of the order of 0.3 enhance significantly the collapse-revivals presented on Figs. 2 and 4 (keeping the same coherent parameter $`|\alpha |^2=10`$ and taking $`\theta =0`$), while higher squeeze parameters tend to destroy them.
## IV Population Trapping
When the atom-field initial state is such that $`\mathrm{sin}(\theta _n)=0`$, we get from Eq. (44) $`\mathrm{\Delta }_n=0`$, whatever the value of $`K_n`$. In this case, we have
$$\delta \sigma _{aa}=\delta \sigma _{bb}=\delta P_n=0,$$
(71)
indicating that the interaction of the atom with the cavity radiation field has no effect on the atomic populations $`\sigma _{ii}`$ ($`i=a,b`$) and on the cavity photon distribution $`P_n`$, whatever the cavity field mode function, whatever the cavity interaction length $`\kappa L`$ and whatever the atomic initial velocity. We conclude that the mazer give rise to the perfect population trapping phenomenon, when considering zero temperature and no dissipation in the high-$`Q`$ cavity. This property holds for the ultracold, intermediate and thermal-atom regimes, as it is completely independent on the external atomic degree of freedom. For the same reason, it holds for any momentum wavefunction $`A(k)`$ of the initial wave packet.
The class of states verifying $`\mathrm{sin}(\theta _n)=0`$, named *perfect trapping states*, are given by
$$|\gamma ^\pm =\frac{\gamma ^m|a\pm |b}{\sqrt{1+|\gamma |^{2m}}}\sqrt{1|\gamma |^2}\underset{n=0}{\overset{\mathrm{}}{}}\gamma ^n|n,$$
(72)
where $`\gamma `$ is a complex number with $`|\gamma |<1`$.
Indeed, rewriting these states in terms of the dressed-state basis, we find
$$|\gamma ^\pm =\sqrt{\frac{1|\gamma |^2}{1+|\gamma |^{2m}}}\left(\underset{n=0}{\overset{\mathrm{}}{}}\sqrt{2}\gamma ^{n+m}|\pm ,n\pm \underset{n=0}{\overset{m1}{}}\gamma ^n|b,n\right).$$
(73)
For each $`n`$ there is only a single dressed-state present in the sum of expression (73). Depending on whether it is $`|+,n`$ or $`|,n`$, we have respectively $`\mathrm{sin}(\theta _n/2)=0`$ or $`\mathrm{cos}(\theta _n/2)=0`$, and so $`\mathrm{sin}(\theta _n)=0`$ in any case.
This give rise to another very interesting feature of the perfect trapping states. The reflection and transmission probabilities (II E) become
$`R`$ $`=`$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}w_n^2|r_n^\pm |^2,`$ (75)
$`T`$ $`=`$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}w_n^2|t_n^\pm |^2+{\displaystyle \underset{n=1}{\overset{m}{}}}w_n^2,`$ (76)
with
$`w_n`$ $`=`$ $`\sqrt{{\displaystyle \frac{1|\gamma |^2}{1+|\gamma |^{2m}}}}\sqrt{2}|\gamma |^{n+m},`$ (78)
$`w_n`$ $`=`$ $`\sqrt{{\displaystyle \frac{1|\gamma |^2}{1+|\gamma |^{2m}}}}|\gamma |^{mn}.`$ (79)
The particle moving along the $`z`$-axis is only sensitive to either a superposition of the potentials $`V_n^+(z)`$ or a superposition of $`V_n^{}(z)`$, but never to both. In principle, it would be possible to imagine an experimental set-up where the particles would encounter only an effective potential well, instead of an effective potential hill.
It is important to emphasize that the perfect trapping states do not make the cavity transparent to the incident atoms, because the reflection coefficient $`R`$ is not nullified.
## V Summary
In this paper, we have studied collapse-revival patterns that appear in the changes of the atomic populations induced by the interaction of ultracold two-level atoms with electromagnetic cavities of various interaction lengths that are in resonance with an $`m`$-photon transition of the atoms. In particular, the sech<sup>2</sup> and gaussian cavity mode profiles have been considered and differences in the collapse-revival patterns are reported. They are stronger in the case of the gaussian potential. With the aim of such studies in view, we have written the quantum theory of the $`m`$-photon mazer by use of the dressed-state coordinate formalism. Simple expressions for the atomic populations, the cavity photon statistics, and the reflection and transmission probabilities have been given for any initial pure state of the atom-field system. The evidence for the population trapping phenomenon has then been very easily given. The trapping states written in Sec. IV have the property to leave, after the atom-field interaction, the cavity field and the internal atomic degrees of freedom at their initial value, independently of the cavity field mode, the cavity interaction length, and the initial atomic velocity.
###### Acknowledgements.
This work has been supported by the Belgian Institut Interuniversitaire des Sciences Nucléaires (IISN) and the Brazilian Conselho Nacional de Desenvolvimento Científico (CNPq). E. S. wants to thank for hospitality to Prof. Werner Vogel and co-workers at Rostock University in Germany.
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# Untitled Document
UT-880 hep-th/0005080
Open-Closed String Field Theory
in the Background $`B`$-Field
Teruhiko Kawano
Department of Physics, University of Tokyo
Hongo, Tokyo 113-0033, Japan
kawano@hep-th.phys.s.u-tokyo.ac.jp
and
Tomohiko Takahashi
Department of Physics, Nara Women’s University
Nara 630-8506, Japan
tomo@phys.nara-wu.ac.jp
In this paper, we study open-closed string field theory in the background B-field in the so-called $`\alpha =p^+`$ formulation. The string field theory in the infrared gives noncommutative gauge theory in the open string sector. Since this theory includes closed string fields as dynamical variables, we can obtain another string field theory in the same background through the condensation of closed string fields, whose low-energy effective action does not show the noncommutativity of spacetime. Although we have two string field theories in the same background, we show that these theories are equivalent. In fact, we give the unitary transformation from string fields in one of them to string fields in the other.
May, 2000
1. Introduction
In our previous paper , we studied Witten’s open string field theory in the background $`B`$-field \[1,,3\]. In the resulting string field theory, the kinetic term has an ordinary form, except that the metric in the BRS charge of the kinetic term is the open string metric $`G_{ij}`$, not the closed string metric $`g_{ij}`$. The open string metric $`G_{ij}`$ is related to the closed string metric $`g_{ij}`$ by $`G_{ij}=g_{ij}b_{ik}g^{kl}b_{lj}`$, where $`b_{ij}=(2\pi \alpha ^{})B_{ij}`$; $`B_{ij}`$ is the background $`B`$-field. The interaction term is also affected by the background $`B`$-field. The interaction term of Witten’s open string field theory is only the three-string vertex $`{}_{321}{}^{}V_3\left|\right|\mathrm{\Psi }_{1}^{}|\mathrm{\Psi }_2|\mathrm{\Psi }_3`$. In the presence of the $`B`$-field, the three-string vertex $`V_3|`$ becomes $`V_3|\mathrm{exp}[(i/4)_{r<s}^3\theta ^{ij}p_i^{(r)}p_j^{(s)}]`$. These results are in agreement with the results in . Therefore, we can expect that the low-energy physics is effectively described by noncommutative gauge theory.
The novelty of the papers \[1,,3\] is that the noncommutative factor $`e^{[(i/4){\scriptscriptstyle \theta ^{ij}p_i^{\left(r\right)}p_j^{\left(s\right)}}]}`$ on the three-string vertex $`V_3|`$ can be expressed as $`_{r=1}^3U^{(r)}`$, where $`U^{(r)}=e^{M^{(r)}}`$ is a unitary operator of the $`r`$-th string. Therefore, it is very tempting to transform open string fields $`\mathrm{\Psi }`$ by $`\mathrm{\Psi }U^1\mathrm{\Psi }`$ to remove the noncommutative factor $`\mathrm{exp}[(i/4)_{r<s}^3\theta ^{ij}p_i^{(r)}p_j^{(s)}]`$ from the three-string vertex, and to expect that the low-energy dynamics is described by an ordinary gauge theory with background field strength $`B_{ij}`$. In fact, the unitary operator $`U^{(r)}`$ cannot be uniquely determined. In the papers \[1,,3\], a unitary operator $`e^{M^{(r)}}`$ was proposed to be $`M^{(r)}=(i/2)\theta ^{ij}p_{L}^{}{}_{i}{}^{(r)}p_{R}^{}{}_{j}{}^{(r)}`$, where $`p_{L}^{}{}_{i}{}^{}`$ is given by the integration of the momentum $`P_i(\sigma )`$ over half a string $`(0,\pi /2)`$ and $`p_{R}^{}{}_{i}{}^{}`$ over $`(\pi /2,\pi )`$. Upon the unitary transformation, we obtain the transformed BRS charge $`e^MQ_\mathrm{B}e^M`$ in the kinetic term and find that it is divergent. It seems that the divergence comes from the midpoint of open strings: $`\sigma =\pi /2`$. Since the transformed BRS charge should not depend on the types of the string interactions, it is natural to seek candidates for the operator $`M`$ which can be used in other string field theories, as well as in Witten’s string field theory. As we argued in our paper , another operator $`M`$ has been proposed to be $`M=\frac{i}{4}_0^\pi 𝑑\sigma _0^\pi 𝑑\sigma ^{}ϵ(\sigma \sigma ^{})\theta ^{ij}P_i(\sigma )P_j(\sigma ^{})`$, which we called $`\stackrel{~}{M}`$ in . By making use of this operator $`M`$, the transformed BRS charge $`e^MQ_\mathrm{B}e^M`$ can be found to be finite, as we will see below.
Thus, if we perform the unitary transformation $`\mathrm{\Psi }e^M\mathrm{\Psi }`$, we may expect to obtain another string field theory in the same background $`B`$-field, which gives an ordinary gauge theory with the background field strength $`B_{ij}`$ in the infrared. The aim of this paper is to show that this is indeed the case. To this end, it is appropriate to use an open-closed string field theory, because we have closed string fields as dynamical variables in such a theory. Since the antisymmetric tensor $`B_{ij}`$ is a component field of closed string fields, it is easy to understand the background $`B`$-field as the condensation of closed string fields. Thus, we can obtain an open-closed string field theory in the background $`B`$-field from the theory in the absence of the $`B`$-field by shifting closed string fields by the vacuum expectation value of the dynamical $`B`$-field contained in closed string fields. In this paper, we will use an open-closed string field theory recently given by Asakawa, Kugo, and one of the authors \[5--7\].
The operator $`M=\frac{i}{4}_0^\pi 𝑑\sigma _0^\pi 𝑑\sigma ^{}ϵ(\sigma \sigma ^{})\theta ^{ij}P_i(\sigma )P_j(\sigma ^{})`$ can be applied to the open-closed string field theory of the papers \[5--7\]. Indeed, we can express all the interactions of the string field theory in the $`B`$-field by using this operator $`M`$ and the vertices of the string field theory without the $`B`$-field, in the same way as we have done in Witten’s string field theory \[1,,3\]. In this case, we also need the interactions between open string fields and closed string fields. Therefore, one of the new results in our construction is the interaction between open strings and closed strings in the presence of the background $`B`$-field. In this paper, one of our main results is to show that this theory can be transformed into the theory obtained by the condensation of closed string fields by the above transformation $`\mathrm{\Psi }e^M\mathrm{\Psi }`$.
This paper is organized as follows: In section $`2`$, we construct the open-closed string field theory in the $`B`$-field by solving the overlapping conditions to obtain the string vertices. In section $`3`$, we obtain another string field theory in the same background as the theory in section $`2`$ by utilizing the condensation of closed string fields. In section $`4`$, we discuss the transformation of the string field theory in section $`2`$ into the theory in section $`3`$ by using the unitary transformation. Section $`5`$ is devoted to discussion. In the appendix, we solve the overlapping conditions for the open-closed string transition vertex and the open-open-closed string vertex in the background $`B`$-field to obtain these vertices.
2. Open-Closed String Field Theory in the Background $`B`$-Field
Some covariant string field theories have been constructed as covariantized light-cone gauge string field theories in \[8--10\]. The theories include unphysical parameters called string length parameters $`\alpha `$, which are necessary to describe the joining-splitting vertices instead of the light-cone momentum $`p^+`$ in the light-cone gauge string field theory \[11,,12\]. In order to avoid such unphysical parameters, though we have to give up a manifestly Lorentz covariant formulation, we can use the so-called $`\alpha =p^+`$ HIKKO theories , where we eliminate $`p^+`$ and identify $`\alpha `$ as $`p^+`$. This formulation still retains gauge invariance and is distinct from the light-cone gauge string field theory in this respect. In the series of papers \[5--7\], an open-closed string field theory has been constructed in the $`\alpha =p^+`$ HIKKO formulation.
In this paper, we will study this theory in a constant background $`B`$-field. In general, the gauge invariance of this system is subtle due to the divergence from the dilaton tadpoles. Although, as shown in \[6,,7\], gauge invariance can be maintained by the cancellation of the dilaton tadpoles from the disk amplitude and the crosscap amplitude, we will only consider the oriented sector in the theory. Since D-branes and orientifold plans are the sources of the disk and the crosscap, respectively, we will restrict our study to D-branes, despite the importance of orientifold planes in maintaining the full gauge invariance of the theory <sup>1</sup> We could consider another approach to a gauge invariant open-closed string field theory , which is constructed in a curved background with gauge invariance ensured by the Fishler-Susskind mechanism . However, such a theory has not yet been fully constructed..
As we will show in this section, the vertices of our string field theory in the $`B`$-field can be obtained by replacing open string fields $`|\mathrm{\Psi }`$ with $`e^M\widehat{|\mathrm{\Psi }}`$ in the string field theory with no $`B`$-field, where $`M`$ is a particular operator. The aim of this section is to explain the derivation of this prescription.
Before we construct the string field theory in our background, let us recall some facts about the first-quantized open string theory in this background. The worldsheet action of open strings is
$$S=\frac{1}{4\pi \alpha ^{}}𝑑\tau _0^\pi 𝑑\sigma \left(g_{ij}\eta ^{ab}_a\widehat{X}^i_b\widehat{X}^j2\pi \alpha ^{}B_{ij}ϵ^{ab}_a\widehat{X}^i_b\widehat{X}^j\right),$$
where $`g_{ij}`$ is a constant background metric and $`B_{ij}`$ is the constant background antisymmetric field. If the Dirichlet boundary condition is not chosen for all the directions of the string coordinates, the boundary condition can be seen to be $`g_{ij}_\sigma \widehat{X}^j+(2\pi \alpha ^{})B_{ij}_\tau \widehat{X}^j=0`$ at $`\sigma =0,\pi `$. For simplicity, we will impose this boundary condition on all the string coordinates $`\widehat{X}^i(\tau ,\sigma )`$. The conjugate momenta $`P_i(\sigma )`$ are given by $`(1/2\pi \alpha ^{})g_{ij}_\tau \widehat{X}^j(\sigma )+B_{ij}_\sigma \widehat{X}^j(\sigma )`$.
Applying the Dirac quantization method to this system \[16,,17\] (see also the appendix in ), we obtain the commutation relations
$$\begin{array}{cc}\hfill [\widehat{X}^i(\sigma ),P_j(\sigma ^{})]& =i\delta _j^i\delta _c(\sigma ,\sigma ^{}),[P_i(\sigma ),P_j(\sigma ^{})]=0,\hfill \\ & [\widehat{X}^i(\sigma ),\widehat{X}^j(\sigma ^{})]=i\theta ^{ij}\left(\delta _{\sigma ,0}\delta _{\sigma ^{},0}\delta _{\sigma ,\pi }\delta _{\sigma ^{},\pi }\right),\hfill \end{array}$$
where the delta function $`\delta _c(\sigma ,\sigma ^{})`$ is $`(1/\pi )_{n=\mathrm{}}^{\mathrm{}}\mathrm{cos}(n\sigma )\mathrm{cos}(n\sigma ^{})`$ and $`\delta _{\sigma ,\sigma ^{}}`$ is the Kronecker delta. The noncommutative parameter $`\theta ^{ij}`$ is given by $`[(g+2\pi \alpha ^{}B)^1B(g2\pi \alpha ^{}B)^1]^{ij}`$. When turning off the $`B`$-field, we have the usual oscillator expansion of the string coordinates and the momenta
$$X^i(\sigma )=x^i+l_s\underset{n0}{}\frac{i}{n}\alpha _n^i\mathrm{cos}(n\sigma ),P_i(\sigma )=\frac{1}{\pi l_s}G_{ij}\alpha _n^j\mathrm{cos}(n\sigma ),$$
with $`l_s=\sqrt{2\alpha ^{}}`$. In the presence of the $`B`$-field, the string coordinates become $`\widehat{X}^i(\sigma )=X^i(\sigma )+\theta ^{ij}Q_j(\sigma )`$, where the open string metric $`G_{ij}`$ is $`[g(2\pi \alpha ^{})^2Bg^1B]_{ij}`$, while the momenta remain unchanged. The change $`Q^i(\sigma )`$ is given by
$$Q_i(\sigma )=\frac{1}{\pi }p_i\left(\sigma \frac{\pi }{2}\right)+\frac{1}{\pi l_s}\underset{n0}{}\frac{1}{n}G_{ij}\alpha _n^j\mathrm{sin}(n\sigma )$$
and is related to $`P_i(\sigma )`$ by
$$Q_i(\sigma )=\frac{1}{2}_0^\pi 𝑑\sigma ^{}ϵ(\sigma \sigma ^{})P_i(\sigma ^{}).$$
From (2.1), we can obtain the commutation relations of these oscillators
$$[x^i,p_j]=i\delta _j^i,[\alpha _m^i,\alpha _n^j]=m\delta _{m+n}G^{ij}.$$
As we mentioned at the beginning of this section, in the resulting string field theory with the $`B`$-field, we will find that open string fields $`|\mathrm{\Psi }`$ of the theory without the $`B`$-field are multiplied by the factor $`e^M`$. Therefore, it is appropriate to define this operator $`M`$ before we give the vertices of the string field theory. The operator $`M`$ is given by
$$M=\frac{i}{4}_0^\pi 𝑑\sigma _0^\pi 𝑑\sigma ^{}ϵ(\sigma \sigma ^{})\theta ^{ij}P_i(\sigma )P_j(\sigma ^{}),$$
where $`ϵ(\sigma )`$ is a step function which is $`1`$ for $`\sigma >0`$ and $`1`$ for $`\sigma <0`$. When solving the overlapping conditions to obtain the vertices, we will find that the usual vertices in \[5--7\] fail to satisfy those conditions in our background. But if we put the operator $`e^M`$ on each leg of the open strings, the resulting vertices satisfy them. This is similar to what we have done for the midpoint interaction of Witten’s theory in \[1,,3\] <sup>2</sup> In Witten’s theory, we can use another operator to make the vertices satisfy the overlapping conditions. See \[1,,3\] for further details. .
For closed strings, since we are considering a flat Minkowski spacetime, the $`B`$-field cannot alter the dynamics of this system. Indeed, though we have the same worldsheet action as (2.1), except that $`\sigma `$ runs from $`0`$ to $`2\pi `$, not to $`\pi `$, the term including the $`B`$-field can be made into a surface term and can thus be discarded. But later it will be convenient to have kept this term in the action. In the language of the Hamiltonian formulation, keeping the surface term means that we have alternative canonical variables $`X^i(\sigma )`$ and $`P_i(\sigma )`$; from the action, the momenta $`P_i(\sigma )`$ can be seen to be $`(1/2\pi \alpha ^{})g_{ij}\dot{X}^i(\sigma )+B_{ij}X^{}_{}{}^{}j(\sigma )`$.
Since the string coordinates obey the usual equation of motion and the periodic boundary condition is imposed on them, they have the usual oscillator expansion
$$X^i(\sigma )=x^i+\frac{l_s}{2}\underset{n0}{}\frac{i}{n}\left(\alpha _n^ie^{in\sigma }+\overline{\alpha }_n^ie^{in\sigma }\right).$$
It is obvious that the combination $`(P_iB_{ij}X^{}_{}{}^{}j)`$ has the ordinary expansion of the usual momenta, therefore we obtain
$$P_i(\sigma )=\frac{1}{2\pi }p_i+\frac{1}{2\pi l_s}\underset{n0}{}\left[\left(g_{ij}+b_{ij}\right)\alpha _n^je^{in\sigma }+\left(g_{ij}b_{ij}\right)\overline{\alpha }_n^je^{in\sigma }\right],$$
where $`b_{ij}=(2\pi \alpha ^{})B_{ij}`$. By using the usual commutation relations of the string coordinates and the momenta, we can verify that the oscillators satisfy
$$[x^i,p_j]=i\delta _j^i,[\alpha _m^i,\alpha _n^j]=m\delta _{m+n}g^{ij},[\overline{\alpha }_m^i,\overline{\alpha }_n^j]=m\delta _{m+n}g^{ij}.$$
We will also find it useful to define for closed strings the counterpart of the operator $`M`$
$$M=\frac{i}{4}_0^{2\pi }𝑑\sigma _0^{2\pi }𝑑\sigma ^{}ϵ(\sigma \sigma ^{})\theta ^{ij}P_i(\sigma )P_j(\sigma ^{}).$$
The BRS charges in this background are needed to obtain the kinetic terms of our string field theory. In the string field theory, we have the kinetic term of open string fields $`\widehat{|\mathrm{\Psi }}`$ and that of closed string fields $`|\mathrm{\Phi }`$. Because we have the term including the $`B`$-field in the action for closed strings, the BRS charge for closed strings in our background is the same form as that for open strings, when we write them in terms of the string coordinates and the momenta. The ghost part of the BRS charges is not different from the usual one. The matter part can be given by using the energy-momentum tensor
$$\begin{array}{cc}\hfill \widehat{T}_{\pm \pm }& =\frac{1}{4\alpha ^{}}g^{ij}\left[(2\pi \alpha ^{})P_i\pm \left(g_{ik}b_{ik}\right)\widehat{X}^k\right]\left[(2\pi \alpha ^{})P_j\pm \left(g_{jl}b_{jl}\right)\widehat{X}^l\right]\hfill \\ & =\frac{1}{4\alpha ^{}}G_{ij}\left[\widehat{X}^i\pm (2\pi \alpha ^{})\left(G^{ik}\frac{\theta ^{ik}}{2\pi \alpha ^{}}\right)P_k\right]\left[\widehat{X}^j\pm (2\pi \alpha ^{})\left(G^{jl}\frac{\theta ^{jl}}{2\pi \alpha ^{}}\right)P_l\right],\hfill \end{array}$$
where the prime denotes the differentiation with respect to $`\sigma `$. Here, the above expression (2.1) can be applied to the energy-momentum tensor only for open strings. In order to obtain the energy-momentum tensor for closed strings, we have to replace the string coordinates $`\widehat{X}^i(\sigma )`$ with $`X^i(\sigma )`$ in (2.1) and take away the hat on the energy-momentum tensor. Noting the combination $`(2\pi \alpha ^{})g^{ij}(P_iB_{ij}X^{}_{}{}^{}j)`$ in the former expression of the energy-momentum tensor in (2.1), we can see that $`T_{\pm \pm }`$ for closed strings has the usual oscillator expansion. We can use the latter expression of $`\widehat{T}_{\pm \pm }`$ for open strings to understand that $`\widehat{T}_{\pm \pm }`$ also has the usual oscillator expression except that the closed string metric $`g_{ij}`$ is replaced by the open string metric $`G_{ij}`$, as we mentioned in the introduction.
Now, let us consider our open-closed string field theory in the background $`B`$-field. In the absence of the $`B`$-field, the action $`S`$ of the oriented open-closed string field theory is
$$\begin{array}{cc}\hfill & \frac{1}{2}\mathrm{\Psi }\left|Q_\mathrm{B}\right|\mathrm{\Psi }\frac{1}{2}\mathrm{\Phi }\left|Q_\mathrm{B}(b_0^{}𝒫)\right|\mathrm{\Phi }+\frac{g}{3}V_3^\mathrm{o}(1,2,3)||\mathrm{\Psi }_{321}+\frac{g^2}{4}V_4^\mathrm{o}(1,2,3,4)||\mathrm{\Psi }_{4321}\hfill \\ & +\frac{g^2}{3!}V_3^\mathrm{c}(1^\mathrm{c},2^\mathrm{c},3^\mathrm{c})||\mathrm{\Phi }_{321}+gU(1,2^\mathrm{c})||\mathrm{\Phi }_2|\mathrm{\Psi }_1+\frac{g^2}{2}U_\mathrm{\Omega }(1,2,3^\mathrm{c})||\mathrm{\Phi }_3|\mathrm{\Psi }_{21},\hfill \end{array}$$
as has been given in \[6\] <sup>3</sup> In this action, the coefficients $`x_\mathrm{c}`$, $`x_u`$ and $`x_\mathrm{\Omega }`$ given in \[6,,7\] are included in their respective vertices. Therefore, the definition of the vertices is different from that in \[6,,7\] by these multiplicative factors.. Here $`g`$ is the open string coupling constant. The sign $`\widehat{|\mathrm{\Psi }}`$ will be used for an open string field in our theory with $`B`$-field. The multiple products $`|\mathrm{\Phi }_n\mathrm{}|\mathrm{\Phi }_2|\mathrm{\Phi }_1`$ of string fields are denoted by $`|\mathrm{\Phi }_{n\mathrm{}21}`$. In (2.1), we omitted the integrations over the zero modes of the momenta. In the $`\alpha =p^+`$ HIKKO formulation, we identify the $`\alpha `$ parameter as $`p^+`$ for closed strings and as $`2p^+`$ for open strings.
We have five interaction vertices in this theory. These vertices can be obtained up to an overall normalization by solving the overlap condition. Similarly, in our background the string vertices are given by the overlapping condition. For example, the open three-string vertex $`\widehat{V}_3^\mathrm{o}(1,2,3)|`$ in the case where $`\alpha _1,\alpha _2>0,\alpha _3<0`$ is explicitly given by the overlapping condition
$$\begin{array}{cc}& \widehat{V}_3^\mathrm{o}(3,2,1)|\left\{\widehat{\varphi }^{(1)}\left(\sigma _1\right)\widehat{\varphi }^{(3)}\left(\sigma _3\right)\right\}=0,(0<\sigma <\pi \alpha _1),\hfill \\ & \widehat{V}_3^\mathrm{o}(3,2,1)|\left\{\widehat{\varphi }^{(2)}\left(\sigma _2\right)\widehat{\varphi }^{(3)}\left(\sigma _3\right)\right\}=0,(\pi \alpha _1<\sigma <\pi |\alpha _3|),\hfill \end{array}$$
where $`\widehat{\varphi _r}^{(r)}(\sigma _r)`$ denotes $`\widehat{X}^{i(r)}(\sigma _r)`$ or $`(1/\alpha _r)P_i^{(r)}(\sigma _r)`$ of the $`r`$-th open string. Here we use the following parameters: $`\sigma _1=\sigma /\alpha _1`$; $`\sigma _2=(\sigma \pi \alpha _1)/\alpha _2`$; $`\sigma _3=(\pi |\alpha _3|\sigma )/|\alpha _3|`$. Likewise, the usual open three-string vertex $`V_3^\mathrm{o}(1,2,3)|`$ is specified by
$$\begin{array}{cc}& V_3^\mathrm{o}(3,2,1)|\left\{\varphi ^{(1)}\left(\sigma _1\right)\varphi ^{(3)}\left(\sigma _3\right)\right\}=0,(0<\sigma <\pi \alpha _1),\hfill \\ & V_3^\mathrm{o}(3,2,1)|\left\{\varphi ^{(2)}\left(\sigma _2\right)\varphi ^{(3)}\left(\sigma _3\right)\right\}=0,(\pi \alpha _1<\sigma <\pi |\alpha _3|),\hfill \end{array}$$
where $`\varphi ^{(r)}(\sigma _r)`$ is the same as $`\widehat{\varphi }^{(r)}(\sigma _r)`$, but it denotes $`X^{i(r)}(\sigma _r)`$, not $`\widehat{X}^{i(r)}(\sigma _r)`$. This vertex $`V_3^\mathrm{o}(1,2,3)|`$ can be used to solve the overlapping condition (2.1) to get our vertex $`\widehat{V}_3^\mathrm{o}(1,2,3)|`$. Replacing $`X^{i(r)}(\sigma _r)`$ with $`\widehat{X}^{i(r)}(\sigma _r)`$ in the left-hand side of (2.1), we find that those string coordinates do not connect on the vertex by the difference $`Q^{i(r)}(\sigma _r)`$ between $`\widehat{X}^{i(r)}(\sigma _r)`$ and $`X^{i(r)}(\sigma _r)`$. Indeed, using (2.1) and (2.1), we find
$$\begin{array}{cc}& V_3^\mathrm{o}(1,2,3)|\left\{Q^{i(1)}\left(\sigma _1\right)Q^{i(3)}\left(\sigma _3\right)\right\}=V_3^\mathrm{o}(1,2,3)|\frac{1}{2}G^{ij}p_j^{(2)},(0<\sigma <\pi \alpha _1)\hfill \\ & V_3^\mathrm{o}(1,2,3)|\left\{Q^{i(2)}\left(\sigma _2\right)Q^{i(3)}\left(\sigma _3\right)\right\}=V_3^\mathrm{o}(1,2,3)|\left(\frac{1}{2}\right)G^{ij}p_j^{(1)}.(\pi \alpha _1<\sigma <\pi |\alpha _3|)\hfill \end{array}$$
This equation (2.1) helps us to find our open three-string vertex
$$\widehat{V}_3^\mathrm{o}(1,2,3)|=V_3^\mathrm{o}(1,2,3)|\mathrm{exp}\left(\frac{i}{2}\underset{r<s}{}\theta ^{ij}p_i^{(r)}p_j^{(s)}\right),$$
which is the same expression as that in Witten’s open string field theory, as has been found in \[3,,1\].
The noncommutative factor $`\mathrm{exp}[(i/2)_{r<s}\theta ^{ij}p_i^{(r)}p_j^{(s)}]`$ in (2.1) can be rewritten with the operator $`M`$, as has been discussed in . In fact, changing the variables $`\sigma `$ in (2.1) to the above $`\sigma _3`$, we find that the operator $`M`$ for the third string turns into
$$\begin{array}{cc}\hfill M^{(3)}=& \frac{i}{2}_0^{\pi \alpha _1}𝑑\sigma _{\pi \alpha _1}^{\pi |\alpha _3|}𝑑\sigma ^{}\theta ^{ij}\frac{1}{\alpha _3}P_i^{(3)}(\sigma _3)\frac{1}{\alpha _3}P_j^{(3)}(\sigma _3^{})\hfill \\ & +\frac{i}{4}_0^{\pi \alpha _1}𝑑\sigma _0^{\pi \alpha _1}𝑑\sigma ^{}ϵ(\sigma \sigma ^{})\theta ^{ij}\frac{1}{\alpha _3}P_i^{(3)}(\sigma _3)\frac{1}{\alpha _3}P_j^{(3)}(\sigma _3^{})\hfill \\ & +\frac{i}{4}_{\pi \alpha _1}^{\pi |\alpha _3|}𝑑\sigma _{\pi \alpha _1}^{\pi |\alpha _3|}𝑑\sigma ^{}ϵ(\sigma \sigma ^{})\theta ^{ij}\frac{1}{\alpha _3}P_i^{(3)}(\sigma _3)\frac{1}{\alpha _3}P_j^{(3)}(\sigma _3^{}).\hfill \end{array}$$
Putting this operator $`M^{(3)}`$ on the vertex $`\widehat{V}_3^\mathrm{o}(1,2,3)|`$, due to the overlapping condition (2.1), we can see that the second term in the right-hand side of (2.1) can be converted to $`M^{(1)}`$ and the third term to $`M^{(2)}`$. Furthermore, the first term becomes $`(i/2)\theta ^{ij}p_i^{(1)}p_j^{(2)}`$, which can be made into $`(i/2)_{r<s}\theta ^{ij}p_i^{(r)}p_j^{(s)}`$ by using the momentum conservation. Thus,
$$\widehat{V}_3^\mathrm{o}(1,2,3)|=V_3^\mathrm{o}(1,2,3)|\underset{r=1}{\overset{3}{}}e^{M^{(r)}}.$$
In the same manner, the open four-string vertex $`\widehat{V}_4^\mathrm{o}(1,2,3,4)|`$ can be found to be $`V_4^\mathrm{o}(1,2,3,4)|\mathrm{exp}[(1/2)_{r<s}\theta ^{ij}p_i^{(r)}p_j^{(s)}]`$ and be rewritten into
$$\widehat{V}_4^\mathrm{o}(1,2,3,4)|=V_4^\mathrm{o}(1,2,3,4)|\underset{r=1}{\overset{4}{}}e^{M^{(r)}}.$$
Also, we can easily verify that the closed three-string vertex remains intact:
$$\widehat{V}_3^\mathrm{c}(3^\mathrm{c},2^\mathrm{c},1^\mathrm{c})|=V_3^\mathrm{c}(3^\mathrm{c},2^\mathrm{c},1^\mathrm{c})|.$$
The vertex $`U(1,2^\mathrm{c})|`$ gives the transition between an open string and a closed one, and $`U_\mathrm{\Omega }(1,2,3^\mathrm{c})|`$ gives the interaction between two open strings and a closed one. Since the coordinates of the ends of open strings are noncommutative on D-branes and the $`B`$-field can only change the dynamics of closed strings through interaction with open strings, the modification of these vertices by the $`B`$-field are not given by the noncommutative factor $`\mathrm{exp}[(i/2)_{r<s}\theta ^{ij}p_i^{(r)}p_j^{(s)}]`$. However, as discussed in the appendix, by taking advantage of the operator $`M`$ to solve the overlap conditions, we can obtain those vertices
$$\begin{array}{cc}& \widehat{U}(2,1^\mathrm{c})|=U(2,1^\mathrm{c})|e^{M^{(2)}},\hfill \\ & \widehat{U}_\mathrm{\Omega }(3,2,1^\mathrm{c})|=U_\mathrm{\Omega }(3,2,1^\mathrm{c})|e^{M^{(2)}}e^{M^{(3)}}.\hfill \end{array}$$
Thus, summing up our results about the vertices in our string field theory, we only have to substitute $`e^M\widehat{|\mathrm{\Psi }}`$ for $`|\mathrm{\Psi }`$ in the vertices in the absence of the $`B`$-field to obtain all the vertices in our background. Therefore, the action of our string field theory is given by
$$\begin{array}{cc}\hfill \widehat{S}=& \frac{1}{2}\widehat{\mathrm{\Psi }|}\widehat{Q}_\mathrm{B}\widehat{|\mathrm{\Psi }}\frac{1}{2}\mathrm{\Phi }\left|\widehat{Q}_\mathrm{B}(b_0^{}𝒫)\right|\mathrm{\Phi }\hfill \\ & +\frac{g}{3}\widehat{V}_3^\mathrm{o}(1,2,3)|\widehat{|\mathrm{\Psi }}_{321}+\frac{g^2}{4}\widehat{V}_4^\mathrm{o}(1,2,3,4)\left|\widehat{|\mathrm{\Psi }}_{4321}+\frac{g^2}{3!}\widehat{V}_3^\mathrm{c}(1^\mathrm{c},2^\mathrm{c},3^\mathrm{c})|\right|\mathrm{\Phi }_{321}\hfill \\ & +g\widehat{U}(1,2^\mathrm{c})||\mathrm{\Phi }_2\widehat{|\mathrm{\Psi }}_1+\frac{g^2}{2}\widehat{U}_\mathrm{\Omega }(1,2,3^\mathrm{c})||\mathrm{\Phi }_3\widehat{|\mathrm{\Psi }}_{21}.\hfill \end{array}$$
In the next section, we will obtain an alternative string field theory in the same background through the condensation of the $`B`$-field included in closed string fields $`|\mathrm{\Phi }`$.
3. The Condensation of the $`B`$-Field
In the previous section, we obtained the string field theory in the background $`B`$-field by constructing the vertices as the solutions of the overlapping conditions. The dynamical variables of this field theory are closed string fields as well as open string ones. Since the $`B`$-field is a component field of closed string fields, the background $`B`$-field can be obtained through the condensation of closed string fields. Therefore, beginning with the open-closed string field theory (2.1) in the absence of the background $`B`$-field, and shifting closed string fields by the vacuum expectation value of the $`B`$-field, we obtain a string field theory in the same background as that in the previous section. The resulting theory is different from the theory in the previous section, because only the BRS charge is affected by the condensation. Unfortunately, it is difficult to shift closed string fields by a finite vacuum expectation value. Therefore, we will demonstrate the shifting by an infinitesimal amount in the string field theory. Then, integrating the infinitesimal change of the BRS charge under the shifting, we can find the string field theory in the finite $`B`$-field.
3.1. The Background Independence of the String Vertices
Before calculating the contraction of the string vertices with the shifted closed string fields, we need to understand what remains unchanged under the change of the background $`B`$-field: $`B_{ij}B_{ij}+\delta B_{ij}`$, as Kugo and Zwiebach have done in . Although they dealt with the closed string field theory compactified on a torus, we will follow their method and extend the assumption to our system of open strings as well as closed strings. Indeed, we assume that the string coordinates $`X^i`$ and their momenta $`P_i`$ of open strings and closed strings are background independent. In the previous section, we have dealt with the string coordinates $`\widehat{X}^i(\sigma )`$. These coordinates satisfy the boundary condition: $`_\sigma \widehat{X}^j\theta ^{ij}P_j=0`$ at $`\sigma =0,\pi `$. If we assumed that $`\widehat{X}^i(\sigma )`$ was background independent, it would be inconsistent with this boundary condition. To avoid this inconsistency, since $`_\sigma \widehat{X}^i=_\sigma X^i+\theta ^{ij}P_j`$, we should assume that, rather than $`\widehat{X}^i(\sigma )`$, the string coordinates $`X^i(\sigma )`$ of open strings are background independent, as is the momenta $`P_j(\sigma )`$.
By applying this assumption to the energy-momentum tensor of closed strings, as we can see from (2.1), under the infinitesimal change of the background $`B`$-field, we can obtain the variation of the energy-momentum tensor
$$\delta T_{\pm \pm }(\sigma )=\frac{2\pi \alpha ^{}}{2\alpha ^{}}g^{ik}\delta B_{kj}\left(P_i(\sigma )B_{il}X_{}^{}{}_{}{}^{l}(\sigma )\right)X_{}^{}{}_{}{}^{j}(\sigma ).$$
Since the BRS charge for closed strings is given by
$$Q_\mathrm{B}=_0^{2\pi }\frac{d\sigma }{2\pi }\left[c(\sigma )T_{++}(\sigma )+\overline{c}(\sigma )T_{}(\sigma )\right],$$
its variation turns out to be
$$\delta Q_\mathrm{B}=_0^{2\pi }𝑑\sigma \frac{1}{2}\left(c(\sigma )+\overline{c}(\sigma )\right)g^{ik}\delta B_{kj}\left(P_i(\sigma )B_{il}X_{}^{}{}_{}{}^{l}(\sigma )\right)X_{}^{}{}_{}{}^{j}(\sigma ).$$
If we had dealt with the BRS charge for open strings in the above and had replaced the string coordinates $`\widehat{X}^i(\sigma )`$ with the background-independent coordinates $`X^i(\sigma )`$, we would have obtained the variation of the BRS charge in the same form as (3.1) except for the integration over $`(0,\pi )`$ instead of $`(0,2\pi )`$.
By using (2.1) and (2.1), in terms of the string coordinates and the momenta, we can express the oscillators for closed strings as
$$\begin{array}{cc}& \alpha _n^i(B)=_0^{2\pi }\frac{d\sigma }{2\pi l_s}e^{in\sigma }g^{ij}\left[(2\pi \alpha ^{})P_j(2\pi \alpha ^{})B_{jk}X_{}^{}{}_{}{}^{k}+g_{jk}X_{}^{}{}_{}{}^{k}\right],\hfill \\ & \overline{\alpha }_n^i(B)=_0^{2\pi }\frac{d\sigma }{2\pi l_s}e^{in\sigma }g^{ij}\left[(2\pi \alpha ^{})P_j(2\pi \alpha ^{})B_{jk}X_{}^{}{}_{}{}^{k}g_{jk}X_{}^{}{}_{}{}^{k}\right].\hfill \end{array}$$
Here we labeled the oscillators with the letter $`B`$ as $`(\alpha _n^i(B),\overline{\alpha }_n^i(B))`$ to explicitly show their dependence on the background $`B`$-field. These oscillators $`(\alpha _n^i(B),\overline{\alpha }_n^i(B))`$ can be related to the oscillators $`(\alpha _n^i(0),\overline{\alpha }_n^i(0))`$ in the absence of the $`B`$-field by using the operator
$$=\frac{i}{2}_0^{2\pi }𝑑\sigma B_{ij}X_{}^{}{}_{}{}^{i}(\sigma )X^j(\sigma ).$$
In fact, since $`e^{}P_i(\sigma )e^{}=P_i(\sigma )B_{ik}X_{}^{}{}_{}{}^{j}(\sigma )`$, we can see from (3.1) that
$$e^{}(\alpha _n^i(B),\overline{\alpha }_n^i(B))e^{}=(\alpha _n^i(0),\overline{\alpha }_n^i(0)).$$
Also, the $`SL(2,\text{ }\mathrm{C})`$ invariant Fock vacuum depends upon the background $`B`$-field. Since the Fock vacuum is defined by $`(\alpha _n^i(B),\overline{\alpha }_n^i(B))|0;0_B=0`$ for $`n1`$, it is easy to verify that
$$|0;0_B=e^{}|0;0_0.$$
Here, we also labeled the Fock vacua with the letter $`B`$ to show its dependence on the $`B`$-field.
In the next subsection, we will try to find the string field theory in the background $`B`$-field by gradually increasing the amount of the background $`B`$-field. We will do this by taking advantage of the condensation of closed string fields. This condensation turns out to only affect the kinetic term, i.e., the BRS charge, but not the interaction terms. Therefore, in this paper, we will assume that the string vertices are background independent. This is also supported by the fact that the string vertices are specified up to an overall normalization factor by the overlapping conditions which are written in terms of the string coordinates $`X^i(\sigma )`$ and their momenta $`P_i(\sigma )`$.
3.2. The Calculation for the Shifting of Closed String Fields
An infinitesimal vacuum expectation value of the $`B`$-field contained in closed string fields $`|\mathrm{\Phi }`$ is given by
$$|\delta \mathrm{\Phi }=\frac{1}{x_cg^2}c_0^{}\delta B_{ij}\alpha _1^i(B)\overline{\alpha }_1^j(B)|1,1_B(2\pi )^{26}\delta ^{26}(p),$$
where $`\delta B_{ij}`$ is the infinitesimal constant background of the $`B`$-field and $`x_c`$ is the numerical factor included in the closed three-string vertex. In addition, $`c_0^{}`$ is the zero mode of the reparametrization ghost and is given by the combination $`c_0\overline{c}_0`$. $`|1,1_B`$ denotes the state $`c_1\overline{c}_1|0,0_B`$. By shifting closed string fields as $`|\mathrm{\Phi }|\mathrm{\Phi }+|\delta \mathrm{\Phi }`$ in (2.1), we can obtain an open-closed string field theory in the background $`B`$-field. To calculate this shifting in (2.1), we will follow the prescription given by LeClair, Peskin, and Preitschopf, which will facilitate our calculation.
The essential point of the prescription is to take advantage of conformal field theory. The LPP vertices give the relation between vertices in string field theory and correlation functions in the corresponding conformal field theory, as we will see below. In the action (2.1), in terms of the LPP vertices $`v(n,\mathrm{},2,1)|`$, we can write the five vertices as
$$\begin{array}{cc}& V_3^\mathrm{o}(1,2,3)|=v_3^\mathrm{o}(1,2,3)|,\hfill \\ & V_3^\mathrm{c}(1^\mathrm{c},2^\mathrm{c},3^\mathrm{c})|=x_cv_3^\mathrm{c}(1^\mathrm{c},2^\mathrm{c},3^\mathrm{c})|\underset{r=1^\mathrm{c},2^\mathrm{c},3^\mathrm{c}}{}(b_0^{}𝒫)^{(r)},\hfill \\ & U(1,2^\mathrm{c})|=x_uu(1,2^\mathrm{c})|(b_0^{}𝒫)^{(2^\mathrm{c})},\hfill \\ & V_4^\mathrm{o}(1,2,3,4)|=x_4_{\sigma _i}^{\sigma _f}𝑑\sigma _0^{(1)}v_4^\mathrm{o}(1,2,3,4;\sigma _0^{(1)})|b_{\sigma _0^{(1)}},\hfill \\ & U_\mathrm{\Omega }(1,2,3^\mathrm{c})|=x_\mathrm{\Omega }_{\sigma _i}^{\sigma _f}𝑑\sigma _0^{(1)}u_\mathrm{\Omega }(1,2,3^\mathrm{c};\sigma _0^{(1)})|b_{\sigma _0^{(1)}}(b_0^{}𝒫)^{(3^\mathrm{c})},\hfill \end{array}$$
where the vertices denoted by the small letters are the LPP vertices. The parameter $`\sigma _0`$ is a moduli parameter of the vertices and $`𝒫`$ is the operator which projects onto the $`L_0=\overline{L}_0`$ sector for closed strings. The necessity of the inserted anti-ghosts operator $`b`$ is explained in \[19,,6\]. In particular, $`b_0^{}`$ is given by $`(b_0\overline{b}_0)/2`$.
These LPP vertices are defined by
$$v(n,\mathrm{},2,1)||A_1_1|A_2_2\mathrm{}|A_n_n=h_1[𝒪_{A_1}]h_2[𝒪_{A_2}]\mathrm{}h_n[𝒪_{A_n}]$$
where the right-hand side denotes the correlation function of the operators $`𝒪_{A_r}`$ transformed by the conformal mappings $`h_r`$ in the conformal field theory. Here, the string fields $`|A_r_r`$ are $`𝒪_A|0_r`$, and each string field $`|A_r_r`$ corresponds to a unit disc $`|w_r|1`$ of which the corresponding operator $`𝒪_{A_r}`$ is inserted at the origin. These unit discs are mapped by $`h_r`$ into a single complex $`z`$-plane to glue the strings. In particular, if the operator $`𝒪`$ is a primary field $`\varphi (w,\overline{w})`$ of weight $`(d,\overline{d})`$, it is transformed by the conformal mapping $`z=h(w)`$ into
$$h[\varphi (w,\overline{w})]=\left(\frac{dh(w)}{dw}\right)^d\left(\frac{d\overline{h}(\overline{w})}{d\overline{w}}\right)^{\overline{d}}\varphi (h(w),\overline{h}(\overline{w})).$$
As a simple example, the reflector for open strings is defined by
$$R^\mathrm{o}(2,1)||A_1|B_2=I[𝒪_A]𝒪_B,$$
where $`I(z)=1/z`$ . For our calculation later, it will be useful to introduce another definition of the reflector
$$R^\mathrm{o}(2,1)||A_1|B_2=h_1^R[𝒪_A]h_2^R[𝒪_B].$$
These conformal mappings $`h_r^R(w_r)`$ are given by the Mandelstam mapping $`\rho =\alpha _1\mathrm{ln}(zZ_1)+\alpha _2\mathrm{ln}(zZ_2)`$ where $`\alpha _1+\alpha _2=0`$. Here we use an intermediate $`\rho `$-plane: $`w_r=e^{\rho _r}`$; $`\rho =\alpha _1\rho _1+\rho _0`$ for the first string and $`\rho =\alpha _2\rho _2+\rho _0+i\pi \alpha _1`$ for the second string. $`\rho _0`$ is the interaction point.
Thus, in order to make use of this technique, we have to go to the Euclidean signature on the worldsheet. In the conformal field theory language, the string coordinates of closed strings are given by $`X^i(w)=(l_s/2i)_{n=\mathrm{}}^{\mathrm{}}\alpha _n^iw^{n1}`$ and $`\overline{}X^i(\overline{w})=(l_s/2i)_{n=\mathrm{}}^{\mathrm{}}\overline{\alpha }_n^i\overline{w}^{n1}`$. For the reparametrization ghost, $`c(w)=_{n=\mathrm{}}^{\mathrm{}}c_nw^{n+1}`$ and $`\overline{c}(\overline{w})=_{n=\mathrm{}}^{\mathrm{}}\overline{c}_n\overline{w}^{n+1}`$. In terms of these variables, $`|\delta \mathrm{\Phi }`$ can be written as
$$|\delta \mathrm{\Phi }=c_0^{}\left(\frac{2}{\alpha ^{}x_cg^2}\right)\delta B_{ij}\underset{w0}{lim}c\overline{c}X^i\overline{}X^j(w,\overline{w})|0,0_B(2\pi )^{26}\delta ^{26}(p).$$
To accurately formulate our calculation, it is useful to introduce a ‘regularized’ state
$$|B_ϵ^{ij}_B=\underset{w0}{lim}c\overline{c}X^i\overline{}X^j(w,\overline{w})|0,0_B(2\pi )^{26}\delta _ϵ^{26}(p),$$
where $`\delta _ϵ^{26}(p)=(1/2)[\delta (p^+ϵ)+\delta (p^++ϵ)]\delta ^{25}(p)`$. Therefore, we can see that $`c_0^{}\delta B_{ij}|B_ϵ^{ij}(\alpha ^{}x_cg^2/2)|\delta \mathrm{\Phi }`$ in the limit $`ϵ0`$.
Now, let us calculate the closed three-string vertex $`V_3^\mathrm{c}(1^\mathrm{c},2^\mathrm{c},3^\mathrm{c})|`$ inserted by the shift $`|\delta \mathrm{\Phi }`$ by following the LPP prescription. The conformal mappings $`h_r`$ for the closed three-string vertex are given by the Mandelstam mapping \[12,,9\]
$$\rho =\underset{r=1}{\overset{3}{}}\alpha _r\mathrm{ln}(zZ_r)=\alpha _r\mathrm{ln}w_r+\rho _0^{(r)},$$
where string length parameters satisfy $`_{r=1}^3\alpha _r=0`$. For our calculation, we only need the vertex with a particular configuration of strings where $`\alpha _2,\alpha _3>0`$ and $`\alpha _1<0`$. In the Madelstam mapping (3.1) of the configuration, $`\rho _0^{(1)}=\rho _0+i\pi \alpha _3`$; $`\rho _0^{(2)}=\rho _0i\pi \alpha _2`$; $`\rho _0^{(3)}=\rho _0`$, where $`\rho _0`$ is the interaction point given by $`\rho _0=\rho (z_0)`$; $`z_0`$ is defined by $`(\rho /z)(z=z_0)=0`$. The variable $`Z_r`$ is the position of the operator of the $`r`$-th string on the complex $`z`$-plane and are specified by $`Z_r=h_r(w_r=0)`$.
Using the property of the LPP vertex, we find that
$$v_3^\mathrm{c}(1^\mathrm{c},2^\mathrm{c},3^\mathrm{c})||B_ϵ^{ij}_3|A_2|B_1=c\overline{c}X^i\overline{}X^j(Z_3,\overline{Z}_3)h_2^ϵ[𝒪_A]h_1^ϵ[𝒪_B],$$
where we introduce the label $`ϵ`$ for the above mappings $`h_r`$ to remind us that the light-cone momentum $`p^+`$ of the third string is infinitesimal: $`\alpha _3ϵ`$. In the right-hand side of (3.1), since the operator $`c\overline{c}X^i\overline{}X^j`$ contracted with the antisymmetric tensor $`\delta B_{ij}`$ in the shift is a primary field with its weight $`(0,0)`$ and $`h_3(w_3=0)=Z_3`$, we have used the fact that $`lim_{w0}h_3[c\overline{c}\delta B_{ij}X^i\overline{}X^j(w,\overline{w})]=c\overline{c}\delta B_{ij}X^i\overline{}X^j(Z_3,\overline{Z}_3)`$.
In the limit $`\alpha _30`$, we can see from (3.1) that the mappings $`h_r^ϵ`$ become $`h_r^R`$, which are the conformal mappings to give the reflector for closed strings: the counterpart of the mapping associated with the reflector (3.1) for open strings. The insertion point $`Z_3`$ goes to the interaction point $`z_0`$ on the $`z`$-plane. Therefore, the operator $`c\overline{c}X^i\overline{}X^j(Z_3,\overline{Z}_3)`$ in the right-hand side of (3.1) becomes $`c\overline{c}X^i\overline{}X^j(z_0,\overline{z}_0)`$. Since the interaction $`z_0`$ corresponds to the point $`w_1=1`$ of the first string, we can regard the operator $`c\overline{c}X^i\overline{}X^j(z_0,\overline{z}_0)`$ as $`h_1^R[c\overline{c}X^i\overline{}X^j(w=1,\overline{w}=1)]`$. Thus, we find that
$$v_3^\mathrm{c}(1^\mathrm{c},2^\mathrm{c},3^\mathrm{c})||B_ϵ^{ij}_3=R^\mathrm{c}(1^\mathrm{c},2^\mathrm{c})|\left(c\overline{c}X^i\overline{}X^j\right)^{(1^\mathrm{c})}(w=1,\overline{w}=1),$$
where we dropped the states $`|A_2`$ and $`|B_1`$, when the both sides are contracted with $`\delta B_{ij}`$. Since we made no assumptions in choosing these states, (3.1) should hold identically.
When we contract the closed three-string vertex $`V_3^\mathrm{c}(1^\mathrm{c},2^\mathrm{c},3^\mathrm{c})|`$ with the shift $`|\delta \mathrm{\Phi }`$, we can see that the zero mode $`b_0^{}`$ from the vertex cancels $`c_0^{}`$ from the shift by their commutation relation $`\{b_0^{},c_0^{}\}=1`$ and that the product of $`v_3^\mathrm{c}(1^\mathrm{c},2^\mathrm{c},3^\mathrm{c})||B_ϵ^{ij}_3`$ and $`(𝒫b_0^{})^{(1)}(𝒫b_0^{})^{(2)}`$ are left. By (3.1), the former factor becomes the reflector times the operator $`(2/\alpha ^{}x_cg^2)\delta B_{ij}c\overline{c}X^i\overline{}X^j`$ of the first string. We make the projection $`𝒫`$ and the anti-ghost zero mode $`b_0^{}`$ of the second string pass through the operator $`c\overline{c}X^i\overline{}X^j`$. After contracting the anti-ghost zero mode $`b_0^{}`$ with the ghost operators $`c(\sigma )\overline{c}(\sigma )`$, we find that $`V_3^\mathrm{c}(1^\mathrm{c},2^\mathrm{c},3^\mathrm{c})||\delta \mathrm{\Phi }_3`$ becomes
$$R^\mathrm{c}(1^\mathrm{c},2^\mathrm{c})|\left(\frac{1}{\alpha ^{}g^2}\right)𝒫^{(1^\mathrm{c})}(c^{(1)}+\overline{c}^{(1)})\delta B_{ij}X^{i(1)}\overline{}X^{j(1)}(w=1,\overline{w}=1)(b_0^{}𝒫)^{(1^\mathrm{c})}.$$
Since the projection operator $`𝒫`$ is given by $`(1/2\pi )_0^{2\pi }𝑑\sigma \mathrm{exp}i\sigma (L_0\overline{L}_0)`$, for an operator $`𝒪(w,\overline{w})`$, we obtain
$$𝒫𝒪(w,\overline{w})𝒫=_0^{2\pi }\frac{d\sigma }{2\pi }𝒪(we^{i\sigma },\overline{w}e^{i\sigma })𝒫.$$
By applying (3.1) to (3.1), we find that $`V_3^\mathrm{c}(1^\mathrm{c},2^\mathrm{c},3^\mathrm{c})||\delta \mathrm{\Phi }_3`$ can be rewritten as the reflector $`R^\mathrm{c}(1^\mathrm{c},2^\mathrm{c})|`$ operated by
$$\left(\frac{1}{\alpha ^{}g^2}\right)_0^{2\pi }\frac{d\sigma }{2\pi }(c^{(1)}(\sigma )+\overline{c}^{(1)}(\sigma ))\delta B_{ij}X^{i(1)}\overline{}X^{j(1)}(e^{i\sigma },e^{i\sigma })(b_0^{}𝒫)^{(1^\mathrm{c})},$$
where $`c(\sigma )=e^{i\sigma }c(e^{i\sigma })=_nc_ne^{in\sigma }`$ and $`\overline{c}(\sigma )=e^{i\sigma }\overline{c}(e^{i\sigma })=_n\overline{c}_ne^{in\sigma }`$. Furthermore, $`X^i\overline{}X^j(e^{i\sigma },e^{i\sigma })`$ in (3.1) can be written in terms of the string coordinates and the momenta as $`A^i(\sigma )\overline{A}^j(\sigma )`$ where
$$\begin{array}{cc}\hfill A^i(\sigma )& =\frac{1}{2i}g^{ij}\left[(2\pi \alpha ^{})P_j(\sigma )+\left(g_{jk}b_{jk}\right)X_{}^{}{}_{}{}^{k}(\sigma )\right],\hfill \\ \hfill \overline{A}^i(\sigma )& =\frac{1}{2i}g^{ij}\left[(2\pi \alpha ^{})P_j(\sigma )\left(g_{jk}+b_{jk}\right)X_{}^{}{}_{}{}^{k}(\sigma )\right].\hfill \end{array}$$
Substituting $`A^i(\sigma )\overline{A}^j(\sigma )`$ into (3.1), we find that (3.1) becomes
$$\frac{1}{g^2}\frac{2\pi \alpha ^{}}{(2\alpha ^{})}_0^{2\pi }\frac{d\sigma }{2\pi }(c^{(1)}(\sigma )+\overline{c}^{(1)}(\sigma ))g^{ik}\delta B_{kj}\left(P_i(\sigma )B_{il}X_{}^{}{}_{}{}^{l}(\sigma )\right)X_{}^{}{}_{}{}^{j}(\sigma )(b_0^{}𝒫)^{(1^\mathrm{c})}.$$
As we can see from (3.1), (3.1) can certainly be interpreted as the change of the BRS charge and we obtain the final result of this calculation
$$V_3^\mathrm{c}(1^\mathrm{c},2^\mathrm{c},3^\mathrm{c})||\delta \mathrm{\Phi }=\frac{1}{g^2}R^\mathrm{c}(1^\mathrm{c},2^\mathrm{c})|\delta Q_\mathrm{B}^{(1^\mathrm{c})}(B)(b_0^{}𝒫)^{(1^\mathrm{c})},$$
where we label the variation of the BRS charge $`Q_\mathrm{B}(B)`$ with the letter $`B`$ as $`\delta Q_\mathrm{B}(B)`$ to indicate its dependence on the background field $`B_{ij}`$. This result (3.1) is in agreement with the result in .
The condensation of the $`B`$-field can also change the kinetic term of open string fields. The change comes from the contraction of the open-open-closed vertex $`U_\mathrm{\Omega }|`$ and the shift $`|\delta \mathrm{\Phi }`$. The corresponding LPP vertex $`u_\mathrm{\Omega }(1,2,3^\mathrm{c};\sigma _0)|`$ is specified by the Mandelstam mapping
$$\begin{array}{cc}\hfill \rho & =\underset{r=1,2,3,3^{}}{}\alpha _r\mathrm{ln}(zZ_r)\hfill \\ & =\alpha _r\mathrm{ln}w_r+\rho _0^{(r)},\hfill \end{array}$$
where the string length parameters $`\alpha _r`$ satisfy $`_r\alpha _r=0`$, and $`\alpha _3=\alpha _3^{}`$. When $`\alpha _1,\alpha _30`$ and $`\alpha _20`$, $`\rho _0^{(2)}=\rho _0i\sigma _0^{(1)}i\pi \alpha _2`$; $`\rho _0^{(3)}=\rho _0`$; $`\rho _0^{(1)}=\rho _0i\sigma _0^{(1)}`$ for $`0\sigma ^{(1)}(\sigma _0^{(1)}/\alpha _1)`$ and $`\rho _0^{(1)}=\rho _0i\sigma _0^{(1)}+2i\pi \alpha _3`$ for $`(\sigma _0^{(1)}/\alpha _1)\sigma ^{(1)}\pi `$, where $`\sigma ^{(1)}=\mathrm{Im}\mathrm{ln}w_1`$. $`\rho _0`$ is the interaction point which is given by $`\rho _0=\rho (z_0)`$; $`z_0`$ is defined by $`(\rho /z)(z=z_0)=0`$. $`\sigma _0^{(1)}`$ is the moduli parameter of this vertex which runs from $`0`$ to $`\pi \alpha _1`$, which is the point where the second open string breaks into the first open string and the closed string. Henceforth, for simplicity, $`\sigma _0^{(1)}`$ will be denoted by $`\sigma _0`$. The Koba-Nielsen variables $`Z_1`$, $`Z_2`$ of the open strings are real, while $`Z_3`$ and $`Z_3^{}`$ of the closed string are complex numbers satisfying that $`Z_3^{}=Z_{3}^{}{}_{}{}^{}`$.
Before calculating the contraction $`U_\mathrm{\Omega }(1,2,3^\mathrm{c})||B^{ij}_B`$, it is useful to introduce the open string counterpart of the operator $``$, which will be denoted by $`_o`$: $`_o=(i/2)_0^\pi 𝑑\sigma B_{ij}X_{}^{}{}_{}{}^{i}(\sigma )X^j(\sigma )`$. Since the vertex $`U_\mathrm{\Omega }(1,2,3^\mathrm{c})|`$ in the case where $`\alpha _1,\alpha _30`$ $`\alpha _20`$ satisfies the overlapping condition
$$\begin{array}{cc}& u_\mathrm{\Omega }(1,2,3^\mathrm{c};\sigma _0)|\left[X^{i(1)}\left(\frac{\sigma }{\alpha _1}\right)X^{i(2)}\left(\frac{\pi |\alpha _2|\sigma }{|\alpha _2|}\right)\right]=0,\left(0\sigma <\sigma _0\right)\hfill \\ & u_\mathrm{\Omega }(1,2,3^\mathrm{c};\sigma _0)|\left[X^{i(3)}\left(\frac{\sigma \sigma _0}{\alpha _3}\right)X^{i(2)}\left(\frac{\pi |\alpha _2|\sigma }{|\alpha _2|}\right)\right]=0,\left(\sigma _0\sigma <\sigma _0+2\pi \alpha _3\right)\hfill \\ & u_\mathrm{\Omega }(1,2,3^\mathrm{c};\sigma _0)|\left[X^{i(1)}\left(\frac{\sigma 2\pi \alpha _3}{\alpha _1}\right)X^{i(2)}\left(\frac{\pi |\alpha _2|\sigma }{|\alpha _2|}\right)\right]=0,\left(\sigma _0+2\pi \alpha _3\sigma \pi |\alpha _2|\right)\hfill \end{array}$$
we can see that
$$U_\mathrm{\Omega }(1,2,3^\mathrm{c})|\left[_o^{(1)}+_o^{(2)}+^{(3)}\right]=0.$$
Since, by our assumption, the vertex $`U_\mathrm{\Omega }(1,2,3^\mathrm{c})|`$ is background independent and $`|B^{ij}_B=e^{}|B^{ij}_0`$, we can see from (3.1) that
$$U_\mathrm{\Omega }(1,2,3^\mathrm{c})||B^{ij}_B=U_\mathrm{\Omega }(1,2,3^\mathrm{c})||B^{ij}_0e^{_o^{(1)}_o^{(2)}}.$$
Thus, in the following, we will calculate the contraction $`U_\mathrm{\Omega }(1,2,3^\mathrm{c})||B^{ij}_0`$, where there is no background $`B`$-field. Hence, for the sake of simplicity, we will omit from the states and the oscillators the label $`0`$, which means no background $`B`$-field. After the calculation, we will restore the label.
In a similar way to what we have done for the closed three-string vertex, using the property of the LPP vertex, we find that
$$\begin{array}{cc}& U_\mathrm{\Omega }(1,2,3^\mathrm{c})||\delta \mathrm{\Phi }_{3^\mathrm{c}}|A_2|B_1\hfill \\ & =\underset{\alpha _30}{lim}\frac{2x_\mathrm{\Omega }}{\alpha ^{}x_cg^2}\delta B_{ij}_{\sigma _i}^{\sigma _f}𝑑\sigma _0b_{\sigma _0}c\overline{c}X^i\overline{}X^j(Z_3,Z_{3}^{}{}_{}{}^{})h_2^ϵ[𝒪_A]h_1^ϵ[𝒪_B].\hfill \end{array}$$
Here the anti-ghost factor is given by
$$b_{\sigma _0}=\frac{d\rho _0}{d\sigma _0}_{C_{\rho _0}}\frac{d\rho }{2\pi i}b(\rho )+\frac{d\rho _0^{}}{d\sigma _0}_{C_{\rho _0^{}}}\frac{d\rho ^{}}{2\pi i}\overline{b}(\rho ^{}).$$
We can rewrite the anti-ghost factor as
$$\begin{array}{cc}\hfill b_{\sigma _0}& =i_{C_{z_0}}\frac{dz}{2\pi i}\left(\frac{d\rho }{dz}\right)^1b(z)i_{C_{z_0^{}}}\frac{dz^{}}{2\pi i}\left(\frac{d\rho ^{}}{dz^{}}\right)^1\overline{b}(z^{})\hfill \\ & =2i\left(\frac{d^2\rho }{dz^2}\right)_{z=z_0}^1b(z_0)2i\left(\frac{d^2\rho ^{}}{dz^2}\right)_{z^{}=z_0^{}}^1\overline{b}(z_0^{}),\hfill \end{array}$$
where we have used
$$\rho (z)=\rho (z_0)+\frac{1}{2}\left(\frac{d^2\rho }{dz^2}\right)_{z=z_0}^2(zz_0)^2+\mathrm{O}((zz_0)^3).$$
As we can see from the Mandelstam mapping of (3.1), when $`\alpha _3`$ goes to zero, we can obtain the following equations:
$$\begin{array}{cc}\hfill z_0& =Z_3\alpha _3\frac{(Z_3Z_1)(Z_3Z_2)}{\alpha _1Z_1+\alpha _2Z_2}+\mathrm{O}(\alpha _3^2),\hfill \\ \hfill \left(\frac{d^2\rho }{dz^2}\right)_{z=z_0}^1& =\alpha _3\left[\frac{(Z_3Z_1)(Z_3Z_2)}{\alpha _1Z_1+\alpha _2Z_2}\right]^2+\mathrm{O}(\alpha _3^2).\hfill \end{array}$$
The first equation means that the interaction point $`z_0`$ goes to $`Z_3`$ in the limit $`\alpha _30`$. Then, we can evaluate the ghost part of the correlation function (3.1),
$$\begin{array}{cc}\hfill \underset{\alpha _30}{lim}b_{\sigma _0}c(Z_3)\overline{c}(Z_3^{})& =\underset{\alpha _30}{lim}\mathrm{\hspace{0.17em}2}i\left[\left(\frac{d^2\rho }{dz^2}\right)_{z=z_0}^1\frac{1}{z_0Z_3}\overline{c}(Z_3)+\left(\frac{d^2\rho ^{}}{dz^2}\right)_{z^{}=z_0^{}}^1\frac{1}{z_0^{}Z_3^{}}c(Z_3)\right]\hfill \\ & =2i\left[\frac{(Z_3Z_1)(Z_3Z_2)}{\alpha _1Z_1+\alpha _2Z_2}\overline{c}(Z_3)+\frac{(Z_3^{}Z_1)(Z_3^{}Z_2)}{\alpha _1Z_1+\alpha _2Z_2}c(Z_3)\right]\hfill \\ & =2i\left[\left(\frac{d\stackrel{~}{\rho }}{dz}\right)_{z=Z_3}^1\overline{c}(Z_3^{})+\left(\frac{d\stackrel{~}{\rho }^{}}{dz^{}}\right)_{z^{}=Z_3^{}}^1c(Z_3)\right],\hfill \end{array}$$
where $`\stackrel{~}{\rho }`$ denotes the mapping (3.1) in the limit $`\alpha _30`$. In the same limit, since $`z_0Z_3`$, the insertion point of the operator $`c\overline{c}X^i\overline{}X^j`$ becomes equal to the interaction point $`\rho _0`$ on the $`\rho `$ plane. Because this point corresponds to the point $`\sigma ^{(1)}=(\sigma _0/\alpha _1)`$ on the worldsheet of the first string, we can see that
$$\begin{array}{cc}& \left(\frac{d\stackrel{~}{\rho }}{dz}\right)^1\overline{c}X^i\overline{}X^j(z_0,z_{0}^{}{}_{}{}^{})=\frac{e^{i(\sigma _0/\alpha _1)}}{\alpha _1}h_1\left[\overline{c}X^i\overline{}X^j(e^{i(\sigma _0/\alpha _1)},e^{i(\sigma _0/\alpha _1)})\right],\hfill \\ & \left(\frac{d\stackrel{~}{\rho }^{}}{dz^{}}\right)^1cX^i\overline{}X^j(z_0,z_{0}^{}{}_{}{}^{})=\frac{e^{i(\sigma _0/\alpha _1)}}{\alpha _1}h_1\left[cX^i\overline{}X^j(e^{i(\sigma _0/\alpha _1)},e^{i(\sigma _0/\alpha _1)})\right],\hfill \end{array}$$
where we have used $`(dz/d\stackrel{~}{\rho })=(w_1/\alpha _1)(dz/dw_1)`$ and the anti-holomorphic counterpart, when the both sides are contracted with $`\delta B_{ij}`$.
Substituting (3.1) and (3.1) into (3.1) and scaling $`\sigma _0`$ as $`\sigma _0\alpha _1\sigma _0`$, we find that $`U_\mathrm{\Omega }(1,2,3^\mathrm{c})||\delta \mathrm{\Phi }_3`$$`|A_2|B_1`$ becomes
$$R^o(1,2)\left|\frac{8\pi ix_\mathrm{\Omega }}{\alpha ^{}x_cg^2}\delta B_{ij}_0^\pi \frac{d\sigma _0}{2\pi }(\rho _0^{}c^{(1)}(\rho _0)+\rho _0\overline{c}^{(1)}(\rho _0^{}))X^{i(1)}\overline{}X^{j(1)}(\rho _0,\rho _0^{})\right|A_2|B_1,$$
where $`\rho _0=e^{i\sigma _0}`$. In the paper , it was shown that the BRS invariance of the action requires that $`8\pi ix_\mathrm{\Omega }=x_c`$, which we will apply to (3.1). Since $`X^i(w)=(l_s/2i)_n\alpha _n^iw^{n1}`$ and $`\overline{}X^i(w^{})=(l_s/2i)_n\alpha _n^i(w^{})^{n1}`$, in terms of the string coordinates $`i_\sigma X^i(\sigma )=l_s_{n0}\alpha _n^i\mathrm{sin}(n\sigma )`$ and the momenta $`P_i(\sigma )=(1/\pi l_s)_ng_{ij}\alpha _n^j\mathrm{cos}(n\sigma )`$ for open strings, the operator $`X^i\overline{}X^j(\rho _0,\rho _0^{})`$ turns out to be $`(1/4)\left(2\pi \alpha ^{}g^{ik}P_k+X_{}^{}{}_{}{}^{i}\right)`$ $`\left(2\pi \alpha ^{}g^{jl}P_lX_{}^{}{}_{}{}^{j}\right)`$. Therefore, we can verify that (3.1) becomes
$$R^o(1,2)\left|\frac{1}{g^2}\frac{2\pi \alpha ^{}}{(2\alpha ^{})}_0^\pi \frac{d\sigma }{2\pi }\left(c^{(1)}(\sigma )+\overline{c}^{(1)}(\sigma )\right)g^{ik}\delta B_{kj}P_i^{(1)}(\sigma )X_{}^{}{}_{}{}^{j(1)}(\sigma )\right|A_2|B_1.$$
Here $`\sigma _0`$ has been renamed $`\sigma `$. Thus, as we can see from (3.1), $`U_\mathrm{\Omega }(1,2,3^\mathrm{c})||\delta \mathrm{\Phi }_{B\mathrm{\hspace{0.17em}3}}`$ in the presence of the $`B`$-field becomes
$$R^o(1,2)|\frac{1}{g^2}e^{_o^{(1)}}_0^\pi 𝑑\sigma \frac{1}{2}\left(c^{(1)}(\sigma )+\overline{c}^{(1)}(\sigma )\right)g^{ik}\delta B_{kj}P_i^{(1)}(\sigma )X_{}^{}{}_{}{}^{j(1)}(\sigma )e^{_o^{(1)}},$$
where the operator $`e^{_o^{(2)}}`$ have been passed through the operator acting on the reflector $`R^o(1,2)|`$ in (3.1) and made into $`e^{_o^{(1)}}`$ by using the reflector.
Now, let us show that the operator
$$e^_o_0^\pi 𝑑\sigma \left(c(\sigma )+\overline{c}(\sigma )\right)g^{ik}\delta B_{kj}P_i(\sigma )X_{}^{}{}_{}{}^{j}(\sigma )e^_o$$
is equal to the variation of the BRS charge for the open strings. We have
$$[_o,P_i(\sigma )]=\frac{1}{2}B_{ij}\left[\frac{}{\sigma }_0^\pi 𝑑\sigma ^{}\delta _s(\sigma ,\sigma ^{})X^j(\sigma ^{})+_0^\pi 𝑑\sigma ^{}\delta _c(\sigma ,\sigma ^{})X_{}^{}{}_{}{}^{j}(\sigma ^{})\right],$$
where we define that $`\delta _s(\sigma ,\sigma ^{})=(2/\pi )_{n=1}^{\mathrm{}}\mathrm{sin}(n\sigma )\mathrm{sin}(n\sigma ^{})`$ and that $`\delta _c(\sigma ,\sigma ^{})=(1/\pi )+(2/\pi )_{n=1}^{\mathrm{}}\mathrm{cos}(n\sigma )\mathrm{cos}(n\sigma ^{})`$. We can also prove that
$$_0^\pi 𝑑\sigma 𝑑\sigma ^{}f(\sigma )\delta _c(\sigma ,\sigma ^{})g(\sigma ^{})=_0^\pi 𝑑\sigma 𝑑\sigma ^{}f(\sigma )\delta _s(\sigma ,\sigma ^{})g(\sigma ^{})$$
where $`f(\sigma )`$ and $`g(\sigma )`$ are $`\mathrm{cos}(n\sigma )`$ or $`\mathrm{sin}(n\sigma )`$ with $`n0`$. Furthermore, even when either $`f(\sigma )`$ or $`g(\sigma )`$ is $`\sigma `$, (3.1) still holds. By using (3.1) and (3.1), we find that (3.1) becomes
$$_0^\pi 𝑑\sigma \left(c(\sigma )+\overline{c}(\sigma )\right)g^{ik}\delta B_{kj}\left[P_i(\sigma )B_{il}X_{}^{}{}_{}{}^{l}(\sigma )\right]X_{}^{}{}_{}{}^{j}(\sigma ),$$
which is equal to $`\delta Q_\mathrm{B}(B)`$; $`\delta Q_\mathrm{B}(B)`$ is the variation of the BRS charge $`Q_\mathrm{B}(B)`$ under $`B_{ij}B_{ij}+\delta B_{ij}`$. Here we assume that the BRS charge for open strings is defined as
$$Q_\mathrm{B}(B)=_0^\pi \frac{d\sigma }{2\pi }\left[c(\sigma )T_{++}(\sigma )+\overline{c}(\sigma )T_{}(\sigma )\right],$$
where the energy-momentum tensor $`T_{\pm \pm }`$ is
$$\frac{1}{4\alpha ^{}}g^{ij}\left[(2\pi \alpha ^{})P_i\pm \left(g_{ik}b_{ik}\right)X^k\right]\left[(2\pi \alpha ^{})P_j\pm \left(g_{jl}b_{jl}\right)X^l\right].$$
Note that this BRS charge $`Q_\mathrm{B}(B)`$ is different from the BRS charge $`\widehat{Q}_B`$ in the previous section in the usage of $`X^i(\sigma )`$ instead of $`\widehat{X}^i(\sigma )`$. Thus, we obtain that
$$U_\mathrm{\Omega }(1,2,3^\mathrm{c})||\delta \mathrm{\Phi }_B=\frac{1}{g^2}R^o(1,2)|\delta Q_\mathrm{B}^{(1)}(B).$$
At last, we can find another string field theory in the finite $`B`$-field by making use of (3.1) and (3.1):
$$\begin{array}{cc}\hfill S(B)=& \frac{1}{2}\mathrm{\Psi }\left|Q_\mathrm{B}(B)\right|\mathrm{\Psi }\frac{1}{2}\mathrm{\Phi }\left|Q_\mathrm{B}(B)(b_0^{}𝒫)\right|\mathrm{\Phi }\hfill \\ & +\frac{g}{3}V_3^\mathrm{o}(1,2,3)||\mathrm{\Psi }_{321}+\frac{g^2}{4}V_4^\mathrm{o}(1,2,3,4)||\mathrm{\Psi }_{4321}+\frac{g^2}{3!}V_3^\mathrm{c}(1^\mathrm{c},2^\mathrm{c},3^\mathrm{c})||\mathrm{\Phi }_{321}\hfill \\ & +gU(1,2^\mathrm{c})||\mathrm{\Phi }_2|\mathrm{\Psi }_1+\frac{g^2}{2}U_\mathrm{\Omega }(1,2,3^\mathrm{c})||\mathrm{\Phi }_3|\mathrm{\Psi }_{21},\hfill \end{array}$$
where $`Q_\mathrm{B}(B)`$ for closed strings is given by (3.1) and for open strings by (3.1).
4. The Unitary Transformation between the Two SFTs
In the previous sections we have constructed two open-closed string field theories: one, given in section $`2`$, has the noncommutative factor $`\mathrm{exp}[_{r<s}(i/2)\theta ^{ij}p_i^{(r)}p_j^{(s)}]`$ in the open three- and four-string vertices and the unitary operators $`e^M`$ put on the legs of open string fields in the other vertices, while the other, shown in section $`3`$, has the same vertices as the theory in the absence of the $`B`$-field. The string coordinates $`\widehat{X}^i(\sigma )`$ in the BRS charge $`\widehat{Q}_\mathrm{B}`$ of the former theory is replaced with the background independent coordinates $`X^i(\sigma )`$ in the BRS charge $`Q_\mathrm{B}`$ of the latter theory. Thus, since we have two descriptions of the same physics, we are led to ask what relation exists between these theories.
Henceforth, we will call the former theory in section $`2`$ ‘noncommutative’ theory and the latter in section $`3`$ ‘commutative’ theory. As we have shown in section 2, in the noncommutative theory, each open string field has the unitary operator $`e^M`$. If the open string fields $`\widehat{\mathrm{\Psi }}`$ are transformed as
$$|\mathrm{\Psi }=e^M\widehat{|\mathrm{\Psi }},$$
all the interactions in the noncommutative theory become the same as those in the commutative theory. Then, the BRS charge $`\widehat{Q}_\mathrm{B}`$ in the noncommutative theory becomes $`e^M\widehat{Q}_\mathrm{B}e^M`$. We will show below that the transformed BRS charge $`e^M\widehat{Q}_\mathrm{B}e^M`$ is equal to the BRS charge $`Q_\mathrm{B}`$ in the commutative theory. Thus, by the unitary transformation, these theories are related to each other.
By the unitary transformation, the ‘noncommutative’ string coordinates $`_\sigma \widehat{X}^i(\sigma )`$ turn into
$$e^M_\sigma \widehat{X}^i(\sigma )e^M=_\sigma X^i(\sigma )+\theta ^{ij}_0^\pi 𝑑\sigma ^{}\left[\delta _c(\sigma ,\sigma ^{})\delta _s(\sigma ,\sigma ^{})\right]P_j(\sigma ^{}).$$
The second term in the right-hand side of (4.1) can be seen from (3.1) to be vanishing, when it is included in an integrand. Therefore, in the transformed BRS charge $`e^M\widehat{Q}_\mathrm{B}e^M`$, all we have to do is to replace the string coordinates $`\widehat{X}^i(\sigma )`$ with the ordinary coordinates $`X^i(\sigma )`$. This means that $`e^M\widehat{Q}_\mathrm{B}e^M`$ is equal to the BRS charge $`Q_\mathrm{B}`$ in the commutative theory. Thus, the noncommutative theory is mapped to the commutative theory by the unitary transformation (4.1).
5. Discussion
In this paper, we have obtained two open-closed string field theories in the same background. One gives a noncommutative gauge theory in the low-energy limit, while the other is expected to give an ordinary gauge theory with the constant background field strength $`B_{ij}`$ in the infrared. In section $`4`$, we have shown that these theories can be transformed into each other by the unitary operator.
As we mentioned in the footnote in section $`2`$, our oriented open-closed string field theory does not maintain the full gauge invariance. Since we need worldsheets to be orientable to include the background $`B`$-field, it would be preferable to have an open-closed string field theory in an appropriate curved background to maintain the gauge invariance by using the Fischler-Susskind mechanism , while the orientability of worldsheets is kept, as has been argued in . Besides this issue, the condensation of closed string fields we have dealt with only satisfies the equation of motion from the string field theory up to the next leading order of the string coupling constant $`g`$, mainly because, at least at present, we do not have the complete action of a fully gauge invariant oriented open-closed string field theory. However, everything we have shown in this paper would remain true at the leading order of the coupling constant $`g`$, even in such a gauge invariant string field theory.
This paper was motivated by the papers \[1,,3\]. Since the kinetic term in Witten’s string field theory is the same as the kinetic term in our string field theory in this paper, it is natural to think that Witten’s theory transformed by the unitary operator could give an ordinary gauge theory with the background field strength in the low-energy limit. Although, in Witten’s open string field theory, we do not have explicit closed string fields, by using the above transformed theory, we might be able to discuss an infinitesimal variation of the background $`B`$-field in a similar way to the discussion in . This may lead us to the understanding of the gauge symmetry $`B_{ij}B_{ij}+_i\mathrm{\Lambda }_j_j\mathrm{\Lambda }_i`$, $`A_iA_i+\mathrm{\Lambda }_i`$ in Witten’s string field theory.
Acknowledgements
The authors would like to thank Tsuguhiko Asakawa, Taichiro Kugo, Katsumi Itoh, Masahiro Maeno, Kazumi Okuyama, Kazuhiko Suehiro, and Seiji Terashima for valuable discussion. T. K. is grateful to the organizers and the participants of Summer Institute ’99 in Yamanashi, Japan for the hospitality and stimulating atmosphere created there, which helped to initiate this work. T. K. was supported in part by a Grant-in-Aid (#11740143) and in part by a Grant-in-Aid for Scientific Research in a Priority Area: “Supersymmetry and Unified Theory of Elementary Particles”(#707), from the Ministry of Education, Science, Sports and Culture. T. T. was supported in part by Research Fellowships of the Japan Society for the Promotion of Science for Young Scientists.
Appendix A. The Derivations of the Vertices in the Background $`B`$-Field
In this appendix, we will derive the open-closed transition vertex $`\widehat{U}(1,2^\mathrm{c})|`$ and the open-open-closed string vertex $`\widehat{U}_\mathrm{\Omega }(1,2,3^\mathrm{c})|`$ in the presence of the background $`B`$-field.
We begin with the open-closed transition vertex $`\widehat{U}(1,2^\mathrm{c})|`$, which is specified up to an overall normalization by the overlapping condition
$$\widehat{U}(2,1^\mathrm{c})|\left\{\varphi ^{(1^\mathrm{c})}(\sigma ^{(1)})\widehat{\varphi }^{(2)}(\sigma ^{(2)})\right\}=0,(0\sigma 2\pi |\alpha _1|)$$
where $`2\alpha _1+\alpha _2=0`$ and $`\sigma ^{(1)}=\sigma /|\alpha _1|`$; $`\sigma ^{(2)}=(\pi |\alpha _2|\sigma )/|\alpha _2|`$. Here $`\varphi ^{(1^\mathrm{c})}`$ denotes $`X^{i(1)}`$ or $`P_i^{(1)}`$ of the first closed string and $`\widehat{\varphi }^{(2)}`$ denotes $`\widehat{X}^{i(2)}`$ or $`P_i^{(2)}`$ of the second open string. Similarly, the open-closed transition vertex $`U(1,2^\mathrm{c})|`$ in the absence of the $`B`$-field is specified by
$$U(2,1^\mathrm{c})|\{\varphi ^{(1^\mathrm{c})}(\sigma ^{(1)})\varphi ^{(2)}(\sigma ^{(2)})\}=0.(0\sigma 2\pi |\alpha _1|)$$
Here, $`\varphi ^{(2)}`$ denotes $`X^{i(2)}`$ or $`P_i^{(2)}`$ of the second open string.
Recalling the definition of the operators $`M`$ for open strings from (2.1) and for closed strings from (2.1), by making use of (A.1) we can verify that
$$U(2,1^\mathrm{c})|\left(M^{(1^\mathrm{c})}+M^{(2)}\right)=0.$$
We define the closed string counterpart of the ‘dual’ coordinates $`Q_j(\sigma )`$ as
$$Q_j(\sigma )=\frac{1}{2}_0^{2\pi }𝑑\sigma ^{}ϵ(\sigma \sigma ^{})P_j(\sigma ^{}).$$
The dual coordinates $`Q_j(\sigma )`$ for closed strings connect to those for open strings on the open-closed transition vertex $`U(2,1^\mathrm{c})|`$ as
$$U(2,1^\mathrm{c})|\{Q_j^{(1^\mathrm{c})}(\sigma ^{(1)})Q_j^{(2)}(\sigma ^{(2)})\}=0.(0\sigma 2\pi |\alpha _1|)$$
The commutation relation between $`M`$ and the string coordinates $`X^i(\sigma )`$ for closed strings turns out to be
$$[M,X^i(\sigma )]=\theta ^{ij}Q_j(\sigma ).$$
From (A.1), we find that
$$U(2,1^\mathrm{c})|e^{M^{(1^\mathrm{c})}}X^{i(1^\mathrm{c})}(\sigma ^{(1)})=U(2,1^\mathrm{c})|e^{M^{(1^\mathrm{c})}}\widehat{X}^{i(2)}(\sigma ^{(2)}),(0\sigma 2\pi |\alpha _1|)$$
and, thus, we obtain
$$\widehat{U}(2,1^\mathrm{c})|=U(2,1^\mathrm{c})|e^{M^{(2)}}.$$
Next we will move to the open-open-closed string interaction $`\widehat{U}_\mathrm{\Omega }(3,2,1^\mathrm{c})|`$ in the presence of the $`B`$-field. The overlapping condition for the case $`\alpha _1,\alpha _30`$ and $`\alpha _20`$ is given by
$$\begin{array}{cc}& \widehat{u}_\mathrm{\Omega }(1,2,3^\mathrm{c};\sigma _0)|\left[\widehat{\varphi }^{i(1)}(\sigma ^{(1)})\widehat{\varphi }^{i(2)}(\sigma ^{(2)})\right]=0,\left(0\sigma <\sigma _0\right)\hfill \\ & \widehat{u}_\mathrm{\Omega }(1,2,3^\mathrm{c};\sigma _0)|\left[\varphi ^{i(3^\mathrm{c})}(\sigma ^{(3)})\widehat{\varphi }^{i(2)}(\sigma ^{(2)})\right]=0,\left(\sigma _0\sigma <\sigma _0+2\pi \alpha _3\right)\hfill \\ & \widehat{u}_\mathrm{\Omega }(1,2,3^\mathrm{c};\sigma _0)|\left[\widehat{\varphi }^{i(1)}\left(\sigma ^{(1)}2\pi \frac{\alpha _3}{\alpha _1}\right)\widehat{\varphi }^{i(2)}(\sigma ^{(2)})\right]=0,\left(\sigma _0+2\pi \alpha _3\sigma \pi |\alpha _2|\right)\hfill \end{array}$$
where the notations for $`\widehat{\varphi }^{i(r)}`$ are the same as those for the vertex $`\widehat{U}(2,1^\mathrm{c})|`$, except $`\sigma ^{(1)}=\sigma /\alpha _1`$; $`\sigma ^{(2)}=(\pi |\alpha _2|\sigma )/\alpha _2`$; $`\sigma ^{(3)}=(\sigma \sigma _0)/\alpha _3`$. The overlapping condition for the vertex $`U_\mathrm{\Omega }(3,2,1^\mathrm{c})|`$ in the absence of the $`B`$-field is given by
$$\begin{array}{cc}& u_\mathrm{\Omega }(1,2,3^\mathrm{c};\sigma _0)|\left[\varphi ^{i(1)}(\sigma ^{(1)})\varphi ^{i(2)}(\sigma ^{(2)})\right]=0,\left(0\sigma <\sigma _0\right)\hfill \\ & u_\mathrm{\Omega }(1,2,3^\mathrm{c};\sigma _0)|\left[\varphi ^{i(3^\mathrm{c})}(\sigma ^{(3)})\varphi ^{i(2)}(\sigma ^{(2)})\right]=0,\left(\sigma _0\sigma <\sigma _0+2\pi \alpha _3\right)\hfill \\ & u_\mathrm{\Omega }(1,2,3^\mathrm{c};\sigma _0)|\left[\varphi ^{i(1)}(\sigma ^{(1)}2\pi \frac{\alpha _3}{\alpha _1})\varphi ^{i(2)}(\sigma ^{(2)})\right]=0.\left(\sigma _0+2\pi \alpha _3\sigma \pi |\alpha _2|\right)\hfill \end{array}$$
As we can see from (A.1), on the vertex $`\widehat{U}_\mathrm{\Omega }(1,2,3^\mathrm{c})|`$, we find that
$$U_\mathrm{\Omega }(1,2,3^\mathrm{c})|\left\{M^{(1)}+M^{(2)}+M^{(3)}\right\}=U_\mathrm{\Omega }(1,2,3^\mathrm{c})|i\theta ^{ij}p_i^{(3)}Q_j^{(1)}\left(\frac{\sigma _0}{|\alpha _1|}\right).$$
In addition, we can verify that
$$\begin{array}{cc}& u_\mathrm{\Omega }(\sigma _0)|\left[Q_j^{(1)}(\sigma ^{(1)})Q_j^{(2)}(\sigma ^{(2)})\right]=u_\mathrm{\Omega }(\sigma _0)|\frac{1}{2}p^{(3)},\left(0\sigma <\sigma _0\right)\hfill \\ & u_\mathrm{\Omega }(\sigma _0)|\left[Q_j^{(3^\mathrm{c})}(\sigma ^{(3)})Q_j^{(2)}(\sigma ^{(2)})\right]=u_\mathrm{\Omega }(\sigma _0)|Q_j^{(1)}\left(\frac{\sigma _0}{|\alpha _1|}\right),\left(\sigma _0\sigma <\sigma _0+2\pi \alpha _3\right)\hfill \\ & u_\mathrm{\Omega }(\sigma _0)|\left[Q_j^{(1)}(\sigma ^{(1)}2\pi \frac{\alpha _3}{\alpha _1})Q_j^{(2)}(\sigma ^{(2)})\right]=u_\mathrm{\Omega }(\sigma _0)|\frac{1}{2}p_j^{(3)},\left(\sigma _0+2\pi \alpha _3\sigma \pi |\alpha _2|\right)\hfill \end{array}$$
where $`u_\mathrm{\Omega }(\sigma _0)|`$ is short for $`u_\mathrm{\Omega }(1,2,3^\mathrm{c};\sigma _0)|`$. By utilizing (A.1) and (A.1), we can find that
$$\widehat{U}_\mathrm{\Omega }(1,2,3^\mathrm{c})|=U_\mathrm{\Omega }(1,2,3^\mathrm{c})|e^{M^{(1)}}e^{M^{(2)}}.$$
References
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# Constraints on the Mass and Mixing of the 4th Generation Quark From Direct CP Violation ϵ'/ϵ and Rare 𝐾 Decays
## 1 Introduction
Although the Standard Model (SM) is very successful for explaining the particle physics experiments, it has to face the difficulties of many interesting open questions, such asCP violation. The new experimental results for $`ϵ^{}/ϵ`$, which measures direct CP violation in $`K\pi \pi `$ decays, have been reported by KTeV collaboration at Fermilab and NA48 collaboration at CERN,
$`\text{Re}(ϵ^{}/ϵ)`$ $`=`$ $`(28.0\pm 4.1)10^4\mathrm{KTeV},`$ (1)
$`\text{Re}(ϵ^{}/ϵ)`$ $`=`$ $`(18.5\pm 7.3)10^3\mathrm{NA48},`$ (2)
while the new world average reads:
$`\text{Re}(ϵ^{}/ϵ)`$ $`=`$ $`(21.1\pm 4.6)10^4`$ (3)
This establishment of direct CP violation rules out old superweak models. Yet while the SM predicts a non-vanishing $`ϵ^{}/ϵ`$, the values in (1), (2) and (3) exceed most theoretical predictions of SM. People have to face and resolve this discrepancy. Some possibilities to accommodate the data in SM have been pointed out .
The SM makes precise assumptions on the mechanism that generates the CP violation. The only source of CP violating phase originates from the elements $`V_{u_id_j}`$ of the CKM matrix with three quark generations. In SM, there are both indirect ($`ϵ`$) and direct ($`ϵ^{}`$) CP violation.The analysis of $`ϵ^{}/ϵ`$ can be divided into the short-distance (perturbative) part and long-distance ( non-perturbative)part. Using the effective Hamiltonian, ($`_W=_iC_i(\mu )Q_i(\mu )`$), one can obtain an expression of $`ϵ^{}/ϵ`$ that involves CKM parameters ($`V_{u_id_j}`$), Wilson coefficients ($`y_i`$) and local operator matrix elements ($`Q_i_I`$). The source of most theoretical uncertainties for $`ϵ^{}/ϵ`$ is mainly from the difficulty in calculating non-perturbative part (local operator matrix elements), comparing with the phenomenological determination of CKM parameters and the calculation of the Wilson coefficients at a NLO leval. For $`ϵ^{}/ϵ`$, one of the goals of SM is to determine the hadronic matrix elements.
The interesting in this note is not in this non-perturbative part but the new effects with the fourth sequential generation particles in the short-distance part.Except for the SM explanation, there are many directions in the search for New Physics beyond the SM to resolve CP violation. Unlike SM, almost any extension of SM has,in general, new CP violating phases. That is to say, they give new CP violation sources. The new physics on CP violation beyond SM includes CP violation in Supersymmetry models and extensions of fermion sector, scalar sector and gauge sector of SM. In extensions of fermion sector, there are many models, such as vector-like quark models, sterile neutrino models, proposed for probing new effects on CP violation.
In this note, like in ref., we consider a sequential fourth generation model, in which an up-type quark $`t^{^{}}`$, adown-type quark $`b^{^{}}`$, a lepton $`\tau ^{^{}}`$, and a heavy neutrino $`\nu ^{^{}}`$ are added into the SM. The properties of these new fermions are all the same as their corresponding counterparts of other three generations except their masses and CKM mixing, see tab.1,
As the SM does not fix the number of generations, so far we don’t know why there is more than one generation and what law of Nature determines their number. On the one hand, the purely sequential 4th generation is constrained, even excluded in many literatures. For example, in refs. the method of Padé approximates is used to show that for a large fermion mass, it is possible to dynamically generate $`p`$-wave resonance and then the S parameter bound can serve to exclude a heavy fourth generation of fermions. Ref. found that there is no violation of the S parameter upper bound for any value of the heavy fermion mass and that elastic unitarity, imposed as a constraint on strong $`W_LW_L`$ scattering, yields no information concerning and sheds no light on the existence of a heavy fourth generation. The ref. compared various precision determinations of the Femi constant $`G_F`$ to get the rather stringent bound of 3rd and 4th generation lepton mixing angle $`\theta _{34}`$. It founds that the fourth charged lepton is too heavy and seems non-existence. The precision electroweak measurements can also give the strong constrains to the sequential 4th generation, in particlar the S parameter excludes it to 99.8% CL if is degenerate, and if not and a small T parameter is allowed then it is excluded to 98.2% CL .
However, on the other hand, experimentally, the LEP determinations of the invisible partial decay width of the $`Z^0`$ gauge boson only show that there are certainly three light neutrinos of the usual type with mass less than $`M_Z/2`$ . But the existence of the fourth generation with a heavy neutrino, i.e., $`m_{\nu _4}M_Z/2`$ is not yet excluded. Perhaps there exists some more deep or mechanism to give the room of the sequential fourth generation. Because we really don’t know why there only three generations. So, it is not invalueble to research these new generation as one of the new physics. Before having a more fundamental reason for three generations, one may investigate phenomenologically whether the existing experimental data allow the existence of the fourth generation. This is also the main perpuse of this note. There are a number of papers for discussing the fourth generation phenomena.
In our previous paper, we have investigated the constraints on the fourth generation from the inclusive decays of $`BX_sl^+l^{}`$ and $`BX_s\gamma `$. In this note, we further study its effects on direct CP-violating parameter $`ϵ^{}/ϵ`$ in $`K\pi \pi `$ decays as well as possible new constraints from $`ϵ^{}/ϵ`$ and rare K decays.We limit ourselves to the non-SUSY case in order to concentrate on the phenomenological implication of the fourth generation and will call this model as SM4 hereafter for thesake of simplicity.
CP-violating parameter $`ϵ^{}/ϵ`$ is a short distance dominated process and is sensitive to new physics. In SM4 model, there are not new operators produced. The new particle involved is only the fourth generation up-type quark$`t^{}`$. The heavy mass of $`t^{}`$ propagating in the loop diagrams of penguin and box enters the Wilson coefficients $`y_i`$, as well as top quark $`t`$ and $`W`$ boson. The effects of the fourth generation particles can only modify $`y_i`$. Each new Wilson coefficients$`y_i^{\mathrm{new}}(\mu )`$ is the sum of $`y_i^{\mathrm{sm}}(\mu )`$ and $`y_i^{(4)}(\mu )`$ contributed by t and $`t^{}`$ correspondingly. We can get $`y_t^{}^{(4)}`$ by taking the mass of $`t^{}`$ as one of the input parameter. Moreover, for obtaining $`ϵ^{}/ϵ`$ in SM4,we must know something about elements $`V_{t^{}d_j}`$ of the fourth generation $`4\times 4`$ CKM matrix which now contains nineparemeters, i.e., six angles and three phases.But there are no any direct experimental measurements of them. So we have to get their information indirectly from some meson decays. We investigate three rare $`K`$ decays, $`K^+\pi ^+\nu \overline{\nu }`$, $`K_L\pi ^0\nu \overline{\nu }`$ and $`K_L\mu ^+\mu ^{}`$, in SM4. These decays can give the constraint of the fourth CKM factor, $`\text{Im}V_{t^{^{}}s}^{}V_{t^{^{}}d}`$ (or $`V_{t^{^{}}s}^{}V_{t^{^{}}d}`$ and a fourth generation phase $`\theta `$), which is need for calculating $`(ϵ^{}/ϵ)^{\mathrm{new}}`$. We shall take it as an additional input parameter. As a consequence, the total $`(ϵ^{}/ϵ)`$ is the sum of $`(ϵ^{}/ϵ)^{\mathrm{sm}}`$ and $`(ϵ^{}/ϵ)^4`$contributed by the SM and the new particle $`t^{}`$ correspondingly. Unlike the SM, when taking the central values of all parameters for $`ϵ^{}/ϵ`$, the new value of $`(ϵ^{}/ϵ)`$ can reach the range of the current experimental results whatever values of the non-perturbertion part, hadronic matrix elements, are taken in all known cases. Also, the experimental values of $`(ϵ^{}/ϵ)^{\mathrm{exp}}`$ impose strong constraints on the parameter space of $`\text{Im}V_{t^{}s}^{}V_{t^{}d}`$and $`m_t^{}`$.
In sec. 2, we give the basic formulae for $`ϵ^{}/ϵ`$ with the fourth up-like quark $`t^{}`$ in SM4. In sec. 3, we analyze the constraints on the fourth generation CKM matrix factor $`\text{Im}V_{t^{^{}}s}^{}V_{t^{^{}}d}`$ which is necessary for calculating$`ϵ^{}/ϵ`$ in SM4. Sec. 4 is devoted to the numerical analysis. Finally, in sec. 5, we give our conclusion.
## 2 Basic formulae for $`ϵ^{}/ϵ`$ and Wilson Coefficients $`y_i^{(4)}(\mu )`$ in SM4
The essential theoretical tool for the calculation of $`ϵ^{}/ϵ`$ is the $`\mathrm{\Delta }S=1`$ effective Hamiltonian,
$$_W=\underset{i}{}\frac{G_F}{\sqrt{2}}V_{ud}V_{us}^{}\left[z_i(\mu )+\tau y_i(\mu )\right]Q_i(\mu )$$
(4)
with $`\tau =V_{ts}^{}V_{td}/(V_{us}^{}V_{ud})`$. The direct CP violation in $`K\pi \pi `$ is described by $`ϵ^{}`$. The parameter $`ϵ^{}`$ is given in terms of the amplitudes $`A_0A(k(\pi \pi )_{I=0})`$ and $`A_2A(k(\pi \pi )_{I=2})`$ as follows
$$ϵ^{}=\frac{1}{\sqrt{2}}\xi (1\mathrm{\Omega })\mathrm{exp}(i\mathrm{\Phi }),$$
(5)
where
$$\xi =\frac{\text{Im}A_0}{\text{Re}A_0},\omega =\frac{\text{Re}A_2}{\text{Re}A_0},\mathrm{\Omega }=\frac{1}{\omega }\frac{\text{Im}A_2}{\text{Im}A_0}$$
(6)
and $`\mathrm{\Phi }=\pi /2+\delta _2\delta _0\pi /4`$. With the effective Hamiltonian (4), we can cast (5) into the form
$$\frac{ϵ^{}}{ϵ}=\text{Im}\lambda _t\left[P^{(1/2)}P^{(3/2)}\right],$$
(7)
where
$`P^{(1/2)}`$ $`=`$ $`{\displaystyle P_i^{(1/2)}}=r{\displaystyle y_iQ_i_0(1\mathrm{\Omega }_{\eta +\eta ^{}})}`$ (8)
$`P^{(3/2)}`$ $`=`$ $`{\displaystyle P_i^{(3/2)}}={\displaystyle \frac{r}{\omega }}{\displaystyle y_iQ_i_2}.`$ (9)
with $`r=G_F\omega /(2|\text{Re}A_0)|)`$. $`y_i`$ are the Wilson coefficients and the hadronic matrix elements are
$$Q_i_I(\pi \pi )_I|Q_i|K$$
(10)
The operators $`Q_i`$ and $`Q_i_I`$ are given explicitly in many reviews
When including the contributions from the fourth generation up-type quark $`t^{}`$, the above equations will be modified. The corresponding effective Hamiltonian can be expressed as
$$_W=\underset{i}{}\frac{G_F}{\sqrt{2}}V_{ud}V_{us}^{}\left[z_i(\mu )+\tau y_i^{\mathrm{SM}}(\mu )+\tau ^{}y_i^{(4)}(\mu )\right]Q_i(\mu ).$$
(11)
with $`\tau =V_{ts}^{}V_{td}`$ and $`\tau ^{}=V_{t^{^{}}s}^{}V_{t^{^{}}d}`$. In comparison with the SM, one may introduce the new effective coefficient functions $`y_i^{\mathrm{new}}(\mu )`$
$$y_i^{\mathrm{new}}(\mu )=y_i^{\mathrm{SM}}+\frac{\tau ^{}}{\tau }y_i^{(4)}(\mu ),$$
(12)
where $`y_i(\mu )`$ are the Wilson coefficient functions in the SM and $`y_i^{(4)}`$ are the ones due to fourth generation quark contributions. The evolution for $`y_i^{(4)}(\mu )`$ is analogy to the one $`y_i^{\mathrm{SM}}(\mu )`$ in SM except replacing the t-quark by $`t^{}`$ quark. The corresponding diagrams of penguin and box are shown in fig. 1.
Using (11) and (12), eq. (7) can be written as
$`({\displaystyle \frac{ϵ^{}}{ϵ}})`$ $`=`$ $`({\displaystyle \frac{ϵ^{}}{ϵ}})^{\mathrm{SM}}+({\displaystyle \frac{ϵ^{}}{ϵ}})^{(4)},`$
$`({\displaystyle \frac{ϵ^{}}{ϵ}})^{(4)}`$ $`=`$ $`\text{Im}\lambda _t^{}\left[P^{(1/2)}P^{(3/2)}\right],`$ (13)
where the definitions of $`P^{(1/2)}`$ and $`P^{(3/2)}`$ are the same as (8) and (9) only by changing $`y_i(\mu )`$ into $`y_i^{(4)}(\mu )`$, and
$$\text{Im}\lambda _t^{}=\text{Im}V_{t^{^{}}s}^{}V_{t^{^{}}d}.$$
(14)
Thus the main test of evaluating $`ϵ^{}/ϵ`$ in the SM4 is to calculate the Wilson coefficients $`y_i^{(4)}(\mu )`$ and to provide the possible constraints on $`\text{Im}\lambda _t^{}`$. The constraints of $`\text{Im}\lambda _t^{}`$ will be discussed in next section. The calculation of $`y_i^{(4)}(\mu )`$ is the same as their counterpart $`y_i^{\mathrm{SM}}(\mu )`$ in SM and can be simply done by changing $`m_t`$ to $`m_t^{^{}}`$, which is easy to be found in any corresponding reviews. Here we repeat the same calculations and only provide the numerical results for $`y_i^{(4)}(\mu )`$ as the functions of the mass $`m_t^{^{}}`$. In the numerical calculations we take a large range for $`t^{^{}}`$-quark mass $`m_t^{}=`$ 50GeV, 100GeV, 150GeV, 200GeV, 250GeV, 300GeV, 400Gev See tab. 2,
## 3 Constraints on CKM Factor $`V_{t^{^{}}s}^{}V_{t^{^{}}d}`$ in SM4
Though we have no direct information for the additional fourth generation CKM matrix elements, while constraints may be obtained from some rare meson decays. In ref., we obtained the values of the fourth CKM factor $`V_{t^{^{}}s}^{}V_{t^{^{}}b}`$ from the decay of $`Bs\gamma `$. In this paper, we shall investigate three rare $`K`$ meson decays: two semi-leptonic decays $`K^+\pi ^+\nu \overline{\nu }`$ and $`K_L\pi ^0\nu \overline{\nu }`$, and one leptonic decay $`K_L\mu ^+\mu ^{}`$ within SM4. These decays can provide certain constraints on the fourth generation CKM factors, $`V_{t^{^{}}s}^{}V_{t^{^{}}d}`$ , $`\text{Im}V_{t^{^{}}s}^{}V_{t^{^{}}d}`$ and $`\text{Re}V_{t^{^{}}s}^{}V_{t^{^{}}d}`$ respectively.
Within SM, the decays $`K\pi \nu \overline{\nu }`$ are loop-induced semileptonic FCNC process determined only by $`Z^0`$-penguin and box diagram. These decays are the theoretically cleanest decays in rare K-decays. The great virtue of $`K_L\pi ^0\nu \overline{\nu }`$ is that it proceeds almost exclusively through direct CP violation which is very important for the investigation of $`ϵ^{^{}}/ϵ`$ in SM4. The precise calculation of these two decays at the NLO in SM can be found in Refs. While experimentally, its branching ratio has not yet been well measured, only an upper bound has be given and is larger by one order of magnitude than the one in SM (see tab. 3)<sup>1</sup><sup>1</sup>1From ref., one can easily derive by means of isospin symmetry the following model independent bound: $`Br(K^0\pi ^+\nu \overline{\nu })<4.4Br(K_L\pi ^+\nu \overline{\nu })`$ which gives $`Br(K^0\pi ^+\nu \overline{\nu })<6.1\times 10^9`$ This bond is much stronger than the direct experimental bound.. This remains allowing the New Physics to dominate their decay amplitude. Moreover, Unlike the previous two semi-leptonic decays, the branching ratio $`Br(K_L\mu ^+\mu ^{})`$ has already been measured with a very good precision. While its experimental result is several times larger than theoretical prediction in SM (see tab. 3). This also provides a window for New Physics.
In the SM4, the branching ratios of the three decay modes mentioned above receive additional contributions from the up-type quark $`t^{}`$
$$Br(K^+\pi ^+\nu \overline{\nu })=\kappa _+\left|\frac{V_{cd}V_{cs}^{}}{\lambda }P_0+\frac{V_{td}V_{ts}^{}}{\lambda ^5}\eta _tX_0(x_t)+\frac{V_{t^{}d}V_{t^{}s}^{}}{\lambda ^5}\eta _t^{}X_0(x_t^{})\right|^2,$$
(15)
$$Br(K_L\pi ^0\nu \overline{\nu })=\kappa _L\left|\frac{\mathrm{Im}V_{td}V_{ts}^{}}{\lambda ^5}\eta _tX_0(x_t)+\frac{\mathrm{Im}V_{t^{}d}V_{t^{}s}^{}}{\lambda ^5}\eta _t^{}X_0(x_t^{})\right|^2,$$
(16)
$$Br(K_L\mu \overline{\mu })_{\mathrm{SD}}=\kappa _\mu \left[\frac{\mathrm{Re}\left(V_{cd}V_{cs}^{}\right)}{\lambda }P_0^{}+\frac{\mathrm{Re}\left(V_{td}V_{ts}^{}\right)}{\lambda ^5}Y_0(x_t)+\frac{\mathrm{Re}\left(V_{t^{}d}V_{t^{}s}^{}\right)}{\lambda ^5}Y_0(x_t^{})\right]^2.$$
(17)
where $`\kappa _+,\kappa _L,\kappa _\mu `$,$`X_0(x_t)`$ , $`X_0(x_t^{})`$, $`Y_0(x_t)`$ ,$`Y_0(x_t^{})`$,$`P_0,P_0^{}`$ may be found in Refs. The QCD correction factors are taken to be $`\eta _t=`$ 0.985 and $`\eta _t^{}=`$ 1.0 .
To solve the constrains of the 4th generation CKM matrix factors $`V_{t^{^{}}s}^{}V_{t^{^{}}d}`$, $`\mathrm{Im}V_{t^{^{}}s}^{}V_{t^{^{}}d}`$ and $`\mathrm{Re}V_{t^{^{}}s}^{}V_{t^{^{}}d}`$, we must conculate the Wilson coefficients $`X_0(x_t^{})`$ and $`Y_0(x_t^{})`$. They are the founctions of the mass of the 4th generation top quark, $`m_t^{}`$. Here we give their numerical results according to several values of $`m_t^{}`$, (see table 4)
We found that the Wilson coefficients $`X_0(x_t^{})`$ and $`Y_0(x_t^{})`$ increase with the $`m_t^{}`$. To get the largest constrain of the factors in eq. (15), (16) and (17), we must use the little value of $`m_t^{}`$. Considering that the 4th generation particles must have the mass larger than $`M_Z/2`$ , we take $`m_t^{}`$ with 50 GeV to get our constrains of those three factors.
Then, from (15), (16) and (17), we arrive at the following constraints
$$|V_{t^{^{}}s}^{}V_{t^{^{}}d}|2\times 10^4,$$
(18)
$$|\text{Im}V_{t^{^{}}s}^{}V_{t^{^{}}d}|1.2\times 10^4,$$
(19)
$$|\text{Re}V_{t^{^{}}s}^{}V_{t^{^{}}d}|1.0\times 10^4.$$
(20)
For the numerical calculations, we will take $`|\text{Im}V_{t^{^{}}s}^{}V_{t^{^{}}d}|1.2\times 10^4`$.
It is easy to check that the equation (18) obeys the CKM matrix unitarity constraint, which states that any pair of rows, or any pair of columns, of the CKM matrix are orthogonal.. The relevant one to those decay channels is
$$V_{us}^{}V_{ud}+V_{cs}^{}V_{cd}+V_{ts}^{}V_{td}+V_{t^{^{}}s}^{}V_{t^{^{}}d}=0.$$
(21)
Here we have taken the average values of the SM CKM matrix elements from Ref. . Considering the fact that the data of CKM matrix is not yet very accurate, there still exists a sizable error for the sum of the first three terms. Using the value of $`V_{t^{^{}}s}^{}V_{t^{^{}}d}`$ obtained from eq. (18), the sum of the four terms in the left hand of (21) can still be close to $`0`$, because the values of $`V_{t^{^{}}s}^{}V_{t^{^{}}d}`$ are about $`10^4`$ order, ten times smaller than the sum of the first three ones in the left of (21). Thus, the values of $`V_{t^{^{}}s}^{}V_{t^{^{}}d}`$ remain satisfying the CKM matrix unitarity constraints in SM4 within the present uncertainties.
## 4 The Numerical Analysis
In the calculation of $`ϵ^{}/ϵ`$, the main source of uncertainty are the hadronic matrix elements $`Q_i_I`$. They depend generally on the renormalization scale $`\mu `$ and on the scheme used to renormlize the operators $`Q_i`$. But the calculation of $`Q_i_I`$ is much beyond the perturbative method. They only can be tread by the non-perturbative methods, like lattice methods, $`1/N`$ expansion, chiral quark models and chiral effective langrangians, which is not sufficient to obtain the high accuracy. We shall present the analysis on $`t^{^{}}`$-quark effects when considering the uncertainties of $`Q_i_I`$ due to model-dependent calculations.
It is customary to express the matrix elements $`Q_i_I`$ in terms of non-perturbative parameters $`B_i^{(1/2)}`$ and $`B_i^{(3/2)}`$ as follows:
$$Q_i_0B_i^{(1/2)}Q_i_0^{(\mathrm{vac})};Q_i_2B_i^{(3/2)}Q_i_2^{(\mathrm{vac})}.$$
(22)
The full list of $`Q_i_I`$ is given in ref.. We take the phenomenological values of $`B_i`$ (see tab.5) except for $`B_6^{(1/2)}`$ and $`B_8^{(3/2)}`$ which are taken as input parameter with values calculated by three different non-perturbative methods. Other numerical input parameters are given in tab.6.
We take the values of $`B_6^{(1/2)}`$and $`B_8^{(3/2)}`$ in three non-perturbitive approaches, lattice methods, $`1/N`$ expansion and chiral quark models (see tab.7) and the figs (see figs.2, 3, 4) in each case respectively.
The numerical results are shown in figs. 2,3,4 which correspond to the three cases of calculating hadronic matrix elements, lattice method, 1/N eapansion, and chiral quark model, respectively. We now present a study for $`ϵ^{^{}}/ϵ`$ as functions of $`\text{Im}\lambda _t^{}`$ and $`m_t^{}`$: $`ϵ^{^{}}/ϵ`$ versus $`\text{Im}\lambda _t^{}`$ with fixing $`m_t^{^{}}`$ is plotted in figs. (a); $`ϵ^{^{}}/ϵ`$ versus $`m_t^{}`$ with fixing $`\text{Im}\lambda _t^{}`$ is ploted in figs. (b); and the allowed parameter space of $`\text{Im}\lambda _t^{}`$ and $`m_t^{^{}}`$ is plotted in figs. (c). We shall analyze each case in detail as follows.
In figs (a) we plot eight lines corresponding to $`m_t^{^{}}=`$50, 100, 150, 200 250, 300, 350, 400GeV respectively. First, we notice that the slope of the line decreases as $`m_t^{^{}}`$ increases. At a value of $`m_t^{^{}}`$, about 230GeV, the slope is zero because the second part in the right hand side of eq. (13) vanishes. The reason is similar to that in SM, i.e., with increasing $`m_t^{^{}}`$ the EW penguins become increasingly important and their contributions to $`ϵ^{^{}}/ϵ`$ are with the opposite sign to those of QCD penguins so that at some values of $`m_t^{^{}}`$ there is a cancellation. The behavior comes essentially once $`m_t^{^{}}`$ becomes larger than 230GeV, the slope is negative. Its absolute value increases with $`m_t^{^{}}`$. Such a behavior comes essentially from the change of the Wilson coefficients $`y_i^{(4)}`$ as $`m_t^{^{}}`$. Second, from figs (a), we found, within the constraints on $`\text{Im}\lambda _t^{}`$ from the three rare $`K`$ meson decays, that $`ϵ^{^{}}/ϵ`$ can generally be consistent with the experimental average except for some ranges of $`m_t^{^{}}`$ once the non-perturbative parameters $`B_6^{(1/2)}`$ and $`B_8^{(3/2)}`$ are taken values calculated based on the lattice gauge theory and $`1/N`$ expansion. Such a range roughly ranges from 170GeV to 300GeV, which can be seen from figs. (2a) and (3a). There is no excluded range for the case of the chiral quark model. This is because in the first two cases, the SM values $`(ϵ^{^{}}/ϵ)^{\mathrm{SM}}`$ are about $`8.8\times 10^4`$, which is much lower than the experimental average. For a large range of $`m_t^{^{}}`$, $`(ϵ^{^{}}/ϵ)^{(4)}`$ is not large enough to make total $`ϵ^{^{}}/ϵ`$ reach the experimental average. But in the chiral quark model, the SM value is about $`18.8\times 10^4`$ which is in the $`1\sigma `$ error range of the present experimental average so that $`ϵ^{^{}}/ϵ`$ can reach the experimental average for all values of $`m_t^{^{}}`$ in the reasonable region. Thus once the non-perturbative method calculations become more reliable and the experimental measurements get more accuracy, it may provide more strong constraints on the forth generation quark from the study on $`ϵ^{^{}}/ϵ`$. Unfortunately, we can’t get any information on the upper bound of $`m_t^{^{}}`$.
We also plot in Figs (b) eight curves corresponding to $`\text{Im}\lambda _t^{}=`$1.0, 0.75, 0.5, 0.25, -0.25, -0.5, -0.75, -1.0$`\times 10^4`$ respectively. Thus similar results as those in the figs (a) are arrived. These curves are divided into two types determined by the sign of the fourth generation CKM factor $`\text{Im}\lambda _t^{}`$. The reason is also similar to the analysis for figs. (a). The figs.(b) also show the constraints on $`\text{Im}\lambda _t^{}`$. It is interesting to see that there is an excluded region from 0 to $`0.6\times 10^4`$ based on lattice gauge theory results and from 0 to $`0.76\times 10^4`$ based on the $`1/N`$ expansion results. While there is no such an excluded region based on the chiral quark model results. The reason is the same as that in the analysis of figs (a). Moreover, it seems that $`\text{Im}\lambda _t^{}`$ favors the negative values which may be interesting since the negative value of $`\text{Im}\lambda _t^{}`$ is better to satisfy the unitarity constraints of the CKM matrix (see eq. (21)). Therefore if there could exist the fourth generation, from both the theoretical and the experimental parts, one might be able get usefull information on the fourth generation CKM matrix elements, such as $`V_{t^{^{}}s}^{}V_{t^{^{}}b}`$ which has been studied in our previous paper.
In figs. (c), we show the correlation between $`\text{Im}\lambda _t^{}`$ and $`m_t^{^{}}`$. The three curves in the figure correspond to the experimental values of the new world average and its $`1\sigma `$ error, respectively. It is seen that the the allowed parameter space is strongly limited for all three cases when the ratio $`ϵ^{^{}}/ϵ`$ is around the present experimental average within 1 $`\sigma `$ error. The allowed parameter space is divided into two pieces except in the chiral quark model. This is in agreement with the analyses in figs. (a) and (b). Such a small parameter space indicates that $`ϵ^{^{}}/ϵ`$ may impose a very strong constraint on the mass and mixing of the fourth generation up-type quark.
## 5 Conclusion
In summary, we have investigated the direct CP-violating parameter $`ϵ^{}/ϵ`$ in $`K^0\overline{K}^0`$ system with considering the up-type quark $`t^{}`$ in SM4. The basic formulae for $`ϵ^{}/ϵ`$ in SM4 has been presented and the Wilson coefficient functions in the SM4 have also been evaluated. The numerical results of the additional Wilson coefficient functions have been given as functions of the mass $`m_t^{}`$. We have also studied the relevant rare $`K`$ meson decays: two semi-leptonic decays, $`K^+\pi ^+\nu \overline{\nu }`$ and $`K_L\pi ^0\nu \overline{\nu }`$, and one leptionic decay $`K_L\mu ^+\mu ^{}`$, which allow us to obtain the bounds on the fourth generation CKM matrix factor $`V_{t^{^{}}s}^{}V_{t^{^{}}d}`$. In particular, we have analyzed the numerical result of $`ϵ^{}/ϵ`$ as the function of $`m_t^{}`$ and imaginary part of the fourth CKM factor, $`\text{Im}V_{t^{^{}}s}^{}V_{t^{^{}}d}`$ (or $`V_{t^{^{}}s}^{}V_{t^{^{}}d}`$ and a fourth generation CKM matrix phase $`\theta `$). The correlation between $`ϵ^{}/ϵ`$ and $`\text{Im}V_{t^{^{}}s}^{}V_{t^{^{}}d}`$ has been studied in detail with different hadronic matrix elements calculated from various approaches, such as lattice gauge method, $`1/N`$ expansion and chiral quark model. It has been seen that, unlike the SM, when taking the central values of all parameters, the values of $`ϵ^{}/ϵ`$ can be easily made to be consistent with the current experimental data for all estimated values of the relevant hadronic matrix elements from various approaches. Especially, we have also investigated the allowed parameter space of $`m_t^{^{}}`$ and $`\text{Im}V_{t^{^{}}s}^{}V_{t^{^{}}d}`$, as a consequence, when considering $`1\sigma `$-error of the current experimental data for $`ϵ^{}/ϵ`$, the allowed parameter space for $`m_t^{^{}}`$ and $`\text{Im}V_{t^{^{}}s}^{}V_{t^{^{}}d}`$ is very small and strongly restricted. This implies that the experimental data in K system can provide strong constraints on the mass of $`t^{}`$-quark and also on the fourth generation quark mixing matrix.
## Acknowledgments
This research was supported in part by the National Science Foundation of China.
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# Bulk dynamics for interfacial growth models
## I Introduction
In the last fifteen years there has been a great interest in the study of the growth of surfaces by dynamic processes based in addition and substraction of particles (see, for example, ). For instance, the understanding of the conditions under which a growing surface shows a rough structure is nowadays of the greatest importance in the production of thin films and/or pure cristals. Surface growth is usually studied by using lattice models in which simple stochastic rules intend to mimic the relevant phenomena. Their extensive computer simulation have been of a major importance in characterizing and understanding the different morphologies that occur in real experiments. However, due to the intrinsic limitations of computers capabilities, some very interesting aspects are usually subject to inconclusive analysis and data interpretation. In particular, let us remark the inherent difficulty in the study of the surface long time behavior and its scaling properties. Nevertheless, for this particular aspect, analytical models seem to give us the answers to the questions that the computers fail to clarify. These are mainly based in postulating a Langevin equation for the height of the interface measured from a reference substrate. Such Langevin equations intend to mimic the system microscopic dynamics and its collective effect at large scales in space and time. A general choice has the structure
$`_th_t(𝐱)`$ $`=\nu _1^2h_t+\nu _2|h_t|^2+\nu _3^2^2h_t`$ (2)
$`+\mathrm{}+\eta _t(𝐱),`$
where $`h_t(𝐱)`$ is the height of the interface at time $`t`$ at the substrate position $`𝐱R^d`$ and $`\eta _t(𝐱)`$ is a white noise term. It is generally assumed a one-to-one correspondence between various terms in (2) and different physical processes (for example see the discusion in concerning a model for epitaxial growth, or a general method in to propose an equation such as (2) by using the reparametrization invariance symmetry). The details of the microscopic processes that are assumed to be irrelevant at this scale of observation are taken into account through the values of the coefficients $`\nu _1`$, $`\nu _2`$, $`\mathrm{}`$ and the properties of the noise term $`\eta _t`$. Once the Langevin equation (2) is defined, one may apply renormalization group procedures to obtain different universal properties and scaling behaviors. The success of this scheme is clearly represented by the definition and analysis of the Kardar-Parisi-Zhang equation (KPZ) which has been a clear breakthrough in the study of the space-time asymptotic behavior of growth models.
Quite often surface growth is a consequence of bulk dynamic processes. A good example of this is provided by the growth of bacterial colonies where bacteria multiply in a nutrient environment, the shape of the colony being the moving interface . In general, the dynamics of the particles before and after its aggregation to a substrate may influence the system interfacial behavior. For instance, it may appear shadowing effects as it happens in DLA processes and in thin film growth , or they may induce different scaling regimes depending on the time interval studied as it is shown in some molecular beam epitaxy models in which a system bulk dynamics is defined . However, interface models are usually expressed in terms of a height field, $`h_t(𝐱)0`$. In doing so, bulklike contributions are neglected since only interfacial degrees of freedom are being considered. Unfortunately, the mathematical hurdles to deriving the phenomenological dynamics of interfaces from stochastic bulk microscopic models is formidable, and a comprehensive theoretical picture is still lacking although a significant body of rather rigorous work has been devoted to the subject .
It is well known that the macroscopic behavior of systems at nonequilibrium states exhibits a strong dependence on the functional structure of its microscopic dynamics (for instance, in the so-called two-temperature Ising model one finds that the phase diagram changes radically depending on the analytic form for the probability of a spin-flip). Nevertheless, one may expect that this strong relation between microscopic dynamics and macroscopic behavior should dissappear near a renormalization group (RG) fixed point or in the scaling regime where universality seems to guarantee that the microscopic details are irrelevant (at least, one knows that this is true when studying dynamic properties of equilibrium systems near a (RG) fixed point ). However, some recent results on the critical behavior of a nonequilibrium driven diffusive system show that the microscopic dynamics may play a relevant role in the determination of the system universality class . The influence of the microsocopic dynamical details into the critical and non-critical properties of a nonequilibrium model implies, in our opinion, that any a priori construction of a Langevin equation as (2) may occasionally disregard important features.
In this paper we introduce a quite general class of nonequilibrium bulk growth models for which the aforementioned problems can be addressed (there are in the literature some efforts in this direction ). We shall define a stochastic bulk local dynamics expressed by an appropiated continuum Master equation in which, for simplicity, only addition of particles is considered. From the Master equation and using a truncated Kramers-Moyal expansion, we shall derive a Langevin equation for the bulk degrees of freedom in which there is an explicit dependence on the analytic form of the rates. In order to study surface properties, in a subsequent section we shall derive, from this bulk equation, an expression for the interfacial height field dynamics. The illustrative example we take is that of the KPZ equation and, for consistency, in that particular case we check our results by means of a Monte Carlo simulation. Some other examples are then briefly commented and our conclusions are given in the final section.
## II Growth models: A general description
Let us consider a particle density field, $`\rho (𝐱,\tau )`$, $`𝐱R^{d+1}`$, and assume that the probability distribution $`P_\tau (\rho )`$ associated with each field configuration obeys the following continuous Master equation
$`_\tau P_\tau (\rho )=`$ $`{\displaystyle _{R^{d+1}}}d𝐫[c(\rho _𝐫\rho )P_\tau (\rho _𝐫)`$ (4)
$`c(\rho \rho ^𝐫)P_\tau (\rho )].`$
Here, $`c(\rho \rho ^{})`$ is the probability per unit time (or transition rate) from one configuration $`\rho `$ to another $`\rho ^{}`$, and $`\rho _𝐫=\rho (𝐱)\mathrm{\Omega }^1\delta (𝐱𝐫)`$, $`\rho ^𝐫=\rho (𝐱)+\mathrm{\Omega }^1\delta (𝐱𝐫)`$. Note that the density field, $`\rho `$, can only grow in steps of size $`\mathrm{\Omega }^1`$. This is consistent with a picture in which $`\rho (𝐫)`$ is a particle density that results after coarse graining over blocks of size $`\mathrm{\Omega }`$ centered around $`𝐫`$ in a lattice. Therefore, the Master equation (4), so defined, could be thought of as if it only described processes that add one particle per block of the lattice per elementary time step. This is schematically represented in figure 1.
Next, we choose the transition rates such that
$$c(\rho \rho ^{})=w(\rho ;𝐫),$$
(5)
namely, they are a function that depends only on the initial configuration $`\rho `$ and the specific point where mass is added. Now, let us assume that $`\mathrm{\Omega }`$ is large enough (tipically, it should be much bigger than any microcopic length scale present in the original physical problem but much smaller than the correlation lenght of the system) and expand the Master equation (4) on invers powers of $`\mathrm{\Omega }`$. Then, using the expansions
$$c(\rho \rho ^𝐫)=w(\rho ;𝐫),$$
(6)
$`c(\rho _𝐫\rho )`$ $`=`$ $`w(\rho \mathrm{\Omega }^1\delta (xr);𝐫)`$ (7)
$`=`$ $`{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(1)^m}{m!\mathrm{\Omega }^m}}{\displaystyle \frac{\delta ^m}{\delta \rho (𝐫)^m}}w(\rho ;𝐫),`$ (8)
$$P_\tau (\rho _𝐫)=\underset{m=0}{\overset{\mathrm{}}{}}\frac{(1)^m}{m!\mathrm{\Omega }^m}\frac{\delta ^m}{\delta \rho (𝐫)^m}P_\tau (\rho ),$$
(9)
we get the so-called Kramers-Moyal expansion of the Master equation (4) ,
$$_\tau P_\tau (\rho )=_{R^{d+1}}𝑑𝐫\underset{l=1}{\overset{\mathrm{}}{}}\frac{(1)^l}{l!\mathrm{\Omega }^l}\frac{\delta ^l}{\delta \rho (𝐫)^l}\left[w(\rho ;𝐫)P_\tau (\rho )\right].$$
(10)
The next step we take is to keep only the first two terms in (10). Then, we can write down the following Fokker-Planck equation
$`_\tau P_\tau (\rho )=`$ $`{\displaystyle _{R^{d+1}}}d𝐫[{\displaystyle \frac{1}{\mathrm{\Omega }}}{\displaystyle \frac{\delta }{\delta \rho (𝐫)}}`$ (12)
$`+{\displaystyle \frac{1}{2\mathrm{\Omega }^2}}{\displaystyle \frac{\delta ^2}{\delta \rho (𝐫)^2}}\left]\right(w(\rho ;𝐫)P_t(\rho )).`$
To control the goodness of such approximation we appeal to Kurtz theorem , by virtue of which when $`\mathrm{\Omega }\mathrm{}`$, and for a given time $`T<\mathrm{}`$ then
$$\underset{\tau <T}{sup}|\rho _\tau \stackrel{~}{\rho }_\tau |\zeta _\mathrm{\Omega }^T\frac{\mathrm{log}\mathrm{\Omega }}{\mathrm{\Omega }}.$$
(13)
where $`\rho _\tau `$ and $`\stackrel{~}{\rho }_\tau `$ are typical time trajectories on phase space which are solutions of the exact Master equation (4) and the Fokker-Planck one (12) respectively. $`\zeta _\mathrm{\Omega }^T`$ is a random variable whose distribution does not depend on $`\mathrm{\Omega }`$ and satisfying $`\mathrm{exp}(\lambda \zeta _\mathrm{\Omega }^T)<\mathrm{}`$ for any constant $`\lambda >0`$. That is, for a given fixed time $`T`$ one can always find a large enough $`\mathrm{\Omega }`$ such that the difference between solving exactly the master equation or solving the truncated version of it, is of order $`\mathrm{log}\mathrm{\Omega }/\mathrm{\Omega }`$ and therefore negligible. Moreover, this bound is the best one and no new terms of the Kramers-Moyal expansion give better results. However, when $`T\mathrm{}`$ one cannot control, in general, the accumulated influence of the neglected terms in the expansion. But, since the study of growth models is manily focused on the undestanding of their evolution properties, due to Kurtz’s theorem the Fokker-Planck equation (12) is a valid theoretical starting point.
Lastly, the Fokker-Planck equation (12) is equivalent in the Stratonovich sense to the Langevin equation
$`_t\rho (𝐫,t)=`$ $`w(\rho ;𝐫){\displaystyle \frac{1}{4\mathrm{\Omega }}}{\displaystyle \frac{\delta }{\delta \rho (𝐫)}}w(\rho ;𝐫)`$ (14)
$`+`$ $`{\displaystyle \frac{1}{\mathrm{\Omega }^{1/2}}}w(\rho ;𝐫)^{1/2}\nu _t(𝐫)`$ (15)
where $`t=\mathrm{\Omega }^1\tau `$, $`\nu _t(𝐫)`$ is a white noise with zero mean and $`\nu _t(𝐫)\nu _t^{}(𝐫^{})=\delta (𝐫𝐫^{})\delta (tt^{})`$. We recall that $`w(\rho ;𝐫)`$ is the probability per unit time of adding a particle (of mass $`\mathrm{\Omega }^1`$) at point $`𝐫`$. If we considered the posibility of particle substractions then we should do the following substitutions in equation (15): $`ww_+w_{}`$ in the first term and and $`ww_++w_{}`$ in the last two terms. $`w_{+()}w(\rho ;𝐫)_{+()}`$ is the probability per unit time of adding (substracting) a particle.
Before we proceed to the next section, let us remark that a) eq. (15) describes the evolution of the bulk density field of a system that growths by addition of particles with, in principle, no other physical assumptions, and b) the Langevin equation (15) depends directly on the functional form of the bulk rates. The election of these bulk rates then provides the physical restrictions for the particular growth model that is going to be specifically modelled. Also, it is remarkable that the influence of the bulk dynamics affects the noise term through a nontrivial factor.
## III Height dynamics: KPZ equation and Monte Carlo simulation.
We proceed to single out the interfacial degrees of freedom of equation (15). In order to achieve this, we place two conditions on the solutions of equation (15). First, let us impose that our bulk dynamics produces a surface perpendicular to the $`z`$-axis without overhangs and bubbles. This condition is necessary to ensure that we have a well defined interface (note that equation (15) contains overhang/vacancy and shadowing effects). Secondly, we shall neglect any interfacial profile. Then, we may assume that the solutions of the Langevin equation (15) have the form
$$\rho (𝐫,t)=\mathrm{\Theta }\left(h_t(𝐱)z\right)$$
(16)
where $`𝐫(𝐱,z)`$, $`𝐱R^d`$ is a point in the substrate and $`h_t(𝐱)`$ is the height of the growing surface at time $`t`$. $`\mathrm{\Theta }(\lambda )=1,1/2,0`$ if $`\lambda >0,=0,<0`$ respectively. That is, for a given point in the bulk $`𝐫`$, if its $`z`$ coordinate is larger, equal or smaller than the actual position of the surface, $`h_t(𝐱)`$, then the density field is $`\rho (𝐫,t)=0`$, $`1/2`$ and $`1`$ respectively. Note that since the $`\mathrm{\Theta }`$ function is not continuous, when differenciating we should use a regularized version of it, e.g. $`\mathrm{\Theta }(x)=\frac{1}{2}[1+\mathrm{tanh}(ax)]`$ with $`a\mathrm{}`$.
We are interested in constructing a dynamical equation for the $`h_t(𝐱)`$ fields. Therefore, let us make a time derivative in equation (16)
$$_t\rho (𝐫,t)=_th_t(𝐱)\delta (h_t(𝐱)z).$$
(17)
Integrating in $`z`$ both sides of equation (17) we find
$$_th_t(𝐱)=_R𝑑z_t\rho (𝐫,t).$$
(18)
Equating this expression for $`_th_t(𝐱)`$, together with (15), will lead us to the desired Langevin equation for the heights.
To make this a bit more concrete, we now introduce a particular set of rates. For instance, we choose the probability of adding mass to the point $`𝐫`$ to be proportional to the square of the gradient of the density field in that point: $`w(\rho ;𝐫)|\rho |^2`$. With this election the unwanted effect of nucleation of bubbles is avoided. After a bit of algebra, we get
$`_th_t=`$ $`\alpha \left(1+(h_t)^2\right)+{\displaystyle \frac{D}{2\mathrm{\Omega }}}\mathrm{\Delta }h`$ (20)
$`+{\displaystyle \frac{1}{\mathrm{\Omega }^{1/2}}}\left(1+(h_t)^2\right)^{1/2}\nu _t,`$
which is the celebrated KPZ equation with a different noise term (a naive power counting argument reduces the relevant part of (20) to the KPZ equation). The coefficient $`D`$ has the proper dimensions and $`\alpha `$ is positive and depends on how the $`\mathrm{\Theta }`$ function is regularized. This comes from our particular ansatz, but, we would like to stress here that, as far as universal properties are concerned, the precise value of the coefficients is immaterial. In fact, it is easy to show (with naive power counting) that for any bulk dynamics given by $`w(\rho ;r)|\rho |^\eta `$ with $`\eta 0`$, gives rise to a height equation falling in the KPZ universality class.
Next, we proceed to check numerically the connection between eq.(4) with $`w(\rho ;𝐫)|\rho |^2`$ and the KPZ equation (20).
Numerical results. The simple bulk rate $`|\rho |^2`$ can be easily implemented in a Monte Carlo experiment. On a two-dimensional square lattice periodic boundary conditions are considered in one of the principal axes. Each lattice site is labeled by an occupation variable $`\rho _𝐫`$ ranging from $`0,1/\mathrm{\Omega },\mathrm{}`$ to 1. A site is empty if $`\rho _𝐫=0`$ and full if $`\rho _𝐫=1`$. Initially the lattice is empty except for a full horizontal bottom line. The growth starts when an empty site $`𝐫`$ is chosen at random from the lattice. Then, $`\rho _𝐫`$ is increased in $`\mathrm{\Omega }^1`$ with a probability that depends on each nearest neighbour of $`𝐫`$ through $`|\rho _𝐫|^2`$. Since symmetrized discrete forms of the gradient operator may give rise to bulk vacancies incompatible with (16), we use the following finite-difference formula for the density derivatives $`(𝐫=(x,z))`$
$$\rho _𝐫=(\rho (x\pm 1,z)\rho (x,z),\rho (x,z)\rho (x,z1)).$$
(21)
That is, left and right derivatives are used alternately to avoid asymmetric effects and, for convenience, a unit lattice spacing is assumed. The height is defined as the distance to the highest occupied lattice site directly above the substrate coordinate $`x`$.
Figure 2 shows the scaling plot for the surface width $`W(L,t)^2=L^1[h_t(x)\overline{h}(t)]^2`$ for different system substrate sizes $`L`$. $`\overline{h}(t)`$ is the mean height of the interface at time $`t`$. The numerical data were averaged over 2000 independent runs for $`L=100,200,500`$ and over 1000 independent runs for $`L=800`$. For reasons of computational efficiency the results shown correnspond to $`\mathrm{\Omega }=1`$. Other values of $`\mathrm{\Omega }`$ yield similar results but the simulations are much more time-consuming. Good data colapse is obtained for a roughness exponent $`\alpha =1/2`$ and a dynamic exponent $`z=3/2`$, in agreement with the KPZ prediction .
The example we have just provided is by no means unique. Our formalism emcompasses many others well known growth equations. Let us just mention that with the simple dynamics given by $`w(\rho ,𝐫)=|\rho |`$ the equation of Golubovic and Wang, related to the anharmonic equilibrium thermal fluctuations of smectics A , is obtained. This is given by
$`_th_t=\left(1+(h_t)^2\right)(\lambda +{\displaystyle \frac{1}{\mathrm{\Omega }}}H)`$ (22)
$`+{\displaystyle \frac{1}{\mathrm{\Omega }^{1/2}}}\left(1+(h_t)^2\right)^{1/4}\nu _t,`$ (23)
where $`\lambda `$ is a coefficient and $`H`$ is the curvature (see ). As we have mentioned before, this kind of dynamics (proportional to $`|\rho |^\eta `$ with $`\eta 0`$), falls in the KPZ universality class. Also, the Edwards-Wilkinson equation , which favours growth at local minima, can also be recovered by considering substraction of particles and a rate of the form $`w=|^2\rho |`$. In this case, the formalism has to be slightly modified by linking the election of additon or substraction of particles to the density field configuration. More explicitely, now the Master equation defining the process reads
$`_\tau P_\tau (\rho )=`$ $`{\displaystyle _{R^{d+1}}}d𝐫[c(\rho ^{}\rho )P_\tau (\rho ^{})`$ (25)
$`c(\rho \rho ^{})P_\tau (\rho )],`$
with $`\rho ^{}=\rho (𝐱)+\alpha \mathrm{\Omega }^1\delta (𝐱𝐫)`$ and $`\alpha =1,0,1`$ for $`^2\rho `$ less, equal and greater than 0, respectively. That is, material is added to those areas where the laplacian of the density field is negative and taken from those where it is positive. In this manner, a balanced distribution of mass is achieved that results in the equilibrium Edwards-Wilkinson universality class.
Many other different rates lead to their corresponding growth equations, sometimes to the same one, showing that growth models with similar surface behavior may not have the same bulk properties.
## IV Conclusions
In this paper we have introduced a class of nonequilibrium models in which a stochastic bulk dynamics is defined. The bulk evolves by an adsorption process represented by a continuum Master equation. From it, we have derived a Langevin equation for the bulk density field whose structure depends on the details of the underlying bulk dynamics. This dependence was then extended to an equation of motion for the interfacial degrees of freedom. In particular, we have examplified the procedure by deriving the KPZ equation from a very simple bulk rate. A Monte Carlo simulation of the bulk process confirms the predicted scaling behavior for the interface. Finally, a number of examples were briefly mentioned. In all cases the bulk dynamics determines the mesoscopic height equation, showing that both scales are related in a non trivial form and that their mutual influence could be far from intuitive.
The strategy developed in this paper is quite general. It includes both local and nonlocal, and equilibrium and nonequilibrium growth processes. Therefore, a great number of growth physical phenomena can be studied, in principle, with our approach. For instance, molecular beam epitaxy models with adatom mobility, driven lattice gases or wetting phenomena by means of lattice gas theories of multilayer adsorption, to name just a few.
###### Acknowledgements.
We acknowledge E. Hernández-García and M.A. Muñoz for a critical reading of the manuscript. C.L. acknowledges the hospitality received at the Centro de Física da Matéria Condensada da Universidade de Lisboa and is supported by CICYT (MAR98-0840). F.S. is supported by the European Comission under grant ERBFMRXCT980183.
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# Anholonomic Soliton–Dilaton and Black Hole Solutions in General Relativity
## I Introduction
### A Soliton–black hole solutions in modern gravity and string/brane theories
One of the main ingredients of the last developments in string theory and gravity is the investigation of connections between black holes and non–perturbative structures of string theory such as Bogomol’nyi–Pasad–Sommerfeld (BPS) solitons or D–branes . New ways to address and propose solutions of old and new fundamental problems of black hole physics are opened via explanation of black hole thermodynamics in terms of microscopic string and membrane physics.
Similar fundamental problems of black hole physics have been recently analyzed in the recent literature using low–dimensional gravity models (see, for instance, Ref. ). Dilaton gravity theories in two spacetime (2D) dimensions have been used as theoretical laboratories of studding 4D black hole physics.
One of 2D theories of particular interest is the so–called Jackiw–Teitelboim theory (JT) because of its connection to the Liouville–Polyakov action in string theory. The black hole solutions to JT–gravity are dimensional reductions of the Banados–Teitelboim–Zanelli (BTZ) black hole solutions and exhibit usual thermodynamic properties with black hole entropy . Such JT black hole solutions describe spacetimes of constant curvature.
The aim of this paper is to develop and apply a new method of construction of 4D solutions of the Einstein equations by considering nonlinear superpositions of 2D and 3D soliton–dilaton–black hole metrics which are treated as some non–perturbative structures in general relativity.
### B Locally anisotropic soliton, black hole and disk/cylinder like solutions
Mathematicians (see, for instance, the review ) have considered for a long time the relationship between Euclidean, $`ϵ=1`$ (pseudo–Euclidean, or Lorentzian, $`ϵ=1),`$ 2D metrics
$$ds^2=ϵ\mathrm{sin}^2\left(\frac{\varphi }{2}\right)dt^2+\mathrm{cos}^2\left(\frac{\varphi }{2}\right)dr^2$$
(1)
with the angle $`\varphi (t,r)`$ solving the Lorentzian (Euclidean) sine–Gordon equation,
$$ϵ\frac{^2\varphi }{t^2}+\frac{^2\varphi }{r^2}=\stackrel{~}{m}^2\mathrm{sin}\varphi ,$$
(2)
which determine some 2D Rimannian geometries with constant negative curvature,
$$\stackrel{~}{R}=2\stackrel{~}{m}^2=const.$$
(3)
The angle $`\varphi (t,r)`$ from (1) describes an embedding of a 2D manifold into a three–dimensional (pseudo) Euclidean space.
The topic of construction of soliton solutions in gravity theories has a long history (see Refs ). For 4D vacuum Einstein gravity the problem was tackled by investigating metrics $`g_{\alpha \beta }`$ of signature $`(,+,+,+)`$ (in this paper permutations of sines will be also considered), with a 2+2 spacetime splitting,
$$ds^2=g_{ij}\left(x^i\right)dx^idx^j+h_{ab}(x^i)dy^ady^b,$$
(4)
where $`g_{ij}=diag[f(x^i),f(x^i)]`$ and $`deth_{ab}<1`$ and the local coordinates are denoted
$$u^\alpha =(\{x^i=(x^1,x^2)\},\{y^a=(y^3=z,y^4)\}),$$
or, in brief, $`u=(x,y).`$ We adopt the convention that the $`x`$–coordinates are provided with Latin indices of type $`i,j,k,\mathrm{}=1,2;`$ the $`y`$–coordinates are with indices of type $`a,b,c,\mathrm{},=3,4`$ and the 4D $`u`$–coordinates will be provided with Greek indices $`\alpha ,\beta ,\mathrm{}=1,2,3,4.`$ The meaning of coordinates (space or time like ones) will depend on the type of construction under consideration. Belinski and Zakharov identified $`y^4`$ with the time like coordinate; Maison treated $`y^a`$ as space variables. It was proved that the vacuum Einstein equations, $`R_{\alpha \beta }=0,`$ are satisfied if the components $`h_{ab}(x^i)`$ are solutions of a generalized (Euclidean) sine–Gordon equation.
In two recent papers Gegenberg and Kunstatter investigated the relationship between black holes of JT dilation gravity and solutions of the sine–Gordon field theory. Their constructions were generalized by Cadoni to soliton solutions of 2D dilaton gravity models which describes spacetimes being of non constant curvature.
In this work we explore the possibility of anholonomic generalizations of the Belinski–Zakharov–Maison soliton constructions (4) in a fashion when the coefficients of the matrix $`h_{ab}`$ could depend on three variables $`(x^i,z).`$ The matrix $`g_{ij}(x^i)`$ will be defined by some Gegenberg–Kunstatter–Cadoni 2D soliton–black hole solutions , their conformal transforms, or by the factor $`f(x^i)`$ from (4) which could be related with solitons for $`h_{ab}.`$ We emphasize that the resulting 4D (pseudo) Riemannian metrics, with generic local anisotropy (in brief, la–metrics) will be found to solve the Einstein equations with energy–momentum tensor.
For definiteness, we consider 4D metrics parametrized by ansatzs of type
$$g_{\alpha \beta }=\left[\begin{array}{cc}g_{ij}+N_i^aN_j^bh_{ab}& N_j^eh_{ae}\\ N_i^eh_{be}& h_{ab}\end{array}\right]$$
(5)
which are given with respect to a local coordinate basis $`du^\alpha =(dx^i,dy^a)`$ being dual to $`/u^\alpha =(/x^i,/y^a).`$ For simplicity, the 2D components $`g_{ij}`$ and $`h_{ab}`$ are considered to be some diagonal matrices (for two dimensions a diagonalization is always possible),
$$g_{ij}(x^k)=\left(\begin{array}{cc}g_1(x^k)& 0\\ 0& g_2(x^k)\end{array}\right)$$
(6)
and
$$h_{ab}(x^k,z)=\left(\begin{array}{cc}h_3(x^k,z)& 0\\ 0& h_4(x^k,z)\end{array}\right).$$
(7)
The components $`N_i^a=N_i^a(x^i,z)`$ will be selected as to satisfy the 4D Einstein gravitational field equations.
The metric (5) can be rewritten in a very simple form
$$g_{\alpha \beta }=\left(\begin{array}{cc}g_{ij}(x^k)& 0\\ 0& h_{ab}(x^k,z)\end{array}\right)$$
(8)
with respect to some 2+2 anholonomic bases (tetrads, or vierbiends) defined
$`\delta _\alpha `$ $`=`$ $`(\delta _i,_a)={\displaystyle \frac{\delta }{u^\alpha }}`$ (9)
$`=`$ $`\left(\delta _i={\displaystyle \frac{\delta }{x^i}}={\displaystyle \frac{}{x^i}}N_i^b(x^j,y){\displaystyle \frac{}{y^b}},_a={\displaystyle \frac{}{y^a}}\right)`$ (10)
and
$`\delta ^\beta `$ $`=`$ $`(d^i,\delta ^a)=\delta u^\beta `$ (11)
$`=`$ $`\left(d^i=dx^i,\delta ^a=\delta y^a=dy^a+N_k^a(x^j,y^b)dx^k\right).`$ (12)
The coefficients $`N_j^a\left(u^\alpha \right)`$ from (9) and (11) could be treated as the components of an associated nonlinear connection, N–connection, structure (see ; in this work we do not consider in details the N–connection geometry).
A specific point of this paper, comparing with another soliton approaches in gravity theories, is to show how anholonomic constructions can be used for generation of 4D soliton–dilaton–black hole non–perturbative structures in general relativity. This way a correspondence between solutions of so–called locally anisotropic (super) gravity and string theories and metrics given with respect to anholonomic frames in general relativity is derived.
Ansatzs of type (6), (7) and (8) can be used for construction of an another class of solutions with generic local anisotropy of the Einstein equations. If the gravitational field equations are written with respect to an anholonomic basis (9) and/or (11), the coefficients $`g_{ij},h_{ab}`$ and $`N_j^a`$ satisfy some very simplified systems of partial differential equations. We can construct various classes of black hole and disk/cylinder like solutions which in the locally isotropic limit are conformally equivalent to some well known BTZ, Schwarzschild and/or Kerr, Weyl cylindrical and Neugebauer–Meinel disk solutions. In general relativity there are admitted singular (in a point, on unclosed infinite lines or on closed curves such as ellipses, circles) solutions with horizons being described by hypersurface equations for rotation ellipsoid, torus, ellipses and so on. A physical treatment of such nonlinear configurations is to consider values like anisotropic mass, oscillation of horizons, variable interaction constants and gravitational non–linear self polarizations.
### C Outline
The paper is organized as follow:
Section II reviews the geometry of anholonomic frames on (pseudo) Riemannian spaces and associated nonlinear connection structures. There are defined the basic geometric objects and written the Einstein equations with respect to anholonomic frames split by nonlinear connections.
In Section III, there are considered the general properties and reductions of basic geometric objects and field equations for 4D metrics constructed as nonlinear superpositions of 2D horizontal (with respect to a nonlinear connection structure) metrics, depending on two horizontal coordinates, and of 2D vertical coordinates depending on three (two horizontal plus one vertical) coordinates. It is given a classification of such 4D metrics depending on signatures of 2D metrics and resulting 4D metrics.
In Section IV, we prove that the Einstein equations admit soliton like 2D (both for horizontal and vertical components of 4D metrics) and 3D (for vertical components of 4D metrics) solutions. There are examined some classes of integral varieties for the Einstein equations admitting non–perturbative structures generated, for instance, by 2D and 3D sine–Gordon and Kadomtsev–Petviashvili equations. Some exact solutions for locally anisotropic deformations of the sine–Gordon systems are constructed.
Section V describes an effective locally anisotropic soliton–dilaton field theory and contains a topological analysis of such models.
Section VI elucidates the interconnection of locally anisotropic 2D soliton and black hole solutions. Nonlinear superpositions to 4D are considered.
In Section VII, we construct 3D black hole solutions with generic local anisotropic. As some simplest examples there are taken configurations when the horizon is parametrized by an ellipse and the possibility of oscillation in time of such horizons is shown.
Section VIII is devoted to the physics of 4D black hole and disk/cylinder solutions with generic local anisotropy. There are analyzed the general properties of metrics describing such solutions and discussed the question of their physical treatment. The construction of singular solutions with various type of horizon hypersurfaces is performed by considering correspondingly the rotation ellipsoid, epllipsoidal cylinder, torus, bipolar and another systems of coordinates. It is shown that in the locally isotropic limit such solutions could be equivalent to some conformal transforms of some static or rotating configurations like for the Schwarzschild–Kerr, Weyl cylindrical and Neugebauer–Meinel disk solutions.
In Section IX, some additional examples of locally anisotropic soliton–dilaton–black hole solutions are given. It is illustrated how in general relativity we can construct two soliton non–perturbative structures, proved that nonlinear connections and non–diagonal energy–momentum tensor components can induce Kadomtsev–Petviashvily soliton like solutions and there are considered new types of two and three coordinate soliton–dilaton vacuum gravitational configurations.
Finally, in Section X, we discuss our results and present conclusions.
## II Anholonomic Frames on (Pseudo) Riemannian Spaces
We outline the geometric background on anholonomic frames modelling 2D local anisotropies (la) in 4D curved spaces (see Refs. and for details on spacetime differential geometry and N–connection structures in vector bundle spaces). We note that a frame anholonomy induces a corresponding local anisotropy. Spacetimes enabled with anholonomic frame (and associated N–connection) structures are also called locally anisotropic spacetimes, in brief la–spacetimes.
In this paper spacetimes are modelled as smooth (i.e class $`C^{\mathrm{}})`$ 4D (pseudo) Riemannian manifolds $`V^{(3+1)}`$ being Hausdorff, paracompact and connected and provided with corresponding geometric structures of symmetric metric $`g_{\alpha \beta }`$ of signature $`(,+,+,+)`$ and of linear, in general nonsymmetric, connection $`\mathrm{\Gamma }_{\beta \gamma }^\alpha `$ defining the covariant derivation $`D_\alpha `$ satisfying the metricity conditions $`D_\alpha g_{\beta \gamma }=0.`$ The indices are given with respect to a tetradic (frame) vector field $`\delta ^\alpha =(\delta ^i,\delta ^a)`$ and its dual $`\delta _\alpha =(\delta _i,\delta _a).`$
A frame (local basis) structure $`\delta _\alpha `$ (11) on $`V^{(3+1)}`$ is characterized by its anholonomy coefficients $`w_{\beta \gamma }^\alpha `$ defined from relations
$$\delta _\alpha \delta _\beta \delta _\beta \delta _\alpha =w_{\alpha \beta }^\gamma \delta _\gamma .$$
(13)
The elongation (by N–coefficients) of partial derivatives in the locally adapted partial derivatives (9) reflects the fact that on the (pseudo) Riemannian spacetime $`V^{(3+1)}`$ it is modelled a generic local anisotropy characterized by anholonomy relations (13) when the anholonomy coefficients are computed as follows
$`w_{ij}^k`$ $`=`$ $`0,w_{aj}^k=0,w_{ia}^k=0,w_{ab}^k=0,w_{ab}^c=0,`$ (14)
$`w_{ij}^a`$ $`=`$ $`\mathrm{\Omega }_{ij}^a,w_{aj}^b=_aN_i^b,w_{ia}^b=_aN_i^b,`$ (15)
where
$$\mathrm{\Omega }_{ij}^a=_iN_j^a_jN_i^a+N_i^b_bN_j^aN_j^b_bN_i^a$$
defines the coefficients of N–connection curvature, in brief, N–curvature. On (pseudo) Riemannian spaces this is a characteristic of a chosen anholonomic system of reference.
A 2+2 anholonomic structure distinguishes (d) the geometrical objects into horizontal (h) and vertical (v) components. Such objects are briefly called d–tensors, d–metrics and/or d–connections. Their components are defined with respect to a la–basis of type (9), its dual (11), or their tensor products (d–linear or d–affine transforms of such frames could also be considered). For instance, a covariant and contravariant d–tensor $`Z,`$ is expressed
$`Z`$ $`=`$ $`Z_\beta ^\alpha \delta _\alpha \delta ^\beta `$ (16)
$`=`$ $`Z_j^i\delta _id^j+Z_a^i\delta _i\delta ^a+Z_j^b_bd^j+Z_a^b_b\delta ^a.`$ (17)
A linear d–connection $`D`$ on la–space $`V^{(3+1)},`$
$$D_{\delta _\gamma }\delta _\beta =\mathrm{\Gamma }_{\beta \gamma }^\alpha (x,y)\delta _\alpha ,$$
is parametrized by non–trivial h–v–components,
$$\mathrm{\Gamma }_{\beta \gamma }^\alpha =(L_{jk}^i,L_{bk}^a,C_{jc}^i,C_{bc}^a).$$
(18)
A metric on $`V^{(3+1)}`$ with 2+2 block coefficients (8) is written in distinguished form, as a metric d–tensor (in brief, d–metric), with respect a la–base (11)
$`\delta s^2`$ $`=`$ $`g_{\alpha \beta }\left(u\right)\delta ^\alpha \delta ^\beta `$ (19)
$`=`$ $`g_{ij}(x,y)dx^idx^j+h_{ab}(x,y)\delta y^a\delta y^b.`$ (20)
Some d–connection and d–metric structures are compatible if there are satisfied the conditions
$$D_\alpha g_{\beta \gamma }=0.$$
For instance, a canonical compatible d–connection
$${}_{}{}^{c}\mathrm{\Gamma }_{\beta \gamma }^{\alpha }=\left({}_{}{}^{c}L_{jk}^{i},^cL_{bk}^a,^cC_{jc}^i,^cC_{bc}^a\right)$$
is defined by the coefficients of d–metric (19), $`g_{ij}(x,y)`$ and $`h_{ab}(x,y),`$ and by the N–coefficients,
$`{}_{}{}^{c}L_{jk}^{i}`$ $`=`$ $`{\displaystyle \frac{1}{2}}g^{in}\left(\delta _kg_{nj}+\delta _jg_{nk}\delta _ng_{jk}\right),`$ (21)
$`{}_{}{}^{c}L_{bk}^{a}`$ $`=`$ $`_bN_k^a+{\displaystyle \frac{1}{2}}h^{ac}\left(\delta _kh_{bc}h_{dc}_bN_i^dh_{db}_cN_i^d\right),`$ (22)
$`{}_{}{}^{c}C_{jc}^{i}`$ $`=`$ $`{\displaystyle \frac{1}{2}}g^{ik}_cg_{jk},`$ (23)
$`{}_{}{}^{c}C_{bc}^{a}`$ $`=`$ $`{\displaystyle \frac{1}{2}}h^{ad}\left(_ch_{db}+_bh_{dc}_dh_{bc}\right)`$ (24)
The coefficients of the canonical d–connection generalize for la–spacetimes the well known Cristoffel symbols.
For a d–connection (18) the components of torsion,
$`T(\delta _\gamma ,\delta _\beta )`$ $`=T_{\beta \gamma }^\alpha \delta _\alpha ,`$ (25)
$`T_{\beta \gamma }^\alpha `$ $`=\mathrm{\Gamma }_{\beta \gamma }^\alpha \mathrm{\Gamma }_{\gamma \beta }^\alpha +w_{\beta \gamma }^\alpha `$ (26)
are expressed via d–torsions
$`T_{.jk}^i`$ $`=`$ $`T_{jk}^i=L_{jk}^iL_{kj}^i,T_{ja}^i=C_{.ja}^i,T_{aj}^i=C_{ja}^i,`$ (27)
$`T_{.ja}^i`$ $`=`$ $`0,T_{.bc}^a=S_{.bc}^a=C_{bc}^aC_{cb}^a,`$ (28)
$`T_{.ij}^a`$ $`=`$ $`\mathrm{\Omega }_{ij}^a,T_{.bi}^a=_bN_i^aL_{.bj}^a,T_{.ib}^a=T_{.bi}^a.`$ (29)
We note that for symmetric linear connections the d–torsion is induced as a pure anholonomic effect.
In a similar manner, putting non–vanishing coefficients (18) into the formula for curvature
$`R(\delta _\tau ,\delta _\gamma )\delta _\beta `$ $`=R_{\beta \gamma \tau }^\alpha \delta _\alpha ,`$ (30)
$`R_{\beta \gamma \tau }^\alpha `$ $`=\delta _\tau \mathrm{\Gamma }_{\beta \gamma }^\alpha \delta _\gamma \mathrm{\Gamma }_{\beta \delta }^\alpha +`$ (32)
$`\mathrm{\Gamma }_{\beta \gamma }^\phi \mathrm{\Gamma }_{\phi \tau }^\alpha \mathrm{\Gamma }_{\beta \tau }^\phi \mathrm{\Gamma }_{\phi \gamma }^\alpha +\mathrm{\Gamma }_{\beta \phi }^\alpha w_{\gamma \tau }^\phi ,`$
we can compute the components of d–curvatures
$`R_{h.jk}^{.i}`$ $`=`$ $`\delta _kL_{.hj}^i\delta _jL_{.hk}^i`$ (34)
$`+L_{.hj}^mL_{mk}^iL_{.hk}^mL_{mj}^iC_{.ha}^i\mathrm{\Omega }_{.jk}^a,`$
$`R_{b.jk}^{.a}`$ $`=`$ $`\delta _kL_{.bj}^a\delta _jL_{.bk}^a`$ (36)
$`+L_{.bj}^cL_{.ck}^aL_{.bk}^cL_{.cj}^aC_{.bc}^a\mathrm{\Omega }_{.jk}^c,`$
$`P_{j.ka}^{.i}`$ $`=`$ $`_kL_{.jk}^i+C_{.jb}^iT_{.ka}^b`$ (38)
$`(_kC_{.ja}^i+L_{.lk}^iC_{.ja}^lL_{.jk}^lC_{.la}^iL_{.ak}^cC_{.jc}^i),`$
$`P_{b.ka}^{.c}`$ $`=`$ $`_aL_{.bk}^c+C_{.bd}^cT_{.ka}^d`$ (40)
$`(_kC_{.ba}^c+L_{.dk}^cC_{.ba}^dL_{.bk}^dC_{.da}^cL_{.ak}^dC_{.bd}^c)`$
$`S_{j.bc}^{.i}`$ $`=`$ $`_cC_{.jb}^i_bC_{.jc}^i+C_{.jb}^hC_{.hc}^iC_{.jc}^hC_{hb}^i,`$ (41)
$`S_{b.cd}^{.a}`$ $`=`$ $`_dC_{.bc}^a_cC_{.bd}^a+C_{.bc}^eC_{.ed}^aC_{.bd}^eC_{.ec}^a.`$ (42)
The Ricci tensor
$$R_{\beta \gamma }=R_{\beta \gamma \alpha }^\alpha $$
has the d–components
$`R_{ij}`$ $`=`$ $`R_{i.jk}^{.k},R_{ia}=^2P_{ia}=P_{i.ka}^{.k},`$ (43)
$`R_{ai}`$ $`=`$ $`{}_{}{}^{1}P_{ai}^{}=P_{a.ib}^{.b},R_{ab}=S_{a.bc}^{.c}.`$ (44)
We point out that because, in general, $`{}_{}{}^{1}P_{ai}^{}{}_{}{}^{2}P_{ia}^{},`$ the Ricci d-tensor is non symmetric.
Having defined a d-metric of type (19) in $`V^{(3+1)}`$ we can compute the scalar curvature
$$\stackrel{}{R}=g^{\beta \gamma }R_{\beta \gamma }.$$
of a d-connection $`D,`$
$$\stackrel{}{R}=G^{\alpha \beta }R_{\alpha \beta }=\widehat{R}+S,$$
(45)
where $`\widehat{R}=g^{ij}R_{ij}`$ and $`S=h^{ab}S_{ab}.`$
Now, by introducing the values (43) and (45) into the Einstein’s equations
$$R_{\beta \gamma }\frac{1}{2}g_{\beta \gamma }\stackrel{}{R}=k\mathrm{{\rm Y}}_{\beta \gamma },$$
we can write down the system of field equations for la–gravity with prescribed anholonomic (N–connection) structure :
$`R_{ij}{\displaystyle \frac{1}{2}}\left(\widehat{R}+S\right)g_{ij}`$ $`=`$ $`k\mathrm{{\rm Y}}_{ij},`$ (46)
$`S_{ab}{\displaystyle \frac{1}{2}}\left(\widehat{R}+S\right)h_{ab}`$ $`=`$ $`k\mathrm{{\rm Y}}_{ab},`$ (47)
$`{}_{}{}^{1}P_{ai}^{}`$ $`=`$ $`k\mathrm{{\rm Y}}_{ai},`$ (48)
$`{}_{}{}^{2}P_{ia}^{}`$ $`=`$ $`k\mathrm{{\rm Y}}_{ia},`$ (49)
where $`\mathrm{{\rm Y}}_{ij},\mathrm{{\rm Y}}_{ab},\mathrm{{\rm Y}}_{ai}`$ and $`\mathrm{{\rm Y}}_{ia}`$ are the components of the energy–momentum d–tensor field $`\mathrm{{\rm Y}}_{\beta \gamma }`$ (which includes possible cosmological constants, contributions of anholonomy d–torsions (28) and matter) and $`k`$ is the coupling constant
## III 4D Anholonomic Superpositions of 2D D–Metrics
Let us consider a 4D spacetime $`V^{(3+1)}`$ provided with a d–metric (19) when $`g_i=g_i(x^k)`$ and $`h_a=h_a(x^k,z)`$ for $`y^a=(z,y^4).`$ The N–connection coefficients are restricted to be some functions on three coordinates $`(x^i,z),`$
$`N_1^3`$ $`=`$ $`q_1(x^i,z),N_2^3=q_2(x^i,z),`$ (50)
$`N_1^4`$ $`=`$ $`n_1(x^i,z),N_2^4=n_2(x^i,z).`$ (51)
For simplicity, we shall use brief denotations of partial derivatives, like $`\dot{a}`$$`=a/x^1,a^{}=a/x^2,`$ $`a^{}=a/z`$ $`\dot{a}^{}`$$`=^2a/x^1x^2,`$ $`a^{}=^2a/zz.`$
The non–trivial components of the Ricci d–tensor (43), for the mentioned type of d–metrics depending on three variables, are
$`R_1^1`$ $`=R_2^2={\displaystyle \frac{1}{2g_1g_2}}\times `$ (53)
$`[(g_1^{^{\prime \prime }}+\ddot{g}_2)+{\displaystyle \frac{1}{2g_2}}\left(\dot{g}_2^2+g_1^{}g_2^{}\right)+{\displaystyle \frac{1}{2g_1}}\left(g_1^{2}+\dot{g}_1\dot{g}_2\right)];`$
$$S_3^3=S_4^4=\frac{1}{h_3h_4}[h_4^{}+\frac{1}{2h_4}(h_4^{})^2+\frac{1}{2h_3}h_3^{}h_4^{}];$$
(54)
$`P_{31}`$ $`=`$ $`{\displaystyle \frac{q_1}{2}}[\left({\displaystyle \frac{h_3^{}}{h_3}}\right)^2{\displaystyle \frac{h_3^{}}{h_3}}+{\displaystyle \frac{h_4^{}}{2h_4^{2}}}{\displaystyle \frac{h_3^{}h_4^{}}{2h_3h_4}}]`$ (56)
$`+{\displaystyle \frac{1}{2h_4}}[{\displaystyle \frac{\dot{h}_4}{2h_4}}h_4^{}\dot{h}_4^{}+{\displaystyle \frac{\dot{h}_3}{2h_3}}h_4^{}],`$
$`P_{32}`$ $`=`$ $`{\displaystyle \frac{q_2}{2}}[\left({\displaystyle \frac{h_3^{}}{h_3}}\right)^2{\displaystyle \frac{h_3^{}}{h_3}}+{\displaystyle \frac{h_4^{}}{2h_4^{2}}}{\displaystyle \frac{h_3^{}h_4^{}}{2h_3h_4}}]`$ (58)
$`+{\displaystyle \frac{1}{2h_4}}[{\displaystyle \frac{h_4^{}}{2h_4}}h_4^{}h_4^{}+{\displaystyle \frac{h_3^{}}{2h_3}}h_4^{}];`$
$`P_{41}`$ $`=`$ $`{\displaystyle \frac{h_4}{2h_3}}n_1^{},`$ (59)
$`P_{42}`$ $`=`$ $`{\displaystyle \frac{h_4}{2h_3}}n_2^{}.`$ (60)
The curvature scalar $`\stackrel{}{R}`$ (45) is defined by two non-trivial components $`\widehat{R}=2R_1^1`$ and $`S=2S_3^3.`$
The system of Einstein equations (46) transforms into
$`R_1^1`$ $`=`$ $`\kappa \mathrm{{\rm Y}}_3^3=\kappa \mathrm{{\rm Y}}_4^4,`$ (61)
$`S_3^3`$ $`=`$ $`\kappa \mathrm{{\rm Y}}_1^1=\kappa \mathrm{{\rm Y}}_2^2,`$ (62)
$`P_{3i}`$ $`=`$ $`\kappa \mathrm{{\rm Y}}_{3i},`$ (63)
$`P_{4i}`$ $`=`$ $`\kappa \mathrm{{\rm Y}}_{4i},`$ (64)
where the values of $`R_1^1,S_3^3,P_{ai},`$ are taken respectively from (53), (54), (56), (59).
By using the equations (63) and (64) we can define the N–coefficients (50), $`q_i(x^k,z)`$ and $`n_i(x^k,z),`$ if the functions $`h_i(x^k,z)`$ are known as solutions of the equations (62).
Now, we discuss the question on possible signatures of generated 4D metrics. There are three classes of la–solutions:
1. The horizontal d–metric is fixed to be of Lorentzian signature, $`sign\left(g_{ij}\right)=(,+),`$ the vertical one is of Euclidean signature, $`sign\left(h_{ab}\right)=(+,+),`$ and the resulting 4D metric $`g_{\alpha \beta }`$ will be considered of signature $`(,+,+,+).`$ The local coordinates are chosen $`u^\alpha =(x^1=t,x^2,y^3=z,y^4),`$ where $`t`$ is the time like coordinate and the d–metrics are parametrized
$$g_{ij}(t,x^1)=\left(\begin{array}{cc}g_1=\mathrm{exp}a_1& 0\\ 0& g_2=\mathrm{exp}a_2\end{array}\right)$$
(65)
and
$$h_{ab}(t,x^1,z)=\left(\begin{array}{cc}h_3=\mathrm{exp}b_3& 0\\ 0& h_4=\mathrm{exp}b_4\end{array}\right).$$
(66)
The energy–momentum d–tensor for the Einstein equations (46) could be considered in diagonal form
$$\mathrm{{\rm Y}}_\beta ^\alpha =diag[\epsilon ,p_2,p_3,p_4]$$
(67)
if the N–coefficients $`N_i^a(t,x^1,z)`$ are chosen to make zero the non–diagonal components of the Ricci d–tensor (see (56) and (59)). Here we note that on la–spacetimes, with respect to anholonomic frames, there are possible nonzero values of pressure, $`p0,`$ even $`\epsilon =0.`$
2. The horizontal d–metric is fixed to be of Euclidean signature, $`sign\left(g_{ij}\right)=(+,+),`$ the vertical one is of Lorentzian signature, $`sign\left(h_{ab}\right)=(+,)`$ and the resulting 4D metric $`g_{\alpha \beta }`$ will be considered to be a static one with signature $`(+,+,+,).`$ The local coordinates are chosen $`u^\alpha =(x^1,x^2,y^3=z,y^4=t)`$ and the d–metrics are parametrized
$$g_{ij}\left(x^k\right)=\left(\begin{array}{cc}g_1=\mathrm{exp}a_1& 0\\ 0& g_2=\mathrm{exp}a_2\end{array}\right)$$
(68)
and
$$h_{ab}(x^k,z)=\left(\begin{array}{cc}h_3=\mathrm{exp}b_3& 0\\ 0& h_4=\mathrm{exp}b_4\end{array}\right).$$
(69)
The energy–momentum d–tensor for the Einstein equations (46) could be considered in diagonal form
$$\mathrm{{\rm Y}}_\beta ^\alpha =diag[p_1,p_2,p_3,\epsilon ]$$
(70)
if the coefficients $`N_i^a(x^i,z)`$ are chosen to diagonalize the Ricci d–tensor (when (56) and (59) are zero).
3. The horizontal d–metric is fixed to be of Euclidean signature, $`sign\left(g_{ij}\right)=(+,+),`$ the vertical one is of Lorentzian signature, $`sign\left(h_{ab}\right)=(,+).`$ The local coordinates are chosen $`u^\alpha =(x^1,x^2,y^3=z=t,y^4)`$ and the d–metrics are parametrized
$$g_{ij}\left(x^k\right)=\left(\begin{array}{cc}g_1=\mathrm{exp}a_1& 0\\ 0& g_2=\mathrm{exp}a_2\end{array}\right)$$
(71)
and
$$h_{ab}(x^k,t)=\left(\begin{array}{cc}h_3=\mathrm{exp}b_3& 0\\ 0& h_4=\mathrm{exp}b_4\end{array}\right).$$
(72)
The energy–momentum d–tensor for the Einstein equations (46) is considered in diagonal form
$$\mathrm{{\rm Y}}_\beta ^\alpha =diag[p_1,p_2,\epsilon ,p_4]$$
(73)
if the N–coefficients $`N_j^a(x^i,t)`$ make zero the non–diagonal components of the Ricci d–tensor (with vanishing (56) and (59)).
The following Sections are devoted to a general study and explicit constructions of 4D solutions of the Einstein equations via type 1–3 nonlinear superpositions of 2D soliton–dilaton–black hole d–metrics.
## IV Locally Anisotropic Soliton Like Equations
We have found in the last Section that the vertical component of energy–momentum d–tensor is the non–vacuum source of the horizontal components of a d–metric (see the equations (61)) and, inversely, following (62), one could conclude that the horizontal component of energy–momentum d–tensor is the non–vacuum source of the vertical components of a d–metric. The horizontal and vertical components of the distinguished Einstein equations (46) are rather different by structures and this is taken into account by choosing the metric ansatz (5) which gives rise in a very simplified form of partial differential equations. If the 2D h–metric depends on two variables, $`g_{ij}=g_{ij}(x^k),`$ with the diagonal components satisfying a second order partial differential equation with respect to ’dot’ and ’prime’ derivatives, the 2D v–metric could be on three variables, $`h_{ab}=h_{ab}(x^k,z),`$ with diagonal components satisfying a second order partial derivation with respect to ’star’ derivatives. The purpose of this Section is to prove that both type of Einstein h– and v–equations admit soliton like solutions.
### A Horizontal La–Deformed Sine–Gordon Equations
Let us parametrize the horizontal part of d–metric (h–metric) as
$`g_1`$ $`=`$ $`ϵ\mathrm{sin}^2\left[v\left(x^i\right)/2\right],ϵ=\pm 1,`$ (74)
$`g_2`$ $`=`$ $`\mathrm{cos}^2\left[v\left(x^i\right)/2\right].`$ (75)
The non–trivial component of the Ricci d–tensor (53) is written
$$ϵR_1^1=\frac{1}{\mathrm{sin}v}\left(\ddot{v}ϵv^{\prime \prime }\right)+\rho \left(x^i\right)$$
(76)
where
$$\rho \left(x^i\right)=\rho (v,\dot{v},v^{})=\frac{\mathrm{cos}v1}{\mathrm{sin}^2v}\left(\dot{v}^2ϵv^{2}\right).$$
(77)
The horizontal Einstein equations (61) are
$$\left(\ddot{v}ϵv^{\prime \prime }\right)+\frac{\mathrm{cos}v1}{\mathrm{sin}v}\left(\dot{v}^2ϵv^{2}\right)=ϵ\kappa \mathrm{{\rm Y}}_3^3\mathrm{sin}v,$$
(78)
being compatible for matter states when $`\mathrm{{\rm Y}}_3^3=\mathrm{{\rm Y}}_4^4.`$ For simplicity, we consider the case of constant energy density or pressure (depending on the type of fixed signature), $`\mathrm{{\rm Y}}_3^3=\mathrm{{\rm Y}}_3=const.`$ The equation (78) defines some components of a 4D metric (8) and is defined by a locally anisotropic deformation of the Euclidean (for $`ϵ=1),`$ or Lorentzian (for $`ϵ=1),`$ sine–Gordon equation (2), which are related with 2D (pseudo) Riemannian metrics (1) and constant curvatures (3).
The first term from (76) transforms into a negative constant, $`\stackrel{~}{m}^2,`$ if the function $`v\left(x^i\right)`$ is chosen to be a soliton type one which solves the sine–Gordon equation (2). The physical interpretation of terms depends of the type of the solutions we try to construct (see below). We note that if $`ϵ=1,`$ the second term, $`\rho (x^i),`$ from (77) is connected with the energy density of the soliton wave
$$H=\frac{1}{2}\left(\dot{v}^2+v^{2}\right)+1\mathrm{cos}v.$$
The aim of this Subsection is to investigate some basic properties of the locally anisotropic sine–Gordon equation, which describes a 2D horizontal soliton–dilaton system (see the next Section) induced anholonomically via a source $`\mathrm{{\rm Y}}_3`$ in the vertical subspace.
#### 1 Lorentzian la–soliton systems of Class 1
We consider the Class 1 of Lorentzian h–metrics (65) which in local coordinates $`x^i=(t,r),`$ $`t`$ is a time like coordinate, and for $`ϵ=1`$ and $`\mathrm{{\rm Y}}_3=p_3+(1/2\kappa )\lambda ,`$ where $`p_3`$ is the anisotropic pressure in the $`z`$–direction and $`\lambda `$ is the cosmological constant, are defined by the equation
$$\left(\ddot{v}+v^{\prime \prime }\right)+\frac{\mathrm{cos}v1}{\mathrm{sin}v}\left(\dot{v}^2+v^{2}\right)=(\kappa p_3+\frac{\lambda }{2})\mathrm{sin}v,$$
(79)
where $`\dot{v}=v/t`$ and $`v^{}=v/r.`$
For $`\mathrm{cos}v1,`$ by neglecting quadratic terms $`v^2,`$ we can approximate the solution of (79) by a solution of the Euclidean 2D sine–Gordon equation (2) with the constant $`\stackrel{~}{m}^2=(\kappa p_3+\frac{\lambda }{2})`$. If we wont to treat $`\stackrel{~}{m}^2`$ as a mass like constant, we must suppose that matter is in a state for which $`(\kappa p_3+\frac{\lambda }{2})<0.`$
#### 2 Euclidean la–soliton systems of Class 2
This type of h–metrics (68), of Class 2 from the previous Section, for $`ϵ=1,`$ $`\mathrm{{\rm Y}}_3=p_3+(1/2\kappa )\lambda `$ and both space like local coordinates $`x^i,`$ is given by the following from (78) equations
$$\left(\ddot{v}v^{\prime \prime }\right)+\frac{\mathrm{cos}v1}{\mathrm{sin}v}\left(\dot{v}^2v^{2}\right)=(\kappa p_3+\frac{\lambda }{2})\mathrm{sin}v,$$
(80)
where $`\dot{v}=v/x^1`$ and $`v^{}=v/x^2,`$ which for $`\mathrm{cos}v1`$ has solutions approximated by the Lorentzian 2D sine–Gordon equation with $`\stackrel{~}{m}^2=(\kappa p_3+\frac{\lambda }{2}).`$
#### 3 Euclidean la–soliton systems of Class 3
In this case the h–metrics (71) of Class 3 from the previous Section, for $`ϵ=1,`$ $`\mathrm{{\rm Y}}_3=\epsilon +(1/2\kappa )\lambda `$ and both space like local coordinates $`x^i,`$ is defined by equations
$$\left(\ddot{v}v^{\prime \prime }\right)+\frac{\mathrm{cos}v1}{\mathrm{sin}v}\left(\dot{v}^2v^{2}\right)=(\kappa \epsilon +\frac{\lambda }{2})\mathrm{sin}v,$$
(81)
where $`\dot{v}=v/x^1`$ and $`v^{}=v/x^2,`$ which for $`\mathrm{cos}v1`$ has solutions approximated by the Lorentzian 2D sine–Gordon equation with $`\stackrel{~}{m}^2=(\kappa ϵ+\frac{\lambda }{2}).`$
#### 4 A static one dimensional exact solution
The la–deformed sine–Gordon equations can be integrated exactly for static configurations. If $`v=v_s(x^2=x),v^{}=dv_s/dx`$ the equation (78) transforms into
$$\frac{d^2v_s}{dx^2}+\frac{\mathrm{cos}v_s1}{\mathrm{sin}v_s}\left(\frac{dv_s}{dx}\right)^2+\kappa \mathrm{{\rm Y}}_3\mathrm{sin}v_s=0.$$
(82)
which does not depend on values of $`ϵ=\pm 1.`$ Introducing a new variable $`y(v_s)=\left(dv_s/dx\right)^2`$ we get a linear first order differential equation
$$\frac{dy}{dv_s}+2\frac{\mathrm{cos}v_s1}{\mathrm{sin}v_s}y+2\kappa \mathrm{{\rm Y}}_3\mathrm{sin}v_s=0$$
which is solved by applying standard methods .
### B Vertical Einstein Equations and Possible Soliton Like Solutions
The basic equation
$`{\displaystyle \frac{^2h_4}{z^2}}{\displaystyle \frac{1}{2h_4}}\left({\displaystyle \frac{h_4}{z}}\right)^2`$ (83)
$`{\displaystyle \frac{1}{2h_3}}\left({\displaystyle \frac{h_3}{z}}\right)\left({\displaystyle \frac{h_4}{z}}\right){\displaystyle \frac{\kappa }{2}}\mathrm{{\rm Y}}_1h_3h_4=0`$ (84)
(here we write down the partial derivatives on $`z`$ in explicit form) follows from (54) and (62) and relates some first and second order partial on $`z`$ derivatives of diagonal components $`h_a(x^i,z)`$ of a v–metric with a source $`\kappa \mathrm{{\rm Y}}_1(x^i,z)=\kappa \mathrm{{\rm Y}}_1^1=\kappa \mathrm{{\rm Y}}_2^2`$ in the h–subspace. We can consider as unknown the function $`h_3(x^i,z)`$ (or, inversely, $`h_4(x^i,z))`$ for some compatible values of $`h_4(x^i,z)`$ (or $`h_3(x^i,z))`$ and source $`\mathrm{{\rm Y}}_1(x^i,z).`$
The structure of equation (83) differs substantially from the horizontal one (61), or (78). In this Subsection we analyze some soliton type integral varieties which solve the partial differential equation (83).
#### 1 Belinski–Zakharov–Maison locally isotropic limits
In the vacuum case $`\mathrm{{\rm Y}}_1^10`$ and arbitrary two functions depending only on variables $`x^i`$ are admitted as solutions of (83). The N–coefficients $`q_i`$ and $`n_i`$ became zero in consequence of (63), (64) and (56), (59). So, in the locally isotropic vacuum limit the 4D metrics (5) will transform into a soliton vacuum solution of Einstein equations if, for instance, we require that $`h_a(x^i)`$ are the components of the diagonalized matrix for the Belinski–Zakharov or Maison gravitational solitons and the h–metric transforms is defined by a conformal factor $`f(x^i)`$ being compatible with the v–metric. Of course, instead of soliton ones we can choose another class of vacuum solutions depending on variables $`x^i`$ to be the locally isotropic limit of anholonomic gravitational systems.
#### 2 Kadomtsev–Petviashvili v–solitons
By straightforward verification we conclude that the v–metric component $`h_4(x^i,z)`$ could be a solution of Kadomtsev–Petviashvili (KdP) equation (the first methods of integration of 2+1 dimensional soliton equations where developed by Dryuma and Zakharov and Shabat )
$$h_4^{}+ϵ\left(\dot{h}_4+6h_4h_4^{}+h_4^{\prime \prime \prime }\right)^{}=0,ϵ=\pm 1,$$
(85)
if the component $`h_3(x^i,z)`$ satisfies the Bernoulli equations
$$h_3^{}+Y(x^i,z)(h_3)^2+F_ϵ(x^i,z)h_3=0,$$
(86)
where, for $`h_4^{}0,`$
$$Y(x^i,z)=\kappa \mathrm{{\rm Y}}_1^1\frac{h_4}{h_4^{}},$$
(87)
and
$$F_ϵ(x^i,z)=\frac{h_4^{}}{h_4}+\frac{2ϵ}{h_4^{}}\left(\dot{h}_4+6h_4h_4^{}+h_4^{\prime \prime \prime }\right)^{}.$$
The three dimensional integral variety of (86) is defined by formulas
$$h_3^1(x^i,z)=h_{3(x)}^1\left(x^i\right)E_ϵ(x^i,z)\times \frac{Y(x^i,z)}{E_ϵ(x^i,z)}𝑑z,$$
where
$$E_ϵ(x^i,z)=\mathrm{exp}F_ϵ(x^i,z)𝑑z$$
and $`h_{3(x)}\left(x^i\right)`$ is a nonvanishing function.
In the vacuum case $`Y(x^i,z)=0`$ and we can write the integral variety of (86)
$$h_3^{(vac)}(x^i,z)=h_{3(x)}^{(vac)}\left(x^i\right)\mathrm{exp}\left[F_ϵ(x^i,z)𝑑z\right].$$
We conclude that a solution of KdP equation (86) could generate a non–perturbative component $`h_4(x^i,z)`$ of a diagonal h–metric if the second component $`h_3(x^i,z)`$ is a solution of Bernoulli equations (86) with coefficients determined both by $`h_4`$ and its partial derivatives and by the $`\mathrm{{\rm Y}}_1^1`$ component of the energy–momentum d–tensor (see (87)). In the non–vacuum case the parameters of (2+1) dimensional KdV solitons are connected with parameters defining the interactions with matter fields and/or by a cosmological constant. The further developments in this direction consist in construction of self–consistent (2+1) KdV soliton solutions induced by some soliton configurations from energy–momentum tensor in hydrodynamical (plasma) approximations.
#### 3 (2+1) sine–Gordon v–solitons
In a symilar manner as in previous paragraph we can prove that solutions $`h_4(x^i,z)`$ of (2+1) sine–Gordon equation (see, for instance, )
$$h_4^{}+h_4^{^{\prime \prime }}\ddot{h}_4=\mathrm{sin}(h_4)$$
also induce solutions for $`h_3(x^i,z)`$ following from the Bernoulli equation
$$h_3^{}+\kappa \mathrm{{\rm Y}}_1(x^i,z)\frac{h_4}{h_4^{}}(h_3)^2+F(x^i,z)h_3=0,h_4^{}0,$$
where
$$F(x^i,z)=\frac{h_4^{}}{h_4}+\frac{2}{h_4^{}}\left[h_4^{^{\prime \prime }}\ddot{h}_4\mathrm{sin}(h_4)\right].$$
The integral varieties (with energy–momentum sources and in vacuum cases) are constructed by a corresponding redefinition of coefficients in formulas from the previous paragraph.
#### 4 On some general properties of h–metrics depending on 2+1 variables
By introducing a new variable $`\beta =h_4^{}/h_4`$ the equation (83) transforms into
$$\beta ^{}+\frac{1}{2}\beta ^2\frac{\beta h_3^{}}{2h_3}2\kappa \mathrm{{\rm Y}}_1h_3=0$$
(88)
which relates two functions $`\beta (x^i,z)`$ and $`h_3(x^i,z).`$ There are two possibilities: 1) to define $`\beta `$ (i. e. $`h_4)`$ when $`\kappa \mathrm{{\rm Y}}_1`$ and $`h_3`$ are prescribed and, inversely 2) to find $`h_3`$ for given $`\kappa \mathrm{{\rm Y}}_1`$ and $`h_4`$ (i. e. $`\beta );`$ in both cases one considers only ”\*” derivatives on $`z`$–variable with coordinates $`x^i`$ treated as parameters.
1. In the first case the explicit solutions of (88) have to be constructed by using the integral varieties of the general Riccati equation which by a corresponding redefinition of variables, $`zz\left(\varsigma \right)`$ and $`\beta \left(z\right)\eta \left(\varsigma \right)`$ (for simplicity, we omit dependencies on $`x^i)`$ could be written in the canonical form
$$\frac{\eta }{\varsigma }+\eta ^2+\mathrm{\Psi }\left(\varsigma \right)=0$$
where $`\mathrm{\Psi }`$ vanishes for vacuum gravitational fields. In vacuum cases the Riccati equation reduces to a Bernoulli equation which (we can use the former variables) for $`s(z)=\beta ^1`$ transforms into a linear differential (on $`z)`$ equation,
$$s^{}+\frac{h_3^{}}{2h_3}s\frac{1}{2}=0.$$
(89)
2. In the second (inverse) case when $`h_3`$ is to be found for some prescribed $`\kappa \mathrm{{\rm Y}}_1`$ and $`\beta `$ the equation (88) is to be treated as a Bernoulli type equation,
$$h_3^{}=\frac{4\kappa \mathrm{{\rm Y}}_1}{\beta }(h_3)^2+\left(\frac{2\beta ^{}}{\beta }+\beta \right)h_3$$
(90)
which can be solved by standard methods. In the vacuum case the squared on $`h_3`$ term vanishes and we obtain a linear differential (on $`z)`$ equation.
#### 5 A class of conformally equivalent h–metrics
A particular interest presents those solutions of the equation (88) which via 2D conformal transforms with a factor $`\omega =\omega (x^i,z)`$ are equivalent to a diagonal h–metric on $`x`$–variables, i.e. one holds the parametrization
$$h_3=\omega (x^i,z)a_3\left(x^i\right)\text{ and }h_4=\omega (x^i,z)a_4\left(x^i\right),$$
(91)
where $`a_3\left(x^i\right)`$ and $`a_4\left(x^i\right)`$ are some arbitrary functions (for instance, we can impose the condition that they describe some 2D soliton like or black hole solutions). In this case $`\beta =\omega ^{}/\omega `$ and for $`\gamma =\omega ^1`$ the equation (88) trasforms into
$$\gamma \gamma ^{}=2\kappa \mathrm{{\rm Y}}_1a_3\left(x^i\right)$$
(92)
with the integral variety determined by
$$z=\frac{d\gamma }{\sqrt{|4k\mathrm{{\rm Y}}_1a_3(x^i)\mathrm{ln}|\gamma |+C_1(x^i)|}}+C_2(x^i),$$
(93)
where it is considered that the source $`\mathrm{{\rm Y}}_1`$ does not depend on $`z.`$
Finally, in this Section, we have shown that a large class of 4D solutions, depending on two or three variables, of the Einstein equations can be constructed as nonlinear superpositions of some 2D h–metrics defined by locally anisotropic deformations of 2D sine–Gordon equations and of some v–metrics generated in particular by solutions of Kadomtsev–Petviashvili equations, or of (2+1) sine–Gordon, and associated Bernoulli type equations. From a general viewpoint the v–metrics are defined by integral varieties of corresponding Riccati and/or Bernoulli equations with respect to $`z`$–variables with the h–coordinates $`x^i`$ treated as parameters.
## V Effective Locally Anisotropic Soliton–Dilaton fields
The formula for the h–component $`\widehat{R}`$ of scalar curvature (see (45) and (53)) of a h–metric (6), written for a la–system, differs from the usual one for computation of curvature of 2D metrics. That why additionally to the first term in (76) it is induced the $`\rho `$–term (77). The aim of this Section is to prove that the la–deformed singe–Gordon equation (76) could be equivalently modelled by solutions of the usual 2D singe–Gordon equation and an additional equation for a corresponding effective dilaton field (in brief, by a soliton–dilaton field). We also analyze the 2D dilaton gravity in connection with the sine–Gordon la–field theory.
### A Generic locally anisotropic dilaton fields
Let $`\stackrel{~}{g}_{ij}^ϵ\left(x^i\right)`$ be a 2D metric of Lorentz (or Euclidean) signature for $`ϵ=1`$ (or $`ϵ=1)`$ with a usual 2D scalar curvature $`\stackrel{~}{R}_{(ϵ)}\left(x^i\right).`$ We also consider a conformally equivalent metric
$$\underset{¯}{g}_{ij}^ϵ\left(x^i\right)=\mathrm{exp}\omega \left(x^i\right)\stackrel{~}{g}_{ij}^ϵ\left(x^i\right).$$
(94)
The scalar curvatures of 2D metrics from (94) are related by the formula
$$e^\omega \underset{¯}{R}=\stackrel{~}{R}_{(ϵ)}+\mathrm{}_{(ϵ)}\omega $$
(95)
where $`\mathrm{}_{(ϵ)}`$ is the d’Alambert, $`ϵ=1,`$ (Laplace,$`ϵ=1`$) operator.
In order to model a locally anisotropic 2D horizontal system via a locally isotropic 2D gravity we consider that
$$\widehat{R}=e^\omega \underset{¯}{R},\stackrel{~}{R}_{(ϵ)}=2\stackrel{~}{m}^2$$
and
$$\mathrm{}_{(ϵ)}\omega =\rho \left(x^i\right).$$
(96)
For a given ’tilded’ metric, for instance, $`\stackrel{~}{g}_{ij}=diag(\stackrel{~}{a},\stackrel{~}{b})`$ being a solution of 2D sine–Gordon equation (2) with negative constant scalar curvature (see (3)), the wave (Poisson) equation can be solved in explicit form by imposing corresponding boundary conditions.
So, a 2D locally anisotropic h-space is equivalently modelled by a usual curved 2D locally isotropic (pseudo) Riemannian space and effective interactions with the generic locally anisotropic dilaton field $`\mathrm{\Phi }_{(\omega )}=\mathrm{exp}\omega .`$
### B Locally anisotropic 2D dilaton gravity and sine–Gordon theory
In the previous Subsection the conformal factor $`\mathrm{\Phi }_{(\omega )},`$ in the h–space, was introduced with the aim to compensate the local anisotropy, induced from the v–space. The 2D h–gravity can be formulated as a dilatonic one related to a generalized, la–deformed, sine–Gordon model.
By using Weyl rescallings of h–metric (6) one can write the general action for, in our case, the h–model (see the isotropic variant in and ),
$$S^{[h]}[g_{ij},\mathrm{\Phi }]=\frac{1}{2\pi }d^2x\sqrt{g}[\mathrm{\Phi }\widehat{R}+\varpi ^2V\left(\mathrm{\Phi }\right)],$$
(97)
where the h–metric $`g_{ij}`$ has signature $`(1,1),`$ $`V\left(\mathrm{\Phi }\right)`$ is an arbitrary function of the dilaton field $`\mathrm{\Phi }`$ and $`\varpi `$ is the connection constant. The la–field equations derived from this action are
$`\widehat{R}=\stackrel{~}{R}_{()}+\mathrm{}_{()}\omega `$ $`=`$ $`\varpi ^2{\displaystyle \frac{dV}{d\mathrm{\Phi }}},`$ (98)
$`D_iD_j\mathrm{\Phi }{\displaystyle \frac{\varpi ^2}{2}}g_{ij}V`$ $`=`$ $`0,`$ (99)
where $`\stackrel{~}{R}_{()}=2\stackrel{~}{R}_1^1`$ is defined by the h–component of scalar curvature of type (76), when $`ϵ=1.`$ In the locally isotropic limit this system of equations describes the Cadoni theory .
In consequence of the fact that the theory is invariant under coordinate h–transforms $`(x^1=t,x^2=x)`$ we can introduce the h–metric
$$g^{[h]}=\mathrm{sin}^2\left(\frac{v}{2}\right)dt^2+\mathrm{cos}^2\left(\frac{v}{2}\right)dx^2,$$
(100)
where $`v=v(t,x)`$ and rewrite the system (98) as a system of nonlinear partial differential equations in 2D Euclidean space,
$`\ddot{v}+v^{\prime \prime }`$ $`=`$ $`\left(\rho +{\displaystyle \frac{\varpi ^2}{2}}{\displaystyle \frac{dV}{d\mathrm{\Phi }}}\right)\mathrm{sin}v,`$ (101)
$`\ddot{\mathrm{\Phi }}+\mathrm{\Phi }^{\prime \prime }`$ $`=`$ $`{\displaystyle \frac{\varpi ^2}{2}}V\mathrm{cos}v,`$ (102)
where the function $`\rho `$ is defined by the formula (77) for $`ϵ=1,`$ or, in equivalent form, by a generic la–dilaton given by (96).
The equation (101) reduces to the deformed sine–Gordon equation (78), for $`ϵ=1,`$ if $`V=\mathrm{\Phi }`$, or for constant configurations $`\mathrm{\Phi }_0`$ for which $`V(\mathrm{\Phi }_0)=0`$ and
$`dV/d\mathrm{\Phi }|_{\mathrm{\Phi }_0}>0.`$
For soliton–dilaton la–configurations it is more convenient to consider the action
$$S=\frac{1}{2}d^2x\left[\mathrm{\Phi }\left(\mathrm{}_{()}v+\rho \right)\frac{\varpi ^2}{2}V\mathrm{sin}v\right],$$
(103)
given in the 2D Minkowski space, where $`\mathrm{}_{()}v=\ddot{v}+v^{\prime \prime }`$ and $`\rho =\mathrm{}_{()}\omega .`$ Extremizing the action (103) we obtain the field equations (101) and (102) as well from this action one follows the energy functional
$`E(v,\omega ,\mathrm{\Phi })`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}dx[\dot{\mathrm{\Phi }}(\dot{v}+\dot{\omega })+\mathrm{\Phi }^{}(v^{}+\omega ^{})`$ (105)
$`+{\displaystyle \frac{\varpi ^2}{2}}V\mathrm{sin}v].`$
We note that instead of Lorentz type 2D h–metrics we can consider Euclidean field equations by performing the Wick rotation $`tit.`$ We conclude this Subsection by the remark that a complete correspondence between locally anisotropic h–metrics and dilaton structures is possible if additionally to trigonometric parametrizations of 2D metrics (100) one introduces hyperbolic parametizations
$$g^{[h]}=\mathrm{sinh}^2\left(\frac{v}{2}\right)dt^2+\mathrm{cosh}^2\left(\frac{v}{2}\right)dx^2$$
which results in sinh–Gordon models.
### C Static locally anisotropic soliton–dilaton configurations
Because the $`\rho `$–term (77) vanishes for constant values of fields $`v=2n\pi ,`$ $`n=0,\pm 1,\pm 2,\mathrm{}`$ and $`\mathrm{\Phi }=\mathrm{\Phi }_0`$ the vacua of the model (97) is singled out like in the locally isotropic case , by conditions
$$V\left(\mathrm{\Phi }_0\right)=0\text{ and }\frac{dV}{d\mathrm{\Phi }}\left(\mathrm{\Phi }_0\right)>0.$$
(106)
In order to focus on static deformations induced by soliton like solutions we require $`E0`$ and
$$\underset{x\pm \mathrm{}.}{lim}v^{}0\text{ and }\underset{x\pm \mathrm{}.}{lim}\mathrm{\Phi }^{}0.$$
(107)
The static la–configurations of (101) and (102) are giving by anholonomic deformation of isotropic ones and are given by the system of equations
$`v^{\prime \prime }+{\displaystyle \frac{\mathrm{cos}v1}{\mathrm{sin}v}}\left(v^{}\right)^2`$ $`=`$ $`{\displaystyle \frac{\varpi ^2}{2}}\mathrm{sin}v{\displaystyle \frac{dV}{d\mathrm{\Phi }}},`$ (108)
$`\mathrm{\Phi }^{\prime \prime }={\displaystyle \frac{\varpi ^2}{2}}V\left(\mathrm{\Phi }\right)\mathrm{cos}v.`$ (109)
The first integrals of (108) and (109) are
$$v^{}=\varpi \frac{a_0}{\sqrt{2}}\mathrm{sin}^2\frac{v}{2}\text{ and }\mathrm{\Phi }^{}=\frac{\varpi \sqrt{2}}{a_0}\mathrm{cot}\frac{v}{2}$$
which, after another integration, results in the solutions
$`\varpi \left(xx_0\right)`$ $`=`$ $`\pm {\displaystyle \frac{a_0^2}{\sqrt{2}}}{\displaystyle 𝑑\mathrm{\Phi }\sqrt{\frac{\mathrm{\Psi }c_0}{1a_0^2\left(\mathrm{\Psi }c_0\right)}}},`$ (110)
$`\mathrm{sin}{\displaystyle \frac{v}{2}}`$ $`=`$ $`\pm a_0\sqrt{\mathrm{\Psi }c_0},`$ (111)
where $`\mathrm{\Psi }=\mathrm{\Psi }\left(\mathrm{\Phi }\right)=_0^\mathrm{\Phi }𝑑\tau V\left(\tau \right);`$ $`a_0`$ and $`c_0=\mathrm{\Psi }[\underset{¯}{\mathrm{\Phi }}\left(\pm \mathrm{}\right)]`$ are integration constants. We emphasize that the formula (110), following from a la–model, differs from that obtained in the Cadoni’s locally isotropic theory .
There are two additional two parameter solutions of (108) and (108) which are not contained in (110). The first type of solutions are those for constant $`v`$ field when
$$v=n\pi \text{ and }\varpi \left(xx_0\right)=\pm 𝑑\mathrm{\Phi }\left[(1)^n\mathrm{\Psi }b_0\right]^{1/2},$$
where $`b_0=const.`$ The second type of solutions are for constant dilaton fields $`\mathrm{\Phi }_0`$ and exists if there is at least one zero $`\mathrm{\Phi }=\mathrm{\Phi }_0`$ for $`V\left(\mathrm{\Phi }\right).`$ For $`dV/d\mathrm{\Phi }\left(\mathrm{\Phi }_0\right)>0`$ the equations reduce to the usual sine–Gordon equations
$$v^{\prime \prime }=\frac{\varpi ^2}{2}\frac{dV}{d\mathrm{\Phi }}_{\mathrm{\Phi }_0}\mathrm{sin}v.$$
Note that the model (103) admits static soliton solutions, approaching for $`x\pm \mathrm{}`$ the constant field configuration $`v=2\pi n;n=\pm 1,`$ with $`\mathrm{\Phi }_0=\mathrm{\Phi }\left(\pm \mathrm{}\right),`$ $`V\left(\mathrm{\Phi }_0\right)=0`$ and $`dV/d\mathrm{\Phi }|_{\mathrm{\Phi }_0}>0.`$
### D Topology of locally anisotropic soliton–dilatons
If we suppose that la–deformations do not change the spacetime topology, the conditions (107) imply that every soliton solution tends asymptotically to one of vacuum configurations (106) which could be considered for both locally anisotropic and isotropic systems (see ). The admissible number of solitons to be la–deformed without changing of topology is determined by the number of ways in which the points $`x=\pm \mathrm{}`$ (the zero sphere) can be mapped into the manifold of the mentioned constant–field configurations (106) characterized by the homotopy group
$$\pi _0\left(\frac{Z\times Z_2}{Z_2}\right)=\pi _0\left(Z\right)=Z,$$
when $`G=Z\times Z_2`$ is the invariance group for a generic $`V,`$ and $`Z`$ and $`Z_2`$ are respectively the infinite discrete group translations and the finite group of inversions of the field $`v,`$ parametrized by $`vv+2\pi n,`$ $`vv.`$ This result holds for usual sine–Gordon systems, as well by la–generalizations given by the action (103) and la–field equations (101) and (102).
For soliton like theories it is possible the definition of conserved currents
$$J_{(v)}^i=ϵ^{ij}\delta _iv\text{ and }J_{(\mathrm{\Phi })}^i=ϵ^{ij}\delta _i\mathrm{\Phi },$$
where $`ϵ^{ij}=ϵ^{ji}`$ and the ’elongated’ (in la–case) partial derivatives $`\delta _i`$ are given by (9). The associated topological charges on a fixed la–background are
$`Q_{(v)}`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}𝑑xJ_{(v)}^1={\displaystyle \frac{1}{2\pi }}\left[v\left(\mathrm{}\right)v\left(\mathrm{}\right)\right],`$ (112)
$`Q_{(\mathrm{\Phi })}`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}𝑑xJ_{(\mathrm{\Phi })}^1={\displaystyle \frac{1}{2\pi }}\left[\mathrm{\Phi }\left(\mathrm{}\right)\mathrm{\Phi }\left(\mathrm{}\right)\right].`$ (113)
The topological properties of la–backgroudns are characterized by the topological current and charge of la–dilaton $`\mathrm{\Phi }_{(\omega )}=\mathrm{exp}\omega `$ defined by a solution of Poisson equation (96). The corresponding formulas are
$`J_{(e^\omega )}^i`$ $`=`$ $`ϵ^{ij}\delta _ie^\omega ,`$ (114)
$`Q_{(e^\omega )}`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}𝑑xJ_{(e^\omega )}^1={\displaystyle \frac{1}{2\pi }}\left[\mathrm{exp}\omega \left(\mathrm{}\right)\mathrm{exp}\omega \left(\mathrm{}\right)\right].`$ (115)
If $`J_{(e^\omega )}^i`$ and $`Q_{(e^\omega )}`$ are non-trivial we can conclude that our soliton–dilaton system was topologically changed under la–deformations.
## VI Locally Anisotropic Black Holes and Solitons
In this Section we analyze the connection between h–metrics describing effective 2D black la–hole solutions (with parameters defined by v–components of a diagonal 4D energy–momentum d–tensor) and 2D soliton la–solutions obtained in the previous two Sections.
### A A Class 1 black hole solutions
Let us consider a static h–metric of type (65), for which $`g_1=\alpha \left(r\right)`$ and $`g_2=1/\alpha \left(r\right)`$ for a function on necessary smooth class $`\alpha `$ on h–coordinates $`x^1=T`$ and $`x^2=r.`$ Putting these values of h–metric into (53) we compute
$$R_1^1=R_2^2=\frac{1}{2}\alpha ^{\prime \prime }.$$
Considering the 2D h–subspace to be of constant negative scalar curvature,
$$\widehat{R}=2R_1^1=\stackrel{~}{m}^2,$$
and that the Einstein la–equations (61) are satisfied we obtain the relation
$$\alpha ^{\prime \prime }=\stackrel{~}{m}^2=\kappa \mathrm{{\rm Y}}_3^3=\kappa \mathrm{{\rm Y}}_4^4,$$
(116)
which for a diagonal energy–momentum d–tensor with $`\kappa \mathrm{{\rm Y}}_3^3=kp_3+\lambda /2`$ transforms into
$$\stackrel{~}{m}^2=kp_3+\frac{\lambda }{2}.$$
The solution of (116) is written in the form $`\alpha =\left(\stackrel{~}{m}^2r^2M\right)`$ which defines a 2D h–metric
$$ds_{(h)}^2=\left(\stackrel{~}{m}^2r^2M\right)dT^2+\left(\stackrel{~}{m}^2r^2M\right)^1dr^2$$
(117)
being similar to a black hole solution in 2D Jackiw–Teitelboim gravity and display many of attributes of black holes with that difference that the constant $`\stackrel{~}{m}`$ is defined by 4D physical values in v–subspace and for definiteness of the theory the h–metric should be supplied with a v–component.
The parameter $`M,`$ the mass observable, is the analogue of the Arnowitt–Deser–Misner (ADM) mass in general relativity. If we associate the h–metric (117) to a 2D model of Jackiw–Teitelboim la–gravity, given by the action
$$I_{JT}[\varphi ,g]=\frac{1}{2G_{(2)}}_Hd^2x\sqrt{\left|g\right|}\varphi \left(\widehat{R}+2\stackrel{~}{m}^2\right)$$
where $`\widehat{R}`$ is the h–component of the Ricci d–curvature, $`\varphi `$ is the dilaton field and $`G_{(2)}`$ is the gravitational coupling constant in 2D, we should add to (116) the field equation for $`\varphi ,`$
$$\left(D_iD_j\stackrel{~}{m}^2g_{ij}\right)\varphi =0$$
which has the solution
$$\varphi =c_1\stackrel{~}{m}r$$
with coupling constant $`c_1,`$ (we can consider $`c_1=1`$ for the vacuum configurations $`\varphi =\stackrel{~}{m}r`$ as $`r\mathrm{}).`$ In this case the mass observable is connected with the dilaton as
$$M=\stackrel{~}{m}^2\left|D\varphi \right|^2+\varphi ^2.$$
(118)
Clearly this model is with local anisotropy because the value $`\widehat{R}`$ is defined in a la–manner and not as a usual scalar curvature in 2D gravity.
### B La–deformed soliton–dilaton systems and black la–holes
Suppose we have a h–metric (100) which must solve the equations
$`\ddot{v}+v^{\prime \prime }`$ $`=`$ $`\left(\rho +\stackrel{~}{m}^2\right)\mathrm{sin}v,`$ (119)
$`\ddot{\varphi }+\varphi ^{\prime \prime }`$ $`=`$ $`\stackrel{~}{m}^2\mathrm{cos}v.`$ (120)
These equations are a particular case of the system (101) and (102). For $`\rho =0`$ such equations were investigated in . In the previous Section we concluded that every 2D la–system can be equivalently modelled in an isotropic space by considering an effective interaction with la–dilaton field. The same considerations hold good for 2D la–spaces with that remark that the dilaton field $`\varphi `$ must be composed from a component satisfying the equation (120) and another component defined from the Poisson equation (96).
Having a dilaton field $`\varphi (t,x)`$ we can introduce a new ”radial coordinate”
$$r(t,x):=\varphi /\stackrel{~}{m}$$
which (being substituted into h–metric (100)) results in the horizontal 2D metric
$$ds_{[h]}^2=\stackrel{~}{m}^2\left|D\varphi \right|^2dT^2+\stackrel{~}{m}^2\left|D\varphi \right|^2dr^2,$$
(121)
where
$$dT\left|D\varphi \right|^2\left(\varphi ^{}\mathrm{tan}\frac{v}{2}dt+\dot{\varphi }\mathrm{cot}\frac{v}{2}dx\right).$$
The metric (121) is the same as (117) because the mass observable is defined by (118).
### C The geometry of black la–holes and deformed one–soliton solutions
The one soliton solution of the Euclidean sine–Gordon equation can be written as
$$v(t,x)=4\mathrm{tan}^1\left\{\mathrm{exp}\left[\pm \stackrel{~}{m}\gamma \left(xstx_{(0)}\right)\right]\right\},$$
(122)
where $`\gamma =1/\sqrt{1+s^2},`$ $`s`$ is the spectral parameter and the integration constant $`x_{(0)}`$ is considered, for simplicity, zero. The solution with signs $`+/`$ gives a soliton/anti–soliton configuration. Let us demonstrate that a corresponding black la–hole can be constructed. Putting (122) into the h–metric (100) we obtain a Lorentzian one–soliton 2D metric
$$ds_{[1sol]}^2=\mathrm{sec}h^2\xi dt^2+\mathrm{tanh}^2\xi dx^2$$
where
$$\xi \stackrel{~}{m}\gamma \left(xst\right).$$
In a similar fashion we can compute (by using the function (122)) the la–deformation $`\rho `$ (77) and effective la–dilaton $`\mathrm{\Phi }_\omega =\mathrm{exp}\omega `$ which follow from the Poisson equation (96). In both cases of dilaton equations we are dealing with linear partial differential equations. A combination of type
$$\varphi =\varphi _{[0]}\dot{v}+\varphi _{[1]}v^{}$$
(123)
for arbitrary constants $`\varphi _{[0]}`$ and $`\varphi _{[1]}`$ satisfies the linearized sine–Gordon equation and because for the function (122)
$$\dot{v}=4\stackrel{~}{m}\gamma s\mathrm{sec}h\xi =sv^{}$$
we can put in (123) $`\varphi _{[1]}=0`$ and following a Hamiltonian analysis (in order to have compatibility with the locally isotropic case ; for a black hole mass $`M=s^2`$ with corresponding ADM energy $`E=\stackrel{~}{m}^2s^2/\left(2G_{(2)}\right))`$ we set $`\varphi _{[0]}=1/\left(4\stackrel{~}{m}\gamma ^2\right)`$ so that
$$\varphi =\left|\frac{s}{\gamma }\right|\mathrm{sec}h\xi $$
is chosen to make $`\varphi `$ positive. In consequence, the black hole coordinates $`(r,T)`$ (la–deformations reduces to reparametrization of such coordinates) are defined by
$$r=\varphi /\stackrel{~}{m}=\frac{s}{\stackrel{~}{m}\gamma }\mathrm{sec}h\xi $$
and
$$dT=\left|s\stackrel{~}{m}\right|^1\left[dt\frac{s\mathrm{tanh}^2\xi }{\stackrel{~}{m}\gamma \left(\mathrm{sec}h^2\xi s^2\mathrm{tanh}^2\xi \right)}d\xi \right].$$
With respect to these coordinates the obtained black hole metric is of the form
$$ds_{[bh]}^2=\left(\stackrel{~}{m}^2r^2\stackrel{~}{m}^2s^4\right)dT^2+\left(\stackrel{~}{m}^2r^2\stackrel{~}{m}^2s^4\right)^1dr^2$$
which describes a Jackiw–Teitelboim black hole with mass parameter (118)
$$M_{1sol}=\stackrel{~}{m}^2s^4,$$
defined by the corresponding component $`\mathrm{{\rm Y}}_3^3=\mathrm{{\rm Y}}_4^4`$ of energy–momentum d–tensor in v–space and spectral parameter $`s`$ of the one soliton background, and event horizon at $`\varphi =\varphi _H=s^2.`$
In a similar fashion we can use instead of the function (122) a two and, even multi–, soliton background. The calculus is similar to the locally isotropic case , having some redefinitions of black hole coordinates if it is considered that la–deformations do not change the h–spaces topology, i.e the la–gravitational topological source and charge (114) vanishes.
## VII 3D Black La–Holes
Let us analyze some basic properties of 3D spacetime $`V^{(2+1)}`$ provided with d–metrics of type
$$\delta s^2=g_1\left(x^k\right)\left(dx^1\right)^2+g_2\left(x^k\right)\left(dx^2\right)^2+h_3(x^i,z)\left(\delta z\right)^2,$$
(124)
where $`x^k`$ are 2D coordinates, $`y^3=z`$ is the anisotropic coordinate and
$$\delta z=dz+N_i^3(x^k,z)dx^i.$$
The N–connection coefficients are given by some functions on $`x^i`$ and $`z,`$
$$N_1^3=q_1(x^i,z),N_2^3=q_2(x^i,z).$$
(125)
The non–trivial components of the Ricci d–tensor (43) are
$`R_1^1`$ $`=R_2^2={\displaystyle \frac{1}{2g_1g_2}}[(g_1^{^{\prime \prime }}+\ddot{g}_2)`$ (127)
$`+{\displaystyle \frac{1}{2g_2}}(\dot{g}_2^2+g_1^{}g_2^{})+{\displaystyle \frac{1}{2g_1}}(g_1^{2}+\dot{g}_1\dot{g}_2)];`$
$$P_{3i}=\frac{q_i}{2}[\left(\frac{h_3^{}}{h_3}\right)^2\frac{h_3^{}}{h_3}]$$
(128)
(for 3D the component $`S_3^30,`$ see (54)).
The curvature scalar $`\stackrel{}{R}`$ (45) is $`\stackrel{}{R}=\widehat{R}=2R_1^1.`$
The system of Einstein equations (46) transforms into
$`R_1^1`$ $`=`$ $`\kappa \mathrm{{\rm Y}}_3^3,`$ (129)
$`P_{3i}`$ $`=`$ $`\kappa \mathrm{{\rm Y}}_{3i},`$ (130)
which is compatible for if the 3D matter is in a state when for the energy–momentum d–tensor $`\mathrm{{\rm Y}}_\beta ^\alpha `$ one holds $`\mathrm{{\rm Y}}_1^1=\mathrm{{\rm Y}}_2^2=0`$ and the values of $`R_1^1,P_{3i},`$ are taken respectively from (127) and (128).
By using the equation (130) we can define the N–coefficients (125), $`q_i(x^k,z),`$ if the function $`h_3(x^k,z)`$ and the components $`\mathrm{{\rm Y}}_{3i}`$ of the energy–momentum d–tensor are given. We note that the equations (128) are solved for arbitrary functions $`h_3=h_3(x^k)`$ and $`q_i=q_i(x^k,z)`$ if $`\mathrm{{\rm Y}}_{3i}=0`$ and in this case the component of d–metric $`h_3(x^k)`$ is not contained in the system of 3D field equations.
### A Static elliptic horizons
Let us consider a class of 3D d-metrics which local anisotropy which are similar to Banados–Teitelboim–Zanelli (BTZ) black holes .
The d–metric is parametrized
$$\delta s^2=g_1(\chi ^1,\chi ^2)(d\chi ^1)^2+\left(d\chi ^2\right)^2h_3(\chi ^1,\chi ^2,t)\left(\delta t\right)^2,$$
(131)
where $`\chi ^1=r/r_h`$ for $`r_h=const,`$ $`\chi ^2=\theta /r_a`$ if $`r_a=\sqrt{|\kappa \mathrm{{\rm Y}}_3^3|}0`$ and $`\chi ^2=\theta `$ if $`\mathrm{{\rm Y}}_3^3=0,`$ $`y^3=z=t,`$ where $`t`$ is the time like coordinate. The Einstein equations (129) and (130) transforms respectively into
$$\frac{^2g_1}{(\chi ^2)^2}\frac{1}{2g_1}\left(\frac{g_1}{\chi ^2}\right)^22\kappa \mathrm{{\rm Y}}_3^3g_1=0$$
(132)
and
$$\left[\frac{1}{h_3}\frac{^2h_3}{z^2}\left(\frac{1}{h_3}\frac{h_3}{z}\right)^2\right]q_i=\kappa \mathrm{{\rm Y}}_{3i}.$$
(133)
By introducing new variables
$$p=g_1^{}/g_1\text{ and }s=h_3^{}/h_3$$
(134)
where the ’prime’ in this subsection denotes the partial derivative $`/\chi ^2,`$ the equations (132) and (133) transform into
$$p^{}+\frac{p^2}{2}+2ϵ=0$$
(135)
and
$$s^{}q_i=\kappa \mathrm{{\rm Y}}_{3i},$$
(136)
where the vacuum case should be parametrized for $`ϵ=0`$ with $`\chi ^i=x^i`$ and $`ϵ=1(1)`$ for the signature $`1(1)`$ of the anisotropic coordinate.
A class of solutions of 3D Einstein equations for arbitrary $`q_i=q_i(\chi ^k,t)`$ and $`\mathrm{{\rm Y}}_{3i}=0`$ is obtained if $`s=s(\chi ^i).`$ After integration of the second equation from (134), we find
$$h_3(\chi ^k,t)=h_{3(0)}(\chi ^k)\mathrm{exp}\left[s_{(0)}\left(\chi ^k\right)t\right]$$
(137)
as a general solution of the system (136) with vanishing right part. Static solutions are stipulated by $`q_i=q_i(\chi ^k)`$ and $`s_{(0)}(\chi ^k)=0.`$
The integral curve of (135), intersecting a point $`(\chi _{(0)}^2,p_{(0)}),`$ considered as a differential equation on $`\chi ^2`$ is defined by the functions
$`p`$ $`=`$ $`{\displaystyle \frac{p_{(0)}}{1+\frac{p_{(0)}}{2}\left(\chi ^2\chi _{(0)}^2\right)}},ϵ=0;`$ (138)
$`p`$ $`=`$ $`{\displaystyle \frac{p_{(0)}2\mathrm{tanh}\left(\chi ^2\chi _{(0)}^2\right)}{1+\frac{p_{(0)}}{2}\mathrm{tanh}\left(\chi ^2\chi _{(0)}^2\right)}},ϵ>0;`$ (139)
$`p`$ $`=`$ $`{\displaystyle \frac{p_{(0)}2\mathrm{tan}\left(\chi ^2\chi _{(0)}^2\right)}{1+\frac{p_{(0)}}{2}\mathrm{tan}\left(\chi ^2\chi _{(0)}^2\right)}},ϵ<0.`$ (140)
Because the function $`p`$ depends also parametrically on variable $`\chi ^1`$ we must consider functions $`\chi _{(0)}^2=\chi _{(0)}^2\left(\chi ^1\right)`$ and $`p_{(0)}=p_{(0)}\left(\chi ^1\right).`$
For simplicity, here we elucidate the case $`ϵ<0.`$ The general formula for the non–trivial component of h–metric is to be obtained after integration on $`\chi ^1`$ of (140) (see formula (134))
$`g_1(\chi ^1,\chi ^2)`$ $`=g_{1(0)}\left(\chi ^1\right)\times `$ (142)
$`\left\{\mathrm{sin}[\chi ^2\chi _{(0)}^2\left(\chi ^1\right)]+\mathrm{arctan}{\displaystyle \frac{2}{p_{(0)}\left(\chi ^1\right)}}\right\}^2,`$
for $`p_{(0)}\left(\chi ^1\right)0,`$ and
$$g_1(\chi ^1,\chi ^2)=g_{1(0)}\left(\chi ^1\right)\mathrm{cos}^2[\chi ^2\chi _{(0)}^2\left(\chi ^1\right)]$$
(143)
for $`p_{(0)}(\chi ^1)=0,`$ where $`g_{1(0)}(\chi ^1),\chi _{(0)}^2(\chi ^1)`$ and $`p_{(0)}(\chi ^1)`$ are some functions of necessary smoothness class on variable $`\chi ^1=x^1/\sqrt{\kappa \epsilon },`$ when $`\epsilon `$ is the energy density. If we consider $`\mathrm{{\rm Y}}_{3i}=0`$ and a non–trivial diagonal components of energy–momentum d–tensor, $`\mathrm{{\rm Y}}_\beta ^\alpha =diag[0,0,\epsilon ],`$ the N–connection coefficients $`q_i(\chi ^i,t)`$ could be arbitrary functions.
For simplicity, in our further considerations we shall apply the solution (143).
The d–metric (131) with the coefficients (143) and (137) gives a general description of a class of solutions with generic local anisotropy of the Einstein equations (46).
Let us construct static black la–hole solutions for $`s_{(0)}\left(\chi ^k\right)=0`$ in (137).
In order to construct an explicit la–solution we choose some coefficients $`h_{3(0)}(\chi ^k),g_{1(0)}\left(\chi ^1\right)`$ and $`\chi _0\left(\chi ^1\right)`$ following some physical considerations. For instance, the Schwarzschild solution is selected from a general 4D metric with some general coefficients of static, spherical symmetry by relating the radial component of metric with the Newton gravitational potential. In this section, we construct a locally anisotropic BTZ like solution by supposing that it is conformally equivalent to the BTZ solution if one neglects anisotropies on angle $`\theta ),`$
$$g_{1(0)}\left(\chi ^1\right)=\left[r\left(M_0+\frac{r^2}{l^2}\right)\right]^2,$$
where $`M_0=const>0`$ and $`1/l^2`$ is a constant (which is to be considered the cosmological from the locally isotropic limit. The time–time coefficient of d–metric is chosen
$$h_3(\chi ^1,\chi ^2)=r^2\lambda _3(\chi ^1,\chi ^2)\mathrm{cos}^2[\chi ^2\chi _{(0)}^2\left(\chi ^1\right)].$$
(144)
If we chose in (144)
$$\lambda _3=(M_0+\frac{r^2}{l^2})^2,$$
when the constant
$$r_h=\sqrt{M_0}l$$
defines the radius of a circular horizon, the la–solution is conformally equivalent, with the factor $`r^2\mathrm{cos}^2[\chi ^2\chi _{(0)}^2\left(\chi ^1\right)],`$ to the BTZ solution embedded into a anholonomic background given by arbitrary functions $`q_i(\chi ^i,t)`$ which are defined by some initial conditions of gravitational la–background polarization.
A more general class of la–solutions could be generated if we put, for instance,
$$\lambda _3(\chi ^1,\chi ^2)=(M\left(\theta \right)+\frac{r^2}{l^2})^2,$$
with
$$M\left(\theta \right)=\frac{M_0}{(1+e\mathrm{cos}\theta )^2},$$
where $`e<1.`$ This solution has a horizon, $`\lambda _3=0,`$ parametrized by an ellipse
$$r=\frac{r_h}{1+e\mathrm{cos}\theta }$$
with parameter $`r_h`$ and eccentricity $`e.`$
We note that our solution with elliptic horizon was constructed for a diagonal energy–momentum d-tensor with non–trivial energy density but without cosmological constant. On the other hand the BTZ solution was constructed for a generic 3D cosmological constant. There is not a contradiction here because the la–solutions can be considered for a d–tensor $`\mathrm{{\rm Y}}_\beta ^\alpha =diag[p_11/l^2,p_21/l^2,\epsilon 1/l^2]`$ with $`p_{1,2}=1/l^2`$ and $`\epsilon _{(eff)}=\epsilon +1/l^2`$ (for $`\epsilon =const`$ the last expression defines the effective constant $`r_a).`$ The locally isotropic limit to the BTZ black hole could be realized after multiplication on $`r^2`$ and by approximations $`e0,`$ $`\mathrm{cos}[\theta \theta _0\left(\chi ^1\right)]1`$ and $`q_i(x^k,t)0.`$
### B Oscillating elliptic horizons
The simplest way to construct 3D solutions of the Einstein equations with oscillating in time horizon is to consider matter states with constant nonvanishing values of $`\mathrm{{\rm Y}}_{31}=const.`$ In this case the coefficient $`h_3`$ could depend on $`t`$–variable. For instance, we can chose such initial values when
$$h_3(\chi ^1,\theta ,t)=r^2\left(M\left(t\right)+\frac{r^2}{l^2}\right)\mathrm{cos}^2[\theta \theta _0\left(\chi ^1\right)]$$
(145)
with
$$M=M_0\mathrm{exp}\left(\stackrel{~}{p}t\right)\mathrm{sin}\stackrel{~}{\omega }t,$$
or, for an another type of anisotropy,
$`h_3(\chi ^1,\theta ,t)`$ $`=`$ $`r^2(M_0+{\displaystyle \frac{r^2}{l^2}})\times `$ (147)
$`\mathrm{cos}^2\theta \mathrm{sin}^2[\theta \theta _0(\chi ^1,t)]`$
with
$$\mathrm{cos}\theta _0(\chi ^1,t)=e^1\left(\frac{r_a}{r}\mathrm{cos}\omega _1t1\right),$$
when the horizon is given parametrically,
$$r=\frac{r_a}{1+e\mathrm{cos}\theta }\mathrm{cos}\omega _1t,$$
where the new constants (comparing with those from the previous subsection) are fixed by some initial and boundary conditions as to be $`\stackrel{~}{p}>0,`$ and $`\stackrel{~}{\omega }`$ and $`\omega _1`$ are treated as some real numbers.
For a prescribed value of $`h_3(\chi ^1,\theta ,t)`$ with non–zero source $`\mathrm{{\rm Y}}_{31},`$ in the equation (130), we obtain
$$q_1(\chi ^1,\theta ,t)=\kappa \mathrm{{\rm Y}}_{31}\left(\frac{^2}{t^2}\mathrm{ln}|h_3(\chi ^1,\theta ,t)|\right)^1.$$
(148)
A solution (124) of the Einstein equations (129) and (130) with $`g_2(\chi ^i)=1`$ and $`g_1(\chi ^1,\theta )`$ and $`h_3(\chi ^1,\theta ,t)`$ given respectively by formulas (143) and (145) describe a 3D evaporating black la–hole solution with circular oscillating in time horizon. An another type of solution, with elliptic oscillating in time horizon, could be obtained if we choose (147). The non–trivial coefficient of the N–connection must be computed following the formula (148).
## VIII 4D Locally Anisotropic Black Holes
### A Basic properties
The purpose of this Section is the construction of d–metrics of Class 2, or 3 (see (68) and (69), or (71) and (72)) which are conformally equivalent to some la–deformations of black hole, disk, torus and cylinder like solutions. We shall analyze 4D d-metrics of type
$`\delta s^2`$ $`=`$ $`g_1\left(x^k\right)\left(dx^1\right)^2+\left(dx^2\right)^2`$ (150)
$`+h_3(x^i,z)\left(\delta z\right)^2+h_4(x^i,z)\left(\delta y^4\right)^2.`$
The Einstein equations (61) with the Ricci h–tensor (53) and energy momentum d–tensor (70), or (73), transforms into
$$\frac{^2g_1}{(x^1)^2}\frac{1}{2g_1}\left(\frac{g_1}{x^1}\right)^22\kappa \mathrm{{\rm Y}}_3^3g_1=0.$$
(151)
By introducing the coordinates $`\chi ^i=x^i/\sqrt{|\kappa \mathrm{{\rm Y}}_3^3|}`$ for the Class 3 (2) d–metrics and the variable $`p=g_1^{}/g_1,`$ where by ’prime’ in this Section is considered the partial derivative $`/\chi ^2,`$ the equation (151) transforms into
$$p^{}+\frac{p^2}{2}+2ϵ=0,$$
(152)
where the vacuum case should be parametrized for $`ϵ=0`$ with $`\chi ^i=x^i`$ and $`ϵ=1(1)`$ for Class 2 (3) d–metrics. The equations (151) and (152) are, correspondingly, equivalent to the equations (132) and (135) with that difference that in this Section we are dealing with 4D coefficients and values. The solutions for the h–metric are parametrized like (138), (139), and (140) and the coefficient $`g_1(\chi ^i)`$ is given by a similar to (143) formula (for simplicity, here we elucidate the case $`ϵ<0)`$ which for $`p_{(0)}\left(\chi ^1\right)=0`$ transforms into
$$g_1(\chi ^1,\chi ^2)=g_{1(0)}\left(\chi ^1\right)\mathrm{cos}^2[\chi ^2\chi _{(0)}^2\left(\chi ^1\right)],$$
(153)
where $`g_1\left(\chi ^1\right),\chi _{(0)}^2\left(\chi ^1\right)`$ and $`p_{(0)}\left(\chi ^1\right)`$ are some functions of necessary smoothness class on variable $`\chi ^1=x^1/\sqrt{\kappa \epsilon },`$ $`\epsilon `$ is the energy density. The coefficients $`g_1(\chi ^1,\chi ^2)`$ (153) and $`g_2(\chi ^1,\chi ^2)=1`$ define a h–metric of Class 3 (71) with energy–momentum d–tensor (73). The next step is the construction of h–components of d–metrics for different classes of symmetries of anisotropies.
Now, let us consider the system of equations (62) with the vertical Ricci d–tensor component (54) which are satisfied by arbitrary functions
$$h_3=a_3(\chi ^i)\text{ and }h_4=a_4(\chi ^i).$$
(154)
If v–metrics depending on three coordinates are introduced, $`h_a=h_a(\chi ^i,z),`$ the v–components of the Einstein equations transforms into (83) which reduces to (88) for prescribed values of $`h_3(\chi ^i,z),`$ and, inversely, to (90) if $`h_4(\chi ^i,z)`$ is prescribed. For h–metrics being conformally equivalent to (154) (see transforms (91)) we are dealing to equations of type (92) with integral varieties (93).
### B Rotation Hypersurfaces Horizons
We proof that there are static black hole and cylindrical like solutions of the Einstein equations with horizons being some 3D rotation hypersurfaces. The space components of corresponding d–metrics are conformally equivalent to some locally anisotropic deformations of the spherical symmetric Schwarzschild and cylindrical Weyl solutions. We note that for some classes of solutions the local anisotropy is contained in non–perturbative anholonomic structures.
#### 1 Rotation ellipsoid configuration
There are two types of rotation ellipsoids, elongated and flattened ones. We examine both cases of such horizon configurations
##### a Elongated rotation ellipsoid coordinates:
An elongated rotation ellipsoid hypersurface is given by the formula
$$\frac{\stackrel{~}{x}^2+\stackrel{~}{y}^2}{\sigma ^21}+\frac{\stackrel{~}{z}^2}{\sigma ^2}=\stackrel{~}{\rho }^2,$$
(155)
where $`\sigma 1`$ and $`\stackrel{~}{\rho }`$ is similar to the radial coordinate in the spherical symmetric case.
The space 3D coordinate system is defined
$`\stackrel{~}{x}`$ $`=`$ $`\stackrel{~}{\rho }\mathrm{sinh}u\mathrm{sin}v\mathrm{cos}\phi ,\stackrel{~}{y}=\stackrel{~}{\rho }\mathrm{sinh}u\mathrm{sin}v\mathrm{sin}\phi ,`$ (156)
$`\stackrel{~}{z}`$ $`=`$ $`\stackrel{~}{\rho }\mathrm{cosh}u\mathrm{cos}v,`$ (157)
where $`\sigma =\mathrm{cosh}u,(0u<\mathrm{},0v\pi ,0\phi <2\pi ).`$ The hypersurface metric is
$`g_{uu}`$ $`=`$ $`g_{vv}=\stackrel{~}{\rho }^2\left(\mathrm{sinh}^2u+\mathrm{sin}^2v\right),`$ (158)
$`g_{\phi \phi }`$ $`=`$ $`\stackrel{~}{\rho }^2\mathrm{sinh}^2u\mathrm{sin}^2v.`$ (159)
Let us introduce a d–metric
$`\delta s^2`$ $`=`$ $`g_1(u,v)du^2+dv^2+`$ (161)
$`h_3(u,v,\phi )\left(\delta t\right)^2+h_4(u,v,\phi )\left(\delta \phi \right)^2,`$
where $`\delta t`$ and $`\delta \phi `$ are N–elongated differentials.
As a particular solution (153) for the h–metric we choose the coefficient
$$g_1(u,v)=\mathrm{cos}^2v.$$
(162)
The $`h_3(u,v,\phi )=h_3(u,v,\stackrel{~}{\rho }(u,v,\phi ))`$ is considered as
$$h_3(u,v,\stackrel{~}{\rho })=\frac{1}{\mathrm{sinh}^2u+\mathrm{sin}^2v}\frac{\left[1\frac{r_g}{4\stackrel{~}{\rho }}\right]^2}{\left[1+\frac{r_g}{4\stackrel{~}{\rho }}\right]^6}.$$
(163)
In order to define the $`h_4`$ coefficient solving the Einstein equations, for simplicity, with a diagonal energy–momentum d–tensor for vanishing pressure, we must solve the equation (88) which transforms into a linear equation (89) if $`\mathrm{{\rm Y}}_1=0.`$ In our case $`s(u,v,\phi )=\beta ^1(u,v,\phi ),`$ where $`\beta =\left(h_4/\phi \right)/h_4,`$ must be a solution of
$$\frac{s}{\phi }+\frac{\mathrm{ln}\sqrt{\left|h_3\right|}}{\phi }s=\frac{1}{2}.$$
After two integrations (see ) the general solution for $`h_4(u,v,\phi ),`$ is
$$h_4(u,v,\phi )=a_4(u,v)\mathrm{exp}\left[\underset{0}{\overset{\phi }{}}F(u,v,z)𝑑z\right],$$
(164)
where
$`F(u,v,z)`$ $`=`$ $`(\sqrt{|h_3(u,v,z)|}[s_{1(0)}(u,v)`$ (166)
$`+{\displaystyle \frac{1}{2}}{\displaystyle \underset{z_0(u,v)}{\overset{z}{}}}\sqrt{|h_3(u,v,z)|}dz])^1,`$
$`s_{1(0)}(u,v)`$ and $`z_0(u,v)`$ are some functions of necessary smooth class. We note that if we put $`h_4=a_4(u,v)`$ the equations (62) are satisfied for every $`h_3=h_3(u,v,\phi ).`$
Every d–metric (161) with coefficients of type (162), (163) and (164) solves the Einstein equations (61)–(64) with the diagonal momentum d–tensor
$$\mathrm{{\rm Y}}_\beta ^\alpha =diag[0,0,\epsilon =m_0,0],$$
when $`r_g=2\kappa m_0;`$ we set the light constant $`c=1.`$ If we choose
$$a_4(u,v)=\frac{\mathrm{sinh}^2u\mathrm{sin}^2v}{\mathrm{sinh}^2u+\mathrm{sin}^2v}$$
our solution is conformally equivalent (if not considering the time–time component) to the hypersurface metric (158). The condition of vanishing of the coefficient (163) parametrizes the rotation ellipsoid for the horizon
$$\frac{\stackrel{~}{x}^2+\stackrel{~}{y}^2}{\sigma ^21}+\frac{\stackrel{~}{z}^2}{\sigma ^2}=\left(\frac{r_g}{4}\right)^2,$$
where the radial coordinate is redefined via relation $`\stackrel{~}{r}=\stackrel{~}{\rho }\left(1+\frac{r_g}{4\stackrel{~}{\rho }}\right)^2.`$ After multiplication on the conformal factor
$$\left(\mathrm{sinh}^2u+\mathrm{sin}^2v\right)\left[1+\frac{r_g}{4\stackrel{~}{\rho }}\right]^4,$$
approximating $`g_1(u,v)=\mathrm{sin}^2v0,`$ in the limit of locally isotropic spherical symmetry,
$$\stackrel{~}{x}^2+\stackrel{~}{y}^2+\stackrel{~}{z}^2=r_g^2,$$
the d–metric (161) reduces to
$$ds^2=\left[1+\frac{r_g}{4\stackrel{~}{\rho }}\right]^4\left(d\stackrel{~}{x}^2+d\stackrel{~}{y}^2+d\stackrel{~}{z}^2\right)\frac{\left[1\frac{r_g}{4\stackrel{~}{\rho }}\right]^2}{\left[1+\frac{r_g}{4\stackrel{~}{\rho }}\right]^2}dt^2$$
which is just the Schwazschild solution with the redefined radial coordinate when the space component becomes conformally Euclidean.
So, the d–metric (161), the coefficients of N–connection being solutions of (63) and (64), describe a static 4D solution of the Einstein equations when instead of a spherical symmetric horizon one considers a locally anisotropic deformation to the hypersurface of rotation elongated ellipsoid.
##### b Flattened rotation ellipsoid coordinates
In a similar fashion we can construct a static 4D black hole solution with the horizon parametrized by a flattened rotation ellipsoid ,
$$\frac{\stackrel{~}{x}^2+\stackrel{~}{y}^2}{1+\sigma ^2}+\frac{\stackrel{~}{z}^2}{\sigma ^2}=\stackrel{~}{\rho }^2,$$
where $`\sigma 0`$ and $`\sigma =\mathrm{sinh}u.`$
The space 3D special coordinate system is defined
$`\stackrel{~}{x}`$ $`=`$ $`\stackrel{~}{\rho }\mathrm{cosh}u\mathrm{sin}v\mathrm{cos}\phi ,\stackrel{~}{y}=\stackrel{~}{\rho }\mathrm{cosh}u\mathrm{sin}v\mathrm{sin}\phi ,`$ (167)
$`\stackrel{~}{z}`$ $`=`$ $`\stackrel{~}{\rho }\mathrm{sinh}u\mathrm{cos}v,`$ (168)
where $`0u<\mathrm{},0v\pi ,0\phi <2\pi .`$
The hypersurface metric is
$`g_{uu}`$ $`=`$ $`g_{vv}=\stackrel{~}{\rho }^2\left(\mathrm{sinh}^2u+\mathrm{cos}^2v\right),`$ (169)
$`g_{\phi \phi }`$ $`=`$ $`\stackrel{~}{\rho }^2\mathrm{sinh}^2u\mathrm{cos}^2v.`$ (170)
In the rest the black hole solution is described by the same formulas as in the previous subsection but with respect to new canonical coordinates for flattened rotation ellipsoid.
#### 2 Cylindrical, Bipolar and Toroidal Configurations
We consider a d–metric of type (150). As a coefficient for h–metric we choose $`g_1(\chi ^1,\chi ^2)=\left(\mathrm{cos}\chi ^2\right)^2`$ which solves the Einstein equations (61). The energy momentum d–tensor is chosen to be diagonal, $`\mathrm{{\rm Y}}_\beta ^\alpha =diag[0,0,\epsilon ,0]`$ with $`\epsilon m_0=m_{(lin)}𝑑l,`$ where $`\epsilon _{(lin)}`$ is the linear ’mass’ density. The coefficient $`h_3(\chi ^i,z)`$ will be chosen in a form similar to (163),
$$h_3\left[1\frac{r_g}{4\stackrel{~}{\rho }}\right]^2/\left[1+\frac{r_g}{4\stackrel{~}{\rho }}\right]^6$$
for a cylindrical elliptic horizon. We parametrize the second v–component as $`h_4=a_4(\chi ^1,\chi ^2)`$ when the equations (62) are satisfied for every $`h_3=h_3(\chi ^1,\chi ^2,z).`$
We note that in this work we only proof the existence of the mentioned type horizon configurations. The exact solutions and physics of so–called ellipsoidal black holes, black torus, black cylinders and black disks with, or not, local anisotropy will be examined in .
##### a Cylindrical coordinates:
Let us construct a solution of the Einstein equation with the horizon having the symmetry of ellipsoidal cylinder given by hypersurface formula
$$\frac{\stackrel{~}{x}^2}{\sigma ^2}+\frac{\stackrel{~}{y}^2}{\sigma ^21}=\rho _{}^2,\stackrel{~}{z}=\stackrel{~}{z},$$
where $`\sigma 1.`$ The 3D radial coordinate $`\stackrel{~}{r}`$ is to be computed from $`\stackrel{~}{\rho }^2=\rho _{}^2+\stackrel{~}{z}^2.`$
The 3D space coordinate system is defined
$$\stackrel{~}{x}=\rho _{}\mathrm{cosh}u\mathrm{cos}v,\stackrel{~}{y}=\rho _{}\mathrm{sinh}u\mathrm{sin}v\mathrm{sin},\stackrel{~}{z}=\stackrel{~}{z},$$
where $`\sigma =\mathrm{cosh}u,(0u<\mathrm{},0v\pi ).`$
The hypersurface metric is
$$g_{uu}=g_{vv}=\rho _{}^2\left(\mathrm{sinh}^2u+\mathrm{sin}^2v\right),g_{zz}=1.$$
(171)
A solution of the Einstein equations with singularity on an ellipse is given by
$`h_3`$ $`=`$ $`{\displaystyle \frac{1}{\rho _{}^2\left(\mathrm{sinh}^2u+\mathrm{sin}^2v\right)}}\times {\displaystyle \frac{\left[1\frac{r_g}{4\stackrel{~}{\rho }}\right]^2}{\left[1+\frac{r_g}{4\stackrel{~}{\rho }}\right]^6}},`$ (172)
$`h_4`$ $`=`$ $`a_4={\displaystyle \frac{1}{\rho _{}^2\left(\mathrm{sinh}^2u+\mathrm{sin}^2v\right)}},`$ (173)
where $`\stackrel{~}{r}=\stackrel{~}{\rho }\left(1+\frac{r_g}{4\stackrel{~}{\rho }}\right)^2.`$ The condition of vanishing of the time–time coefficient $`h_3`$ parametrizes the hypersurface equation of the horizon
$$\frac{\stackrel{~}{x}^2}{\sigma ^2}+\frac{\stackrel{~}{y}^2}{\sigma ^21}=\left(\frac{\rho _{(g)}}{4}\right)^2,\stackrel{~}{z}=\stackrel{~}{z},$$
where $`\rho _{(g)}=2\kappa m_{(lin)}.`$
By multiplying the d–metric on the conformal factor
$$\rho _{}^2\left(\mathrm{sinh}^2u+\mathrm{sin}^2v\right)\left[1+\frac{r_g}{4\stackrel{~}{\rho }}\right]^4,$$
where $`r_g=\rho _{(g)}𝑑l`$ (the integration is taken along the ellipse), for $`\rho _{}1,`$ in the local isotropic limit, $`\mathrm{sin}v0,`$ the space component transforms into (171).
##### b Bipolar coordinates:
Let us construct 4D solutions of the Einstein equation with the horizon having the symmetry of the bipolar hypersurface given by the formula
$$\left(\sqrt{\stackrel{~}{x}^2+\stackrel{~}{y}^2}\frac{\stackrel{~}{\rho }}{\mathrm{tan}\sigma }\right)^2+\stackrel{~}{z}^2=\frac{\stackrel{~}{\rho }^2}{\mathrm{sin}^2\sigma },$$
which describes a hypersurface obtained under the rotation of the circles
$$\left(\stackrel{~}{y}\frac{\stackrel{~}{\rho }}{\mathrm{tan}\sigma }\right)^2+\stackrel{~}{z}^2=\frac{\stackrel{~}{\rho }^2}{\mathrm{sin}^2\sigma }$$
around the axes $`Oz`$; because $`|c\mathrm{tan}\sigma |<|\mathrm{sin}\sigma |^1,`$ the circles intersect the axes $`Oz.`$ The 3D space coordinate system is defined
$`\stackrel{~}{x}`$ $`=`$ $`{\displaystyle \frac{\stackrel{~}{\rho }\mathrm{sin}\sigma \mathrm{cos}\phi }{\mathrm{cosh}\tau \mathrm{cos}\sigma }},\stackrel{~}{y}={\displaystyle \frac{\stackrel{~}{\rho }\mathrm{sin}\sigma \mathrm{sin}\phi }{\mathrm{cosh}\tau \mathrm{cos}\sigma }},`$ (174)
$`\stackrel{~}{z}`$ $`=`$ $`{\displaystyle \frac{\stackrel{~}{r}\mathrm{sinh}\tau }{\mathrm{cosh}\tau \mathrm{cos}\sigma }},`$ (176)
$`\left(\mathrm{}<\tau <\mathrm{},0\sigma <\pi ,0\phi <2\pi \right).`$
The hypersurface metric is
$$g_{\tau \tau }=g_{\sigma \sigma }=\frac{\stackrel{~}{\rho }^2}{\left(\mathrm{cosh}\tau \mathrm{cos}\sigma \right)^2},g_{\phi \phi }=\frac{\stackrel{~}{\rho }^2\mathrm{sin}^2\sigma }{\left(\mathrm{cosh}\tau \mathrm{cos}\sigma \right)^2}.$$
(177)
A solution of the Einstein equations with singularity on a circle is given by
$$h_3=\left[1\frac{r_g}{4\stackrel{~}{\rho }}\right]^2/\left[1+\frac{r_g}{4\stackrel{~}{\rho }}\right]^6\text{ and }h_4=a_4=\mathrm{sin}^2\sigma ,$$
where $`\stackrel{~}{r}=\stackrel{~}{\rho }\left(1+\frac{r_g}{4\stackrel{~}{\rho }}\right)^2.`$ The condition of vanishing of the time–time coefficient $`h_3`$ parametrizes the hypersurface equation of the horizon
$$\left(\sqrt{\stackrel{~}{x}^2+\stackrel{~}{y}^2}\frac{r_g}{2}c\mathrm{tan}\sigma \right)^2+\stackrel{~}{z}^2=\frac{r_g^2}{4\mathrm{sin}^2\sigma },$$
where $`r_g=\rho _{(g)}𝑑l`$ (the integration is taken along the circle), $`\rho _{(g)}=2\kappa m_{(lin)}.`$
By multiplying the d–metric on the conformal factor
$$\frac{1}{\left(\mathrm{cosh}\tau \mathrm{cos}\sigma \right)^2}\left[1+\frac{r_g}{4\stackrel{~}{\rho }}\right]^4,$$
(178)
for $`\rho _{}1,`$ in the local isotropic limit, $`\mathrm{sin}v0,`$ the space component transforms into (177).
##### c Toroidal coordinates:
Let us consider solutions of the Einstein equations with toroidal symmetry of horizons. The hypersurface formula of a torus is
$$\left(\sqrt{\stackrel{~}{x}^2+\stackrel{~}{y}^2}\stackrel{~}{\rho }c\mathrm{tanh}\sigma \right)^2+\stackrel{~}{z}^2=\frac{\stackrel{~}{\rho }^2}{\mathrm{sinh}^2\sigma }.$$
The 3D space coordinate system is defined
$`\stackrel{~}{x}`$ $`=`$ $`{\displaystyle \frac{\stackrel{~}{\rho }\mathrm{sinh}\tau \mathrm{cos}\phi }{\mathrm{cosh}\tau \mathrm{cos}\sigma }},\stackrel{~}{y}={\displaystyle \frac{\stackrel{~}{\rho }\mathrm{sin}\sigma \mathrm{sin}\phi }{\mathrm{cosh}\tau \mathrm{cos}\sigma }},`$ (179)
$`\stackrel{~}{z}`$ $`=`$ $`{\displaystyle \frac{\stackrel{~}{\rho }\mathrm{sinh}\sigma }{\mathrm{cosh}\tau \mathrm{cos}\sigma }},`$ (181)
$`\left(\pi <\sigma <\pi ,0\tau <\mathrm{},0\phi <2\pi \right).`$
The hypersurface metric is
$$g_{\sigma \sigma }=g_{\tau \tau }=\frac{\stackrel{~}{\rho }^2}{\left(\mathrm{cosh}\tau \mathrm{cos}\sigma \right)^2},g_{\phi \phi }=\frac{\stackrel{~}{\rho }^2\mathrm{sin}^2\sigma }{\left(\mathrm{cosh}\tau \mathrm{cos}\sigma \right)^2}.$$
(182)
This, another type of solution of the Einstein equations with singularity on a circle, is given by
$$h_3=\left[1\frac{r_g}{4\stackrel{~}{\rho }}\right]^2/\left[1+\frac{r_g}{4\stackrel{~}{\rho }}\right]^6\text{ and }h_4=a_4=\mathrm{sinh}^2\sigma ,$$
where $`\stackrel{~}{r}=\stackrel{~}{\rho }\left(1+\frac{r_g}{4\stackrel{~}{\rho }}\right)^2.`$ The condition of vanishing of the time–time coefficient $`h_3`$ parametrizes the hypersurface equation of the horizon
$$\left(\sqrt{\stackrel{~}{x}^2+\stackrel{~}{y}^2}\frac{r_g}{2\mathrm{tanh}\sigma }c\right)^2+\stackrel{~}{z}^2=\frac{r_g^2}{4\mathrm{sinh}^2\sigma },$$
where $`r_g=\rho _{(g)}𝑑l`$ (the integration is taken along the circle), $`\rho _{(g)}=2\kappa m_{(lin)}.`$
By multiplying the d–metric on the conformal factor (178), for $`\rho _{}1,`$ in the local isotropic limit, $`\mathrm{sin}v0,`$ the space component transforms into (182).
### C Disks with Local Anisotropy
The d–metric is of type (161)
$`\delta s^2`$ $`=`$ $`g_1(\rho ,\zeta )d\rho ^2+d\zeta ^2+`$ (184)
$`h_3(\rho ,\zeta ,\phi )\left(\delta t\right)^2+h_4(\rho ,\zeta ,\phi )\left(\delta \phi ^{}\right)^2,`$
where the 4D coordinates are parametrized $`x^1=\rho `$ (the coordinate radius), $`x^2=\zeta ,y^3=t,y^4=\stackrel{~}{\phi }`$ (like for the disk solution in general relativity ) and $`\delta t`$ and $`\delta \stackrel{~}{\phi }`$ are N–elongated differentials). One uses also primed coordinates given with respect to corotating frame of reference, $`\rho ^{}=\rho ,\zeta ^{}=\zeta ,\phi ^{}=\phi \mathrm{\Omega }t,t^{}=t,`$ where $`\mathrm{\Omega }`$ is the angular velocity as measured by an observer at $`\mathrm{}.`$ The h–coordinates run respectively values $`0\rho <\mathrm{}`$ and $`\mathrm{}<\zeta <\mathrm{}.`$ We consider a disk defined by conditions $`\zeta =0`$ and $`\rho \rho _0.`$
As a particular solution (153) for the h–metric we choose the coefficient
$$g_1(\rho ,\zeta )=\left(\mathrm{cos}\zeta \right)^2.$$
(185)
The explicit form of coefficients $`h_3`$ and $`h_4`$ are defined by using functions
$$A=\rho ^2\mathrm{exp}[2\left(Uk\right)],B=\mathrm{exp}\left(4U\right)$$
(186)
and
$$\stackrel{~}{\phi }=\phi \frac{Ba}{ABa^2}t$$
where $`U,k,`$ and $`a`$ are some functions on $`(\rho ,\zeta ,\phi ).`$ For the locally isotropic disk solutions we consider only dependencies on $`(\rho ,\zeta ),`$ in this case we shall write
$$U=U_0(\rho ,\zeta ),k=k_0(\rho ,\zeta ),a=a_0(\rho ,\zeta )$$
and
$$A=A_0(\rho ,\zeta ),B=B_0(\rho ,\zeta ),$$
where the values with the index $`0`$ are computed by using the Neugebauer and Meinel disk solution . The (energy) mass density is taken
$$\epsilon (\rho ,\zeta ,\phi )=\delta \left(\zeta \right)\mathrm{exp}\left(Uk\right)\sigma _p(\rho ,\phi ),$$
where $`\sigma _p(\rho ,\phi )`$ is the (proper) surface mass density which is (in the la–case) non–uniformly distributed on the disk; for locally isotropic distributions $`\sigma _p=\sigma _p\left(\rho \right).`$ After the problem is solved one computes $`\sigma _p`$ as
$$\sigma _p=\frac{1}{2\pi }e^{Uk}\frac{U^{}}{\zeta }_{\zeta =0^+}.$$
The time–time component $`h_3`$ is chosen in the form
$$h_3(\rho ,\zeta ,\phi )=\frac{AB}{ABa^2}$$
and, for simplicity, we state the second v–component of d–metric $`h_4`$ to depend only h–coordinates as
$$h_4=a_4(\rho ,\zeta )=A_0B_0a_0^2.$$
As in the locally isotropic case one introduces the complex Ernst potential
$$f(\rho ,\zeta ,\phi )=e^{2U(\rho ,\zeta ,\phi )}+ib(\rho ,\zeta ,\phi ),$$
which depends additionally on coordinate $`\phi .`$ If the real and imaginary part of this potential are defined the coefficients (186) are computed
$`a(\rho ,\zeta ,\phi )`$ $`=`$ $`{\displaystyle \underset{0}{\overset{\rho }{}}}\stackrel{~}{\rho }e^{4U}b,_\zeta d\stackrel{~}{\rho },`$ (187)
$`k(\rho ,\zeta ,\phi )`$ $`=`$ $`{\displaystyle \underset{0}{\overset{\rho }{}}}\stackrel{~}{\rho }[U,_{\stackrel{~}{\rho }}^2U,_\zeta ^2+{\displaystyle \frac{1}{4}}e^{4U}(b,_{\stackrel{~}{\rho }}^2b,_\zeta ^2)]d\stackrel{~}{\rho }.`$ (188)
\[In the integrands, one has $`U=U(\stackrel{~}{\rho },\zeta ,\phi )`$ and $`b=b(\stackrel{~}{\rho },\zeta ,\phi )`$.\]
The Ernst potential is computed as in the locally isotropic limit with that difference that the values also depend on angular parameter $`\phi ,`$
$$f=\mathrm{exp}\left\{\underset{K_1}{\overset{K_a}{}}\frac{K^2dK}{Z}+\underset{K_2}{\overset{K_b}{}}\frac{K^2dK}{Z}v_2\right\},$$
(189)
with
$`Z`$ $`=`$ $`\sqrt{(K+iz)(Ki\overline{z})(K^2K_1^2)(K^2K_2^2)},`$ (190)
$`K_1`$ $`=`$ $`\rho _0\sqrt{{\displaystyle \frac{i\mu }{\mu }}}(\mathrm{}K_1<0),K_2=\overline{K}_1,`$ (191)
where $`\mathrm{}`$ denotes the real part. The real (positive) parameter $`\mu `$ is given by
$$\mu =2\mathrm{\Omega }^2\rho _0^2e^{2V_0}$$
where $`V_0=const.`$ The upper integration limits $`K_a`$ and $`K_b`$ in (189) are calculated from
$$\underset{K_1}{\overset{K_a}{}}\frac{dK}{Z}+\underset{K_2}{\overset{K_b}{}}\frac{dK}{Z}=v_0,\underset{K_1}{\overset{K_a}{}}\frac{KdK}{Z}+\underset{K_2}{\overset{K_b}{}}\frac{KdK}{Z}=v_1,$$
(192)
where the functions $`v_0`$, $`v_1`$ and $`v_2`$ in (192) and (189) are given by
$`v_0`$ $`=`$ $`{\displaystyle \underset{i\rho _0}{\overset{+i\rho _0}{}}}{\displaystyle \frac{H}{Z_1}}𝑑K,v_1={\displaystyle \underset{i\rho _0}{\overset{+i\rho _0}{}}}{\displaystyle \frac{H}{Z_1}}K𝑑K,`$ (193)
$`v_2`$ $`=`$ $`{\displaystyle \underset{i\rho _0}{\overset{+i\rho _0}{}}}{\displaystyle \frac{H}{Z_1}}K^2𝑑K,`$ (194)
$`H`$ $`=`$ $`{\displaystyle \frac{\mu \mathrm{ln}\left[\sqrt{1+\mu ^2(1+K^2/\rho _0^2)^2}+\mu (1+K^2/\rho _0^2)\right]}{\pi i\rho _0^2\sqrt{1+\mu ^2(1+K^2/\rho _0^2)^2}}}`$ (195)
$`(\mathrm{}H`$ $`=`$ $`0),`$ (196)
$$Z_1=\sqrt{(K+iz)(Ki\overline{z})},$$
where $`\mathrm{}Z_1<0`$ for $`\rho `$ and $`\zeta `$ outside the disk. In (193) one has to integrate along the imaginary axis. The integrations from $`K_1`$ to $`K_a`$ and $`K_2`$ to $`K_b`$ in (189) and (192) have to be performed along the same paths in the two–seethed Riemann surface associated with $`Z(K)`$. The problem of finding $`K_a`$ and $`K_b`$ from (192) is a special case of Jacobi’s inversion problem.
So, we have constructed a locally anisotropic generalization of the Neugebauer–Meinel disk solution in general relativity, with an additional dependence on angle $`\phi `$. In the locally isotropic limit, $`g_1=\left(\mathrm{cos}\zeta \right)^21,`$ when the d–metric (184) is conformally equivalent (with the factor $`\mathrm{exp}[2(U_0(\rho ,\zeta )k_0(\rho ,\zeta ))]`$) to the disk solution from ).
### D Locally Anisotropic generalizations of the Schwarzschild and Kerr solutions
#### 1 A Schwarzschild like la–solution
The d–metric of type (161) is taken
$`\delta s^2`$ $`=`$ $`g_1(\chi ^1,\theta )d(\chi ^1)^2+d\theta ^2+`$ (198)
$`h_3(\chi ^1,\theta ,\phi )\left(\delta t\right)^2+h_4(\chi ^1,\theta ,\phi )\left(\delta \phi \right)^2,`$
where on the horizontal subspace $`\chi ^1=\rho /r_a`$ is the undimensional radial coordinate (the constant $`r_a`$ will be defined below), $`\chi ^2=\theta `$ and in the vertical subspace $`y^3=z=t`$ and $`y^4=\phi .`$ The energy–momentum d–tensor is taken to be diagonal $`\mathrm{{\rm Y}}_\beta ^\alpha =diag[0,0,\epsilon ,0].`$ The coefficient $`g_1`$ is chosen to be a solution of type (153)
$$g_1(\chi ^1,\theta )=\mathrm{cos}^2\theta .$$
For $`h_4=\mathrm{sin}^2\theta `$ and $`h_3(\rho )=[1r_a/4\rho ]^2/[1+r_a/4\rho ]^6,`$ where $`r=`$ $`\rho (1+\frac{r_g}{4\rho })^2,`$ $`r^2=x^2+y^2+z^2,`$ $`r_a\dot{=}r_g`$ is the Schwarzschild gravitational radius, the d–metric (198) describes a la–solution of the Einstein equations which is conformally equivalent, with the factor $`\rho ^2\left(1+r_g/4\rho \right)^2,`$ to the Schwarzschild solution (written in coordinates $`(\rho ,\theta ,\phi ,t)),`$ embedded into a la–background given by non–trivial values of $`q_i(\rho ,\theta ,t)`$ and $`n_i(\rho ,\theta ,t).`$ In the anisotropic case we can extend the solution for anisotropic (on angle $`\theta )`$ gravitational polarizations of point particles masses, $`m=m\left(\theta \right),`$ for instance in elliptic form, when
$$r_a\left(\theta \right)=\frac{r_g}{\left(1+e\mathrm{cos}\theta \right)}$$
induces an ellipsoidal dependence on $`\theta `$ of the radial coordinate,
$$\rho =\frac{r_g}{4\left(1+e\mathrm{cos}\theta \right)}.$$
We can also consider arbitrary solutions with $`r_a=r_a(\theta ,t)`$ of oscillation type, $`r_a\mathrm{sin}\left(\omega _1t\right),`$ or modelling the mass evaporation, $`r_a\mathrm{exp}[st],s=const>0.`$
So, fixing a physical solution for $`h_3(\rho ,\theta ,t),`$ for instance,
$$h_3(\rho ,\theta ,t)=\frac{\left[1r_a\mathrm{exp}[st]/4\rho \left(1+e\mathrm{cos}\theta \right)\right]^2}{\left[1+r_a\mathrm{exp}[st]/4\rho \left(1+e\mathrm{cos}\theta \right)\right]^6},$$
where $`e=const<1,`$ and computing the values of $`q_i(\rho ,\theta ,t)`$ and $`n_i(\rho ,\theta ,t)`$ from (63) and (64), corresponding to given $`h_3`$ and $`h_4,`$ we obtain a la–generalization of the Schwarzschild metric.
We note that fixing this type of anisotropy, in the locally isotropic limit we obtain not just the Schwarzschild metric but a conformally transformed one, multiplied on the factor $`1/\rho ^2\left(1+r_g/4\rho \right)^4.`$
#### 2 A Kerr like la–solution
The d–metric is of type (184) is taken
$`\delta s^2`$ $`=`$ $`g_1(r/r_g,\theta )dr^2+d\theta ^2+`$ (200)
$`h_3(r/r_g,\theta ,\stackrel{~}{\phi })\left(\delta t\right)^2+h_4(r/r_g,\theta ,\stackrel{~}{\phi })\left(\delta \stackrel{~}{\phi }\right)^2.`$
In the locally isotropic limit this metric is conformally equivalent to the Kerr solution, with the factor $`r_{}^2=r^2+a_{}^2,`$ $`a_{}=const`$ is associated to the rotation momentum, if
$$g_1^{[i]}=1/\left(r\right),$$
where $`(r)=r^2rr_g+a_{}^2,`$ $`r_g`$ is the gravitational radius and the index $`[i]`$ points to locally isotropic values,
$$h_4^{[i]}=A(r,\theta ),$$
where $`A(r,\theta )=\frac{\mathrm{sin}^2\theta }{r_{}^2}\left(r^2+a_{}^2+\frac{rr_ga_{}^2}{r_{}^2}\mathrm{sin}^2\theta \right)`$ and
$$h_3^{[i]}=\left(\frac{Q^2}{A}+B\right),$$
where $`Q=\frac{rr_ga_{}}{r_{}^4}\mathrm{sin}^2\theta `$ and $`B=\frac{1}{r_{}^2}\left(1\frac{rr_g}{r_{}^2}\right).`$ The tilded angular variable $`\stackrel{~}{\phi }`$ is introduced with the aim to get a diagonal d–metric, $`\stackrel{~}{\phi }=\phi \frac{Q}{A}t.`$
A locally anisotropic generalization is to be found if we consider, for instance, that $`r_gr_g\left(\theta \right)`$ is defined by an anisotropic mass $`\stackrel{~}{m}(\theta )`$ and the locally isotropic values with $`r_g=cons`$ are changed into those with variable $`r_g\left(\theta \right).`$ The d–metric coefficient $`h_4(r,\theta ,\stackrel{~}{\phi })`$ and the corresponding N–connection components are taken as to solve the equations (62) and (63).
## IX Some Additional Examples
### A Two–soliton locally anisotropic solutions
Instead of one soliton solutions we can also consider la–deformations of multi–soliton configurations (as a review see ). In this Subsection we give an example of anholonomic, for simplicity, two–soliton configuration in general relativity. The d–metric to be constructed is of Class 1 with the h–component being a solution (type (79)) of the equation (78) for $`ϵ=1`$ and $`x^i=(t,x).`$
The horizontal component of d–metric is induced, via a conformal transform (94), from the so–called two soliton 2D Lorentz metric (here we follow the denotations from being adapted to locally anisotropic constructions)
$$d\stackrel{~}{s}^2=2\frac{FG}{F^2+G^2}dt^2+\frac{G^2F^2}{F^2+G^2}dx,$$
where for the two soliton solution
$`F`$ $`=`$ $`\mathrm{cot}\mu \mathrm{sinh}\left[\stackrel{~}{m}\mathrm{sin}\mu \gamma \left(t+vx\right)\right],`$ (201)
$`G`$ $`=`$ $`\mathrm{sinh}\left[\stackrel{~}{m}\mathrm{cos}\mu \gamma \left(xvt\right)\right]`$ (202)
or, for the soliton–anti-soliton solution,
$`F`$ $`=`$ $`\mathrm{cot}\mu \mathrm{cosh}\left[\stackrel{~}{m}\mathrm{sin}\mu \gamma \left(t+vx\right)\right],`$ (203)
$`G`$ $`=`$ $`\mathrm{cosh}\left[\stackrel{~}{m}\mathrm{cos}\mu \gamma \left(xvt\right)\right],`$ (204)
with constant parameters $`\mu ,\gamma `$ and $`v.`$ The function
$$u(t,x)=4\mathrm{tan}^1\left(F/G\right)$$
is a solution of 2D Euclidean sine–Gordon equation
$$_t^2u+_x^2u=\stackrel{~}{m}^2\mathrm{sin}u.$$
The locally anisotropic deformation is described by a la–dilaton fild $`\omega \left(x^i\right)`$ chosen to solve the Poisson equation (96) with the source $`\rho `$ (77) computed by using the two soliton function $`u(t,x).`$ In consequence, the h–metric is of the form
$$g=\mathrm{exp}[\omega \left(x^i\right)]\times \left[2\frac{FG}{F^2+G^2}dt^2+\frac{G^2F^2}{F^2+G^2}\right]dx.$$
(205)
The next step is the construction of a soliton like v–metric. Let, for simplicity,
$$h_3=a_3\left(x^i\right)=\mathrm{exp}[\omega \left(x^i\right)]\frac{G^2F^2}{F^2+G^2}$$
(206)
and $`h_4=h_4(x^i,z)`$ is to be defined by the equation (83), which for $`h_3/z=0`$ transforms into
$$h_4\frac{^2h_4}{z^2}\frac{1}{2}\left(\frac{h_4}{z}\right)^2\frac{k\mathrm{{\rm Y}}_1}{2}a_3\left(x^i\right)\left(h_4\right)^2=0,$$
(207)
where we consider a diagonal energy–momentum d–tensor $`\mathrm{{\rm Y}}_\beta ^\alpha =diag[\epsilon ,0,0,0].`$ Introducing a new variable $`h_4=\xi ^2,`$ the equation (207) transform into a linear second order differential equation on $`z`$ when coordinates $`x^i`$ are treated as parameters,
$$\frac{^2\xi }{z^2}+\lambda \left(x^i\right)\xi =0,$$
where $`\lambda \left(x^i\right)=\kappa \epsilon a_3\left(x^i\right)/2.`$ The general solution
$`\xi (x^i,z)`$ is
$$\xi =\{\begin{array}{ccc}\hfill c_1\mathrm{cosh}(z\sqrt{|\lambda |})+c_2\mathrm{sinh}(z\sqrt{|\lambda |})& ,& \lambda <0;\hfill \\ \hfill c_1+c_2z& ,& \lambda =0;\hfill \\ \hfill c_1\mathrm{cos}(z\sqrt{\lambda })+c_2\mathrm{sin}(z\sqrt{\lambda })& ,& \lambda >0,\hfill \end{array}$$
where $`c_1`$ and $`c_2`$ are some functions on $`x^i.`$
The coefficients $`h_3=a_3(x^i)`$ (206) and $`h_4=\xi ^2(x^i,z)`$ define a h–metric induced by a 2D two soliton equation. The complete d–metric solving the Einstein equations (46) is defined by considering v–coefficients of type (205).
### B Kadomtsev–Petviashvily structures and non–diagonal energy–momentum d–tensors
Such structures, for diagonal energy–momentum d–tensors and vacuum Einstein equations, where proven to exist in Subsection IVB, paragraph 2. Here we show that another type of three dimensional soliton structures could be generated by nondiagonal components $`\mathrm{{\rm Y}}_{31}`$ and $`\mathrm{{\rm Y}}_{32}.`$
For $`\mathrm{{\rm Y}}_1^1=\mathrm{{\rm Y}}_2^2=0`$ every function $`h_4=a_4\left(x^i\right)`$ solves the v–component of Einstein equations (62). Let us consider a function $`h_3=h_3(x^i,z).`$ If the anholonomic constraints on the system of reference are imposed by N–connection coefficients $`N_i^3=q_i,`$ when
$$q_i=2\kappa \mathrm{{\rm Y}}_{3i}h_3\left[h_3\left(h_3^{}\right)^2+ϵ\left(\dot{h}_3+6h_3h_3^{}+h_3^{\prime \prime \prime }\right)^{}\right]^1,$$
where $`ϵ=\pm 1,`$ the system of equations (63) reduces to the Kadomtsev–Petviashvili equation for $`h_3,`$
$$h_3^{}+ϵ\left(\dot{h}_3+6h_3h_3^{}+h_3^{\prime \prime \prime }\right)^{}=0.$$
The solution of Einstein equations is to completed by considering some functions $`N_i^4=n_i`$ satisfying (64) and a h–metric $`g_{ij}(x^k)`$ solving (61).
### C Anholonomic soliton like vacuum configurations
The main result of Belinski–Zakharov–Maison works was the proof that vacuum gravitational soliton like structures could be defined in the framework of general relativity with $`h_{ab}\left(x^i\right)`$ (from 4D metric (4)) being a solution of a generalized type of sine–Gordon equations. The function $`f(x^i)`$ (from (4)) is to be determined by some integral relations after the components $`h_{ab}\left(x^i\right)`$ have been constructed.
By reformulating the problem of definition of soliton like integral varieties of vacuum Einstein equations from the viewpoint of anholonomic frame structures, there are possible further generalizations and constructions of new classes of solutions.
For vanishing energy–momentum d–tensors the Einstein equations (61)–(64) transform into
$`2(g_1^{^{\prime \prime }}`$ $`+\ddot{g}_2){\displaystyle \frac{1}{g_2}}(\dot{g}_2^2+g_1^{}g_2^{}){\displaystyle \frac{1}{g_1}}(g_1^{2}+\dot{g}_1\dot{g}_2)=0;`$ (208)
$`h_4^{}`$ $`{\displaystyle \frac{1}{2h_4}}(h_4^{})^2{\displaystyle \frac{1}{2h_3}}h_3^{}h_4^{}=0;`$ (210)
$`2q_1h_4\left[\left({\displaystyle \frac{h_3^{}}{h_3}}\right)^2{\displaystyle \frac{h_3^{}}{h_3}}+{\displaystyle \frac{h_4^{}}{2h_4^{2}}}{\displaystyle \frac{h_3^{}h_4^{}}{2h_3h_4}}\right]+`$ (211)
$`\left[{\displaystyle \frac{\dot{h}_4}{h_4}}h_4^{}2\dot{h}_4^{}+{\displaystyle \frac{\dot{h}_3}{h_3}}h_4^{}\right]`$ $`=`$ $`0,`$ (212)
$`2q_2h_4\left[\left({\displaystyle \frac{h_3^{}}{h_3}}\right)^2{\displaystyle \frac{h_3^{}}{h_3}}+{\displaystyle \frac{h_4^{}}{2h_4^{2}}}{\displaystyle \frac{h_3^{}h_4^{}}{2h_3h_4}}\right]+`$ (213)
$`\left[{\displaystyle \frac{h_4^{}}{h_4}}h_4^{}2h_4^{}+{\displaystyle \frac{h_3^{}}{h_3}}h_4^{}\right]`$ $`=`$ $`0;`$ (214)
$$n_1^{}=0\text{ and }n_2^{}=0,$$
(215)
where we suppose that $`g_1,g_2,h_3`$ and $`h_4`$ are not zero.
The equation (208), relates two components and their first and second order partial derivatives of a diagonal h–metric $`g_1(x^i)`$ and $`g_2(x^i).`$ We can prescribe one of the components in order to find the second one by solving a second order partial differential equation. For instance, we can consider the h–metric to be induced by a soliton–dilaton solution (like in the Section IV, but for vacuum solitons the constants will be not defined by any components of the energy–momentum d–tensor).
Let us fix a soliton 2D solution with diagonal auxiliary metric
$`\stackrel{~}{g}_{ij}=diag\{\stackrel{~}{g}_1=ϵ\mathrm{sin}^2\left[v\left(x^i\right)/2\right],\stackrel{~}{g}_2=\mathrm{cos}^2\left[v\left(x^i\right)/2\right]\},`$ (216)
$`ϵ=\pm 1,`$ (217)
and model the local anisotropy by a la–dilaton field $`\omega \left(x^i\right)`$ relating the metric $`\stackrel{~}{g}_{ij}`$ with the h–components $`g_{ij}`$ via a conformal transform (94). The la–dilaton is to be found as a solution of the equations (96) where the source $`\rho \left(x^i\right)`$ is computed by using the formula (77).So, we conclude that vacuum h–metrics can be described by corresponding soliton–dilaton systems.
The equation (210) relates two components and their first and second order partial derivatives on $`z`$ of a diagonal v–metric $`h_3(x^i,z)`$ and $`h_4(x^i,z)`$ which depends on three variables. We also can prescribe one of these components (for instance, as was shown in details in Section IV B to be a solution of the Kadomtsev–Patviashvili, or (2+1) dimensional sine–Gordon equation; the Belinski–Zakharov–Maison solutions can be considered as some particular case soliton vacuum configurations which do not depend on variable $`z)`$ the second v–component being defined after solution of the resulted partial differential equation on $`z,`$ with the h–coordinates $`x^i`$ treated as parameters.
If the values $`h_3(x^i,z)`$ and $`h_4(x^i,z)`$ are defined, we have algebraic equations (212) for calculation of coefficients $`q_1(x^i,z)`$ and $`q_2(x^i,z).`$ The equations (215) are satisfied by arbitrary $`n_1(x^i,z)`$ and $`n_2(x^i,z)`$ depending linearly on the third variable $`z.`$
## X Conclusions
In this paper, we have elaborated a new method of construction of exact solutions of the Einstein equations by using anholonomic frames with associated nonlinear connection structures.
We analyzed 4D metrics
$$ds^2=g_{\alpha \beta }du^\alpha du^\beta $$
when $`g_{\alpha \beta }`$ are parametrized by matrices of type
$$\left[\begin{array}{cccc}g_1+q_1{}_{}{}^{2}h_{3}^{}+n_1{}_{}{}^{2}h_{4}^{}& 0& q_1h_3& n_1h_4\\ 0& g_2+q_2{}_{}{}^{2}h_{3}^{}+n_2{}_{}{}^{2}h_{4}^{}& q_2h_3& n_2h_4\\ q_1h_3& q_2h_3& h_3& 0\\ n_1h_4& n_2h_4& 0& h_4\end{array}\right]$$
(218)
with coefficients being some functions of necessary smooth class $`g_i=g_i(x^j),q_i=q_i(x^j,z),n_i=n_i(x^j,z),`$ $`h_a=h_a(x^j,z).`$ Latin indices run respectively $`i,j,k,\mathrm{}`$ $`=1,2`$ and $`a,b,c,\mathrm{}=3,4`$ and the local coordinates are denoted $`u^\alpha =(x^i,y^3=z,y^4).`$ A metric (218) can be diagonalized,
$$\delta s^2=g_i(x^j)\left(dx^i\right)^2+h_a(x^j,z)\left(\delta y^a\right)^2,$$
with respect to anholonomic frames (9) and (11), here we write down only the ’elongated’ differentials
$$\delta z=dz+q_i(x^j,z)dx^i,\delta y^4=dy^4+n_i(x^j,z)dx^i.$$
The key result of this paper is the proof that for the introduced ansatz the Einstein equations simplify substantially for 3D and 4D spacetimes, the variables being separated:
* The equation (61) with the non–trivial component of the Ricci tensor (53) relates two (so–called, horizontal) components of metric $`g_i`$ with the (so–called, vertical) values of the diagonal energy–momentum tensor. We proved that such components of metric could be described by soliton–dilaton and black hole like solutions with parameters being determined by vertical sources.
* Similarly, the equation (62) with the non–trivial component of the Ricci tensor (53) relates two vertical components of metric $`h_a`$ with the horizontal values of the diagonal energy–momentum tensor. The vertical coefficients of metric could depend on three variables $`(x^i,z)`$ and this equation contains their first and second derivatives on $`z,`$ the dependence on horizontal coordinates $`x^i`$ being parametric.
* As to the rest of equations (63) and (64) with corresponding non–trivial Ricci tensors (56) and (59), they form an algebraic system for definition of the nonlinear connection coefficients $`q_i(x^i,z)`$ and second order differential equation on $`z`$ for the nonlinear connection coefficients $`n_i(x^i,z)`$ after the functions $`h_a(x^i,z)`$ have been defined and non–diagonal components of energy–momentum tensor are given.
The Einstein equations consist a system of second order nonlinear partial differential equations whose particular solutions are selected from the general integral variety by imposing some physical motivated conditions on the type of singularities, horizon hypersurfaces, perturbative and/or non–perturbative behavior of background configurations, limit correspondences with some well known solutions, physical laws, symmetries and so on.
We investigated the conditions when from the class of solutions of 4D and 3D gravitational field equations parametrized by metric ansatzs of type (218) we can obtain some locally anisotropic generalizations of well known soliton–dilaton, black hole, cylinder and disk solutions.
In this paper we have shown that one can use solutions of generalized sine–Gordon equations in two and three dimensions to generate 4D solutions of Einstein gravity with soliton–dilaton parameters being related to 4D energy–momentum values. We have found a broad class of 2D, 3D and 4D black hole configurations with generic local anisotropy. Our results seem to indicate that there is a deep connection between black hole and soliton–dilaton states in gravitational theories of lower and 4D dimensions. Via nonlinear superpositions the lower dimensional locally anisotropic configurations induce similar structures in higher dimensions. We conclude that if the former direct applications of the 2D soliton–dilaton–black hole models (more naturally treated in the framework of 2D gravity and string theory) are very rough approximations for general relativity, after introducing of some well defined principles of nonlinear superposition, the lower dimensional solutions could be considered as some building blocks for construction of non–perturbative solutions in four dimensions.
We presented a series of computations involving the dynamics of locally anisotropic gravitational soliton deformations, black hole dynamics and constructed exact 4D and 3D solutions of the Einstein equations with horizons being (under corresponding dimension) of elliptic, rotation ellipsoidal, bipolar, elliptic cylinder and toroidal configuration. We showed that such solutions are naturally contained in general relativity and defined by corresponding anholonomic constraints, anisotropic distributions of masses and energy densities and could model some anisotropic nonlinear self–gravitational polarizations and renormalizations of gravitational and cosmological constants.
Our approach represents just a first step in the differential geometric and nonlinear analysis of the role that solitons and singular configurations with local anisotropy plays in 2D, 3D and 4D gravity. The natural developments of our approach would be to use nonlinear superpositions to describe the semiclassical and quantum dynamics of extremal black holes induced from string theory, the corresponding nonequilibrium thermodynamics of such black holes. One would be of interest supersymmetric extensions of the method and investigation of the mentioned non–perturbative structures in the framework of string theory. Work is in progress to address these issues.
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# 1. PRELIMINARIES
## 1. PRELIMINARIES
A derivation of an algebra $``$ over $`𝐂`$ (the complex numbers) is defined as a $`𝐂`$-linear map $`\mathrm{d}`$ from $``$ into a $``$-bimodule satisfying the Leibniz rule
$$\mathrm{d}(ab)=a\mathrm{d}b+\mathrm{d}ab\text{for all }a,b.$$
In this paper, $`\mathrm{d}ab`$ means $`(\mathrm{d}a)b`$. We set $`\mathrm{\Gamma }(\mathrm{d})=\mathrm{Lin}_𝐂\{a\mathrm{d}b|a,b\}`$ (the $`𝐂`$-linear span). We write $`\mathrm{d}^{}\mathrm{d}`$, if $`\mathrm{d}^{}`$ and $`\mathrm{d}`$ are derivations of $``$ and the $`𝐂`$-linear map $`\mathrm{\Gamma }(\mathrm{d})\mathrm{\Gamma }(\mathrm{d}^{}):a\mathrm{d}ba\mathrm{d}^{}b`$ is well-defined, and consider derivations $`\mathrm{d}`$, $`\mathrm{d}^{}`$ of $``$ identical, if $`\mathrm{d}^{}\mathrm{d}`$ and $`\mathrm{d}\mathrm{d}^{}`$. The set of derivations of $``$ with $``$ is a complete lattice, this follows from Prop. 1.1. If $``$ is a $``$-algebra, then $`a\mathrm{d}^{}b:=\mathrm{d}(b^{})a^{}`$ defines an involution on the set of derivations of $``$.
We denote by $`𝒜`$ a Hopf algebra over $`𝐂`$ with comultiplication $`\mathrm{\Delta }`$, counit $`\epsilon `$ and antipode $`S`$, cf. . We set $`=_𝐂`$ and
$$a_{(1)}a_{(2)}=\mathrm{\Delta }(a),a_{(1)}\mathrm{}a_{(n+1)}=a_{(1)}\mathrm{}a_{(n1)}\mathrm{\Delta }(a_{(n)})$$
for $`n=2,\mathrm{\hspace{0.17em}3},\mathrm{}`$ (Sweedler’s notation) and use the map $`{}_{}{}^{+}:𝒜𝒜`$ defined by $`a^+=a\epsilon (a)\mathrm{\hspace{0.17em}1}`$. We assume that $``$ is a subalgebra of $`𝒜`$ and a left coideal, i.e. $`\mathrm{\Delta }()𝒜`$, and call a derivation $`\mathrm{d}`$ of $``$ equivariant if and only if the $`𝐂`$-linear map
$$\mathrm{\Gamma }(\mathrm{d})𝒜\mathrm{\Gamma }(\mathrm{d}):a\mathrm{d}ba_{(1)}b_{(1)}a_{(2)}\mathrm{d}b_{(2)},a,b,$$
is well-defined. The notion of (equivariant) derivation is the same as that of (covariant) first order differential calculus introduced in , .
The algebra $``$ may be viewed as the function algebra of a quantum homogeneous space associated to the quantum group with the function algebra $`𝒜`$. Accordingly, as proposed in and , , the equivariant derivations may be considered constituents of flexibilized (deformed) laws of nature with differential operations which are supposed to be still invariant under the quantum group action.
In Section 2, a way is prepared for determining the equivariant derivations e.g. for the quantizations of symmetric spaces in , . It is used in Section 3 in the case of the quantum 2-sphere of Podleś for classifying the 2-dimensional covariant first order differential calculi. Their existence is proved by construction in Section 4.
## 2. ONE-TO-ONE CORRESPONDENCES
###### Theorem 1.
Let $`𝒜`$ be a Hopf algebra, $``$ a left coideal subalgebra of $`𝒜`$.
(i) Let $``$ be a right ideal of $`^+`$. Let $`p:^+^+/`$ be the canonical projection. Then $`a\mathrm{d}b:=ab_{(1)}p(b_{(2)}^+)`$ uniquely determines an equivariant derivation $`\mathrm{d}_{}:=\mathrm{d}`$ of $``$. Let $`\overline{\text{}}:𝒜𝒜/(^+𝒜)`$ be the canonical projection, $`\overline{\mathrm{\Delta }}:=(\overline{\text{}}\mathrm{id})\mathrm{\Delta }`$. Then $`^{}:=\overline{\mathrm{\Delta }}^1(\overline{𝒜})`$ is a right ideal of $`^+`$, $`^{}`$, such that $`\overline{\mathrm{\Delta }}(^{})\overline{𝒜}^{}`$ and $`\mathrm{d}_{}=\mathrm{d}_{^{}}`$.
(ii) Let $`\mathrm{d}`$ be a derivation of $``$. Then $`_\mathrm{d}:=\left\{_i\epsilon (a_i)b_i^+\right|_ia_i\mathrm{d}b_i=0\}`$ is a right ideal of $`^+`$. If $`\mathrm{d}`$ is equivariant, then $`\overline{\mathrm{\Delta }}(_\mathrm{d})\overline{𝒜}_\mathrm{d}`$.
(iii) The maps $`\mathrm{d}_{}`$, $`\mathrm{d}_\mathrm{d}`$ establish a one-to-one correspondence of
$``$ $`\{_\mathrm{d}|\mathrm{d}\text{ an equivariant derivation of }\}`$ and
$``$ $`\{\mathrm{d}_{}|\text{ a right ideal of }^+\}`$,
where $``$ for the derivations corresponds to $``$ for the right ideals. Furthermore, $`_\mathrm{d}_{}`$, if $``$ is a right ideal of $`^+`$, and $`\mathrm{d}__\mathrm{d}\mathrm{d}`$, if $`\mathrm{d}`$ is an equivariant derivation of $``$.
###### Proof.
(i) The left module operation on $`\mathrm{\Gamma }(\mathrm{d})=\mathrm{Lin}_𝐂\{a\mathrm{d}b|a,b\}`$ is determined by $`c(a\mathrm{d}b)=(ca)\mathrm{d}b`$ and the right module operation by $`(a\mathrm{d}b)c=a\mathrm{d}(bc)ab\mathrm{d}c`$. This proves uniqueness. To prove existence, we must show that the right module operation is well-defined (a) and satisfies the right module axioms (b), furthermore that $`\mathrm{d}`$ is equivariant (c). Well-definedness and axioms of the left module operation, the bimodule axiom and the Leibniz rule clearly hold true.
(a) If $`_ia_i\mathrm{d}b_i=0`$, that is, $`_ia_ib_{i(1)}p(b_{i(2)}^+)=0`$, then
$$\begin{array}{c}\left(_ia_i\mathrm{d}b_i\right)c\hfill \\ =_i(a_i\mathrm{d}(b_ic)a_ib_i\mathrm{d}c)\hfill \\ =_i(a_ib_{i(1)}c_{(1)}p((b_{i(2)}c_{(2)})^+)a_ib_ic_{(1)}p(c_{(2)}^+))\hfill \\ =_i(a_ib_{i(1)}c_{(1)}p(b_{i(2)}c_{(2)}\epsilon (b_{i(2)}c_{(2)}))\hfill \\ a_ib_{i(1)}c_{(1)}p(\epsilon (b_{i(2)})c_{(2)}\epsilon (b_{i(2)}c_{(2)})))\hfill \\ =_ia_ib_{i(1)}c_{(1)}p(b_{i(2)}^+c_{(2)})\hfill \\ =\left(_ia_ib_{i(1)}p(b_{i(2)}^+)\right)(c_{(1)}c_{(2)})=0,\hfill \end{array}$$
since $``$ is a right ideal of $``$. Hence, the specified right module operation is well-defined.
(b) We calculate that
$$\begin{array}{cc}\hfill ((a\mathrm{d}b)c)d& =(a\mathrm{d}(bc))d(ab\mathrm{d}c)d\hfill \\ & =a\mathrm{d}(bcd)abc\mathrm{d}dab\mathrm{d}(cd)+abc\mathrm{d}d=(a\mathrm{d}b)(cd)\hfill \end{array}$$
and, since $`\mathrm{d1}=1p(1^+)=0`$, $`(a\mathrm{d}b)\mathrm{\hspace{0.17em}1}=a\mathrm{d}(b\mathrm{\hspace{0.17em}1})ab\mathrm{d1}=a\mathrm{d}b`$.
(c) The map $`a\mathrm{d}ba_{(1)}b_{(1)}a_{(2)}\mathrm{d}b_{(2)}`$ is given by $`\mathrm{\Delta }\mathrm{id}`$.
We still have to show the assertions about $`^{}`$. Since $`\epsilon (^+𝒜)=\{0\}`$ and $`\mathrm{\Delta }(^+𝒜)^+𝒜𝒜+𝒜^+𝒜`$, the maps $`\epsilon _{\overline{𝒜}}:\overline{𝒜}𝐂:\overline{a}\epsilon (a)`$ and $`\overline{\mathrm{\Delta }}_{\overline{𝒜}}:\overline{𝒜}\overline{𝒜}\overline{𝒜}:\overline{a}\overline{a_{(1)}}\overline{a_{(2)}}`$ are well-defined, and $`\overline{𝒜}`$ with $`\overline{\mathrm{\Delta }}_{\overline{𝒜}}`$ is a coalgebra. If $`a^{}`$, then $`a=\epsilon _{\overline{𝒜}}(\overline{a_{(1)}})a_{(2)}`$, thus $`^{}`$. Because $``$ is a right ideal of $``$, $`\overline{\mathrm{\Delta }}(ab)\overline{𝒜}`$ for all $`b^+`$, so $`^{}`$ is a right ideal of $`^+`$. From $`\overline{a_{(1)}}\overline{\mathrm{\Delta }}(a_{(2)})=\overline{\mathrm{\Delta }}_{\overline{𝒜}}(\overline{a_{(1)}})a_{(2)}\overline{𝒜}\overline{𝒜}`$ follows $`\overline{a_{(1)}}a_{(2)}\overline{𝒜}\overline{\mathrm{\Delta }}^1(\overline{𝒜})=\overline{𝒜}^{}`$, thus $`\overline{\mathrm{\Delta }}(^{})\overline{𝒜}^{}`$. Let $`_ia_i\mathrm{d}b_i=0`$ with $`\mathrm{d}=\mathrm{d}_{}`$, that is, $`_ia_ib_{i(1)}b_{i(2)}^+𝒜`$. Then
$$\begin{array}{c}_ia_ib_{i(1)}\overline{\mathrm{\Delta }}(b_{i(2)}^+)\hfill \\ =_i\left(a_{i(1)}b_{i(1)}\overline{\epsilon (a_{i(2)})b_{i(2)}}b_{i(3)}^++a_ib_{i(1)}\overline{b_{i(2)}^+}1\right)\hfill \\ =_ia_{i(1)}b_{i(1)}\overline{a_{i(2)}b_{i(2)}}b_{i(3)}^+𝒜\overline{𝒜}.\hfill \end{array}$$
This implies that $`_ia_ib_{i(1)}b_{i(2)}^+𝒜\overline{\mathrm{\Delta }}^1(\overline{𝒜})=𝒜^{}`$, therefore $`_ia_i\mathrm{d}b_i=0`$ with $`\mathrm{d}=\mathrm{d}_{^{}}`$. Hence, $`\mathrm{d}_{^{}}\mathrm{d}_{}`$, while $`\mathrm{d}_{}\mathrm{d}_{^{}}`$ follows from $`^{}`$.
(ii) If $`_ia_i\mathrm{d}b_i=0`$, then $`_ia_i\mathrm{d}b_ic=_i(a_i\mathrm{d}(b_ic)a_ib_i\mathrm{d}c)=0`$ and
$$\begin{array}{c}_i(\epsilon (a_i)(b_ic)^+\epsilon (a_ib_i)c^+)\hfill \\ =_i(\epsilon (a_i)b_ic\epsilon (a_i)\epsilon (b_ic)\epsilon (a_ib_i)c+\epsilon (a_ib_i)\epsilon (c))\hfill \\ =_i\epsilon (a_i)b_i^+c_\mathrm{d}.\hfill \end{array}$$
Therefore, $`_\mathrm{d}`$ is a right ideal of $`^+`$. If $`\mathrm{d}`$ is equivariant and $`_ia_i\mathrm{d}b_i=0`$, then $`_if(a_{i(1)}b_{i(1)})a_{i(2)}\mathrm{d}b_{i(2)}=0`$ for any $`𝐂`$-linear functional $`f`$ on $`𝒜`$. This implies that $`_if(a_ib_{i(1)})b_{i(2)}^+_\mathrm{d}`$. From this and
$$\begin{array}{cc}\hfill \overline{\mathrm{\Delta }}\left(_i\epsilon (a_i)b_i^+\right)& =_i\left(\overline{\epsilon (a_i)b_{i(1)}}b_{i(2)}\overline{\epsilon (a_ib_i)\mathrm{\hspace{0.17em}1}}1\right)\hfill \\ & =_i\left(\overline{a_ib_{i(1)}}b_{i(2)}\overline{a_ib_i}1\right)\hfill \\ & =_i\overline{a_ib_{i(1)}}b_{i(2)}^+\hfill \end{array}$$
we conclude that $`\overline{\mathrm{\Delta }}(_\mathrm{d})\overline{𝒜}_\mathrm{d}`$.
(iii) Directly from the definitions, we have
(a) $`(^{})(\mathrm{d}_{}\mathrm{d}_{^{}})`$ for right ideals $``$, $`^{}`$ of $`^+`$, and
(b) $`(\mathrm{d}^{}\mathrm{d})(_\mathrm{d}_\mathrm{d}^{})`$ for derivations $`\mathrm{d}`$, $`\mathrm{d}^{}`$ of $``$.
We show that, in addition,
(c) $`_\mathrm{d}_{}`$, if $``$ is a right ideal of $`^+`$, and
(d) $`\mathrm{d}__\mathrm{d}\mathrm{d}`$, if $`\mathrm{d}`$ is an equivariant derivation of $``$.
(c) Let $`_ia_i\mathrm{d}b_i=0`$ with $`\mathrm{d}=\mathrm{d}_{}`$, that is, $`_ia_ib_{i(1)}p(b_{i(2)}^+)=0`$ with the canonical projection $`p:^+^+/`$. Application of $`\epsilon \mathrm{id}`$ leads to $`p\left(_i\epsilon (a_i)b_i^+\right)=0`$, therefore $`_i\epsilon (a_i)b_i^+`$.
(d) If $`_ia_i\mathrm{d}b_i=0`$, then $`_if(a_{i(1)}b_{i(1)})a_{i(2)}\mathrm{d}b_{i(2)}=0`$ for any $`𝐂`$-linear functional $`f`$ on $`𝒜`$, since $`\mathrm{d}`$ is equivariant. Thus $`_if(a_ib_{i(1)})b_{i(2)}^+_\mathrm{d}`$ and $`_if(a_ib_{i(1)})p(b_{i(2)}^+)=0`$, where $`p:^+^+/_\mathrm{d}`$ is the canonical projection. This implies that $`_ia_i\mathrm{d}__\mathrm{d}b_i=_ia_ib_{i(1)}p(b_{i(2)}^+)=0`$.
If $`\mathrm{d}`$ is an equivariant derivation of $``$, then $`_{\mathrm{d}__\mathrm{d}}=_\mathrm{d}`$ by (c) and (d), (b). If $``$ is a right ideal of $`^+`$, then $`\mathrm{d}__\mathrm{d}_{}=\mathrm{d}_{}`$ by (d) and (c), (a). This proves assertion (iii). ∎
In particular, the trivial derivation $`\mathrm{d}_^+`$, the universal derivation $`\mathrm{d}_{\{0\}}`$ and, if $`𝒜`$ is commutative, the commutative universal derivation $`\mathrm{d}_{(^+)^2}`$ occur in the one-to-one correspondence. The equivariant derivations of $``$ induced from those of $`𝒜`$ and e.g. the calculi in , , , also correspond to right ideals of $`^+`$ in this way. For $`=𝒜`$, Woronowicz shows that all equivariant derivations have this property. We generalize this result using a theorem of Takeuchi which requires further notations.
Let $`𝒞`$ be a coalgebra, $`𝒲`$ a right $`𝒞`$-comodule and $`𝒱`$ a left $`𝒞`$-comodule. The cotensor product $`𝒲\mathrm{}_𝒞𝒱`$ is defined as the subspace
$$\left\{_iw_iv_i𝒲𝒱\right|_iw_{i(1)}w_{i(2)}v_i=_iw_iv_{i(1)}v_{i(2)}\}$$
of $`𝒲𝒱`$; the Sweedler notation is used for the $`𝒞`$-comodule operations $`𝒲𝒲𝒞`$ and $`𝒱𝒞𝒱`$. The category of the left $``$-modules $``$ with left $`𝒜`$-comodule structure such that $`(bm)_{(1)}(bm)_{(2)}=b_{(1)}m_{(1)}b_{(2)}m_{(2)}`$ for all $`m`$ and $`b`$, together with the $``$-linear, $`𝒜`$-colinear maps, is denoted by $`{}_{}{}^{𝒜}\mathrm{M}`$. The category of the left $`\overline{𝒜}`$-comodules, together with the $`\overline{𝒜}`$-colinear maps, is denoted by $`{}_{}{}^{\overline{𝒜}}\mathrm{M}`$. Generally, if $``$ is a left $``$-module, $`\overline{\text{}}:/(^+)`$ denotes the canonical projection. We have shown in the proof of Theorem 1 (i) that $`\overline{𝒜}`$ is a coalgebra. Correspondingly, $`\overline{}`$ is an object of $`{}_{}{}^{\overline{𝒜}}\mathrm{M}`$, if $``$ is an object of $`{}_{}{}^{𝒜}\mathrm{M}`$. Moreover, $`𝒜\mathrm{}_{\overline{𝒜}}𝒱`$ with the induced structure of $`𝒜`$ is an object of $`{}_{}{}^{𝒜}\mathrm{M}`$, if $`𝒱`$ is an object of $`{}_{}{}^{\overline{𝒜}}\mathrm{M}`$. That a right $``$-module $`𝒜`$ is faithfully flat means that the functor $`𝒜_{}`$ from the category of left $``$-modules to the category of $`𝐂`$-vector spaces preserves and reflects exact sequences.
The following result, actually the equivalent one with the opposite multiplication and comultiplication, is contained in , proof of Theorem 1.
Theorem (Takeuchi). Let $`𝒜`$ be a Hopf algebra and $``$ a left coideal subalgebra of $`𝒜`$. If $`𝒜`$ is faithfully flat as a right $``$-module, then the maps
$$\begin{array}{c}\mathrm{\Xi }:𝒜\mathrm{}_{\overline{𝒜}}\overline{}:mm_{(1)}\overline{m_{(2)}}\text{and}\hfill \\ \mathrm{\Theta }:\overline{𝒜\mathrm{}_{\overline{𝒜}}𝒱}𝒱:\overline{_ia_iv_i}_i\epsilon (a_i)v_i\hfill \end{array}$$
are bijective for all objects $``$ of $`{}_{}{}^{𝒜}\mathrm{M}`$ and all objects $`𝒱`$ of $`{}_{}{}^{\overline{𝒜}}\mathrm{M}`$. ∎
###### Theorem 2.
Let $`𝒜`$ be a Hopf algebra, $``$ a left coideal subalgebra of $`𝒜`$. If $`𝒜`$ is faithfully flat as a right $``$-module, then $`\mathrm{d}_{}`$, $`\mathrm{d}_\mathrm{d}`$ as in Theorem 1 establish a one-to-one correspondence between
$``$ the right ideals $``$ of $`^+`$ with $`\overline{\mathrm{\Delta }}()\overline{𝒜}`$ and
$``$ the equivariant derivations of $``$.
###### Proof.
If $`\mathrm{d}`$ is an equivariant derivation of $``$, then $`\overline{\mathrm{\Delta }}(_\mathrm{d})\overline{𝒜}_\mathrm{d}`$ according to Theorem 1 (ii), so it remains to show that
(a) $`\mathrm{d}__\mathrm{d}=\mathrm{d}`$, if $`\mathrm{d}`$ is an equivariant derivation of $``$, and
(b) $`_\mathrm{d}_{}=`$, if $``$ is a right ideal of $`^+`$ with $`\overline{\mathrm{\Delta }}()\overline{𝒜}`$.
(a) Let $`\mathrm{d}`$ be an equivariant derivation of $``$. Then $`\mathrm{\Gamma }(\mathrm{d})`$ is an object of $`{}_{}{}^{𝒜}\mathrm{M}`$. According to Takeuchi’s Theorem, the map
$$\mathrm{\Xi }:\mathrm{\Gamma }(\mathrm{d})𝒜\mathrm{}_{\overline{𝒜}}\overline{\mathrm{\Gamma }(\mathrm{d})}:a\mathrm{d}ba_{(1)}b_{(1)}\overline{a_{(2)}\mathrm{d}b_{(2)}}=ab_{(1)}\overline{\mathrm{d}b_{(2)}}$$
is bijective. The kernel of $`\iota :^+\overline{\mathrm{\Gamma }(\mathrm{d})}:c\overline{\mathrm{d}c}`$ is $`_\mathrm{d}`$: If $`_ia_i\mathrm{d}b_i=0`$, then $`\overline{\mathrm{d}\left(_i\epsilon (a_i)b_i^+\right)}=\overline{_i\epsilon (a_i)\mathrm{d}b_i}=\overline{_ia_i\mathrm{d}b_i}=0`$, and if $`\overline{\mathrm{d}c}=0`$, then $`c_{ij},a_i,b_j`$ exist such that $`\mathrm{d}c_{ij}c_{ij}^+a_i\mathrm{d}b_j=0`$, therefore $`c^+_\mathrm{d}`$. The $`𝐂`$-linear map $`(\mathrm{id}\iota ^1)\mathrm{\Xi }:\mathrm{\Gamma }(\mathrm{d})\mathrm{\Gamma }(\mathrm{d}__\mathrm{d}):a\mathrm{d}ba\mathrm{d}__\mathrm{d}b`$ is injective and surjective, thus $`\mathrm{d}__\mathrm{d}=\mathrm{d}`$.
(b) Let $``$ be a right ideal of $`^+`$ with $`\overline{\mathrm{\Delta }}()\overline{𝒜}`$. Then $`^+/`$ is an object of $`{}_{}{}^{\overline{𝒜}}\mathrm{M}`$. Furthermore $`\mathrm{\Gamma }(\mathrm{d}_{})𝒜\mathrm{}_{\overline{𝒜}}(^+/)`$, which follows from $`a\mathrm{d}_{}b=ab_{(1)}p(b_{(2)}^+)`$, where $`p:^+^+/`$ is the canonical projection, and $`a_{(1)}b_{(1)}\overline{a_{(2)}b_{(2)}}p(b_{(3)}^+)=ab_{(1)}\overline{b_{(2)}}p(b_{(3)}^+)`$. According to Takeuchi’s Theorem, the map
$$\mathrm{\Theta }|_{\overline{\mathrm{\Gamma }(\mathrm{d}_{})}}:\overline{\mathrm{\Gamma }(\mathrm{d}_{})}^+/:\overline{a\mathrm{d}_{}b}=\overline{ab_{(1)}p(b_{(2)}^+)}\epsilon (a)p(b^+)$$
is injective. The kernel of $`\iota :^+\overline{\mathrm{\Gamma }(\mathrm{d}_{})}:c\overline{\mathrm{d}_{}c}`$ is $`_\mathrm{d}_{}`$, see the proof of (a), thus $`\mathrm{ker}(\mathrm{\Theta }\iota )=_\mathrm{d}_{}`$. Since $`\mathrm{\Theta }\iota =p`$, we obtain $`_\mathrm{d}_{}=`$. ∎
In Section 4, we give examples of equivariant derivations which do not arise from a right ideal as in Theorem 1 (i), so the statement of the theorem without the condition of faithful flatness is false. However, due to Müller and Schneider this condition is verified for the quantizations of symmetric spaces by Noumi, Dijkhuizen and Sugitani , and for the quantized flag manifolds .
## 3. CLASSIFICATION
We call $`dim_{\epsilon ,\mathrm{l}}\mathrm{d}:=dim_𝐂(\mathrm{\Gamma }(\mathrm{d})/(^+\mathrm{\Gamma }(\mathrm{d})))`$ the left dimension and analogously $`dim_{\epsilon ,\mathrm{r}}\mathrm{d}:=dim_𝐂(\mathrm{\Gamma }(\mathrm{d})/(\mathrm{\Gamma }(\mathrm{d})^+))`$ the right dimension of a first order differential calculus $`\mathrm{d}`$ over $``$ at the classical point $`\epsilon `$. If $``$ is the algebra of regular functions on a nonsingular affine algebraic variety, then $`dim_{\epsilon ,\mathrm{l}}\mathrm{d}_{(^+)^2}=dim_{\epsilon ,\mathrm{r}}\mathrm{d}_{(^+)^2}`$ is the dimension of it. If a basis of $`\mathrm{\Gamma }(\mathrm{d})`$ as a left $``$-module exists, e.g. if $`=𝒜`$ and $`\mathrm{d}`$ is left-covariant, cf. , then $`dim_{\epsilon ,\mathrm{l}}\mathrm{d}`$ is the number of its elements.
We assume that $`q𝐂\{0\}`$ and $`q^n1`$ for all $`n=1,\mathrm{\hspace{0.17em}2},\mathrm{}`$. For the quantum 2-sphere of Podleś we may equivalently choose $`𝒜`$ to be one of the quantum group function algebras $`𝒪(\mathrm{SL}_q(2))`$, $`𝒪(\mathrm{SO}_{q^2}(3))`$ and $`𝒪(\mathrm{Sp}_{q^{1/2}}(2))`$ which are described in . The function algebras $`𝒪(\mathrm{S}_{qc}^2)`$ of the quantum 2-sphere, parameterized by $`c\mathrm{𝐂𝐏}^1`$, are the $`𝒜`$-comodule algebras (except for one) which are isomorphic as a comodule to the classical case and generated as an algebra by the spin 1 subcomodule, cf. . They are isomorphic to right coideal subalgebras $`_c`$ of $`𝒜`$, so the equivalents of our theorems with the opposite multiplication and comultiplication are applicable (we silently assume this exchange of left and right). The algebras $`_c`$ are generated by three elements $`e_1`$, $`e_0`$, $`e_1`$ with the relations
$$\begin{array}{c}(q^2+1)e_1e_1+e_0^2+(q^2+1)e_1e_1=\rho \mathrm{\hspace{0.17em}1},\hfill \\ q^2e_1e_0+e_0e_1=\lambda e_1,\hfill \\ (q^2+1)e_1e_1(q^21)e_0^2(q^2+1)e_1e_1=\lambda e_0,\hfill \\ q^2e_0e_1+e_1e_0=\lambda e_1,\hfill \end{array}\rho ,\lambda 𝐂,()$$
such that $`c=\epsilon (e_1)\epsilon (e_1):\epsilon (e_0)^2`$ and $`\mathrm{\Delta }(e_i)=_je_j\pi _i^j`$, cf. ($`q`$ is $`\mu `$, $`\pi _j^i`$ is $`d_{1,ij}`$). Special values of $`c`$ are $`c(n)=q^{2n}/(q^{2n}+1)^2`$.
We classify the equivariant derivations $`\mathrm{d}`$ of $`_c`$ with $`dim_{\epsilon ,\mathrm{r}}\mathrm{d}=2`$ which arise from a left ideal as in Theorem 1 (i). If $`cc(n)`$ for all $`n=0,\mathrm{\hspace{0.17em}1},\mathrm{}`$, then $`𝒜`$ is faithfully flat as a left $`_c`$-module, cf. , , , and according to Theorem 2 our classification includes all equivariant derivations with $`dim_{\epsilon ,\mathrm{r}}\mathrm{d}=2`$. We denote by $`\overline{\text{}}`$, $`\overline{\mathrm{\Delta }}`$ and $`_\mathrm{d}`$ the equivalents of the previously used structures with left and right reversed. The map $`\chi :_c\overline{\mathrm{\Gamma }(\mathrm{d})}:b\overline{\mathrm{d}b}`$ induces a $`𝐂`$-linear bijection between $`_c^+/_\mathrm{d}`$ and $`\overline{\mathrm{\Gamma }(\mathrm{d})}`$, see the proof of Theorem 2, i.e. we must determine the left ideals $``$ of $`_c^+`$ with $`dim_𝐂(_c^+/)=2`$ and $`\overline{\mathrm{\Delta }}()\overline{𝒜}`$. The Hochschild coboundary maps of the quotient $`_c`$-bimodule $`\overline{\mathrm{\Gamma }(\mathrm{d})}`$ are defined as
$$\begin{array}{c}\delta ^0:\overline{\mathrm{\Gamma }(\mathrm{d})}\mathrm{Hom}_𝐂(_c,\overline{\mathrm{\Gamma }(\mathrm{d})}):\overline{\omega }(bb\overline{\omega }\overline{\omega }b)\text{and}\hfill \\ \delta ^n:\mathrm{Hom}_𝐂(_c^n,\overline{\mathrm{\Gamma }(\mathrm{d})})\mathrm{Hom}_𝐂(_c^{(n+1)},\overline{\mathrm{\Gamma }(\mathrm{d})}),\hfill \end{array}$$
$`\delta ^nf(b_0\mathrm{}b_n)=b_0f(b_1\mathrm{}b_n)+\underset{i=1}{\overset{n}{}}(1)^if(b_0\mathrm{}b_{i1}b_i\mathrm{}b_n)+(1)^{n+1}f(b_0\mathrm{}b_{n1})b_n`$
for $`n=1,\mathrm{\hspace{0.17em}2},\mathrm{}`$. Hence, $`\chi `$ is a 1-cocycle: $`a\chi (b)\chi (ab)+\chi (a)b=0`$ for all $`a,b_c`$. Let $`\tau :_c\mathrm{Hom}_𝐂(\overline{\mathrm{\Gamma }(\mathrm{d})},\overline{\mathrm{\Gamma }(\mathrm{d})})`$ be the representation of $`_c`$ on the quotient left $`_c`$-module $`\overline{\mathrm{\Gamma }(\mathrm{d})}`$. Using coordinates, this says
$$\begin{array}{c}\chi _i(ab)=_k\tau _{ik}(a)\chi _k(b)+\chi _i(a)\epsilon (b),\chi _i(1)=0,\hfill \\ \tau _{ij}(ab)=_k\tau _{ik}(a)\tau _{kj}(b),\tau _{ij}(1)=\delta _{ij}\hfill \end{array}$$
for all $`a,b_c`$ and $`i,j=1,\mathrm{\hspace{0.17em}2}`$. Given a representation $`\tau `$, the 1-coboundaries $`\chi _i=_k\beta _k\tau _{ik}\beta _i\epsilon `$ with $`\beta _1,\beta _2𝐂`$ are solutions to these equations, and further solutions exist exactly if the first cohomology group $`(\mathrm{ker}\delta ^1)/(\mathrm{im}\delta ^0)`$ is not $`\{0\}`$. The equations imply that the functions $`\chi _i`$ and $`\tau _{ij}`$ are uniquely determined by their values on $`e_1`$, $`e_0`$, $`e_1`$ and exist for given values if they are compatible with the relations $`()`$. The solutions in suitable coordinates are $`\chi _i=_k\beta _k\tau _{ik}\beta _i\epsilon +_n\beta _n^{}\xi _i^n`$ with $`\beta _i,\beta _n^{}𝐂`$,
(a) $`\tau (e_1,e_0,e_1)=x^1\left(\begin{array}{cc}q^2\alpha _1\alpha _1& \alpha _0\\ 0& \alpha _1\alpha _1\end{array}\right),\left(\begin{array}{cc}\alpha _0& \left(q^2+1\right)\\ 0& \alpha _0\end{array}\right),x\left(\begin{array}{cc}q^2& 0\\ 0& 1\end{array}\right)`$,
$`\xi ^1(e_1,e_0,e_1)=\left(\begin{array}{c}\alpha _1\\ 0\end{array}\right),\left(\begin{array}{c}0\\ 0\end{array}\right),\left(\begin{array}{c}\alpha _1\\ 0\end{array}\right)`$ if $`x=q^2\alpha _1`$ and
$`\xi ^1(e_1,e_0,e_1)=\left(\begin{array}{c}1\\ \alpha _0\end{array}\right),\left(\begin{array}{c}0\\ \left(q^2+1\right)\alpha _1\end{array}\right),\left(\begin{array}{c}0\\ 0\end{array}\right)`$ if $`x=q^2\alpha _1`$,
(b) $`\tau (e_1,e_0,e_1)=x^1\alpha _1\alpha _1\left(\begin{array}{cc}1& 1\\ 0& 1\end{array}\right),\left(\begin{array}{cc}\alpha _0& 0\\ 0& \alpha _0\end{array}\right),x\left(\begin{array}{cc}1& 1\\ 0& 1\end{array}\right)`$,
$`\xi ^1(e_1,e_0,e_1)=\left(\begin{array}{c}\alpha _1\\ \alpha _1\end{array}\right),\left(\begin{array}{c}0\\ 0\end{array}\right),\left(\begin{array}{c}0\\ \alpha _1\end{array}\right)`$ if $`x=\alpha _1`$ and
$`\xi ^1(e_1,e_0,e_1)=\left(\begin{array}{c}\alpha _0\\ 0\end{array}\right),\left(\begin{array}{c}\left(q^2+1\right)\alpha _1\\ 0\end{array}\right),\left(\begin{array}{c}0\\ 0\end{array}\right)`$ if $`x=q^2\alpha _1`$,
(c) $`\tau (e_1,e_0,e_1)=\alpha _1\alpha _1\left(\begin{array}{cc}x^1& 0\\ 0& y^1\end{array}\right),\left(\begin{array}{cc}\alpha _0& 0\\ 0& \alpha _0\end{array}\right),\left(\begin{array}{cc}x& 0\\ 0& y\end{array}\right)`$,
$`\xi ^1(e_1,e_0,e_1)=\left(\begin{array}{c}\alpha _1\\ 0\end{array}\right),\left(\begin{array}{c}0\\ 0\end{array}\right),\left(\begin{array}{c}\alpha _1\\ 0\end{array}\right)`$ if $`x=\alpha _1`$ and
$`\xi ^1(e_1,e_0,e_1)=\left(\begin{array}{c}\alpha _0\\ 0\end{array}\right),\left(\begin{array}{c}\left(q^2+1\right)\alpha _1\\ 0\end{array}\right),\left(\begin{array}{c}0\\ 0\end{array}\right)`$ if $`x=q^2\alpha _1`$,
$`\xi ^2(e_1,e_0,e_1)=\left(\begin{array}{c}0\\ \alpha _1\end{array}\right),\left(\begin{array}{c}0\\ 0\end{array}\right),\left(\begin{array}{c}0\\ \alpha _1\end{array}\right)`$ if $`y=\alpha _1`$ and
$`\xi ^2(e_1,e_0,e_1)=\left(\begin{array}{c}0\\ \alpha _0\end{array}\right),\left(\begin{array}{c}0\\ \left(q^2+1\right)\alpha _1\end{array}\right),\left(\begin{array}{c}0\\ 0\end{array}\right)`$ if $`y=q^2\alpha _1`$,
where $`x,y𝐂\{0\}`$, $`\alpha _i=\epsilon (e_i)`$ and, if not specified otherwise, $`\xi ^j=0`$; in addition, if $`c=c(1)`$,
(d) $`\tau (e_1,e_0,e_1)=\alpha _1\alpha _1\left(\begin{array}{cc}x^1& 0\\ 0& 0\end{array}\right),\left(\begin{array}{cc}\alpha _0& 0\\ 0& 0\end{array}\right),\left(\begin{array}{cc}x& 0\\ 0& 0\end{array}\right)`$,
$`\xi ^1(e_1,e_0,e_1)=\left(\begin{array}{c}\alpha _1\\ 0\end{array}\right),\left(\begin{array}{c}0\\ 0\end{array}\right),\left(\begin{array}{c}\alpha _1\\ 0\end{array}\right)`$ if $`x=\alpha _1`$ and
$`\xi ^1(e_1,e_0,e_1)=\left(\begin{array}{c}\alpha _0\\ 0\end{array}\right),\left(\begin{array}{c}\left(q^2+1\right)\alpha _1\\ 0\end{array}\right),\left(\begin{array}{c}0\\ 0\end{array}\right)`$ if $`x=q^2\alpha _1`$,
(e) $`\tau (e_1,e_0,e_1)=\left(\begin{array}{cc}0& 0\\ 0& 0\end{array}\right),\left(\begin{array}{cc}0& 0\\ 0& 0\end{array}\right),\left(\begin{array}{cc}0& 0\\ 0& 0\end{array}\right)`$
and if $`c=c(2)`$,
(f) $`\tau (e_1,e_0,e_1)=\frac{q^21}{q^4+1}\alpha _0\left(\begin{array}{cc}0& 1\\ 0& 0\end{array}\right),\frac{q^21}{q^4+1}\alpha _0\left(\begin{array}{cc}1& 0\\ 0& q^2\end{array}\right),\frac{q^21}{q^4+1}\alpha _0\left(\begin{array}{cc}0& 0\\ q^2& 0\end{array}\right)`$
and if $`c=0`$,
(a’) $`\tau (e_1,e_0,e_1)=x\left(\begin{array}{cc}q^2& 0\\ 0& 1\end{array}\right),\left(\begin{array}{cc}\alpha _0& \left(q^2+1\right)\\ 0& \alpha _0\end{array}\right),x^1\left(\begin{array}{cc}0& \alpha _0\\ 0& 0\end{array}\right)`$,
$`\xi ^1(e_1,e_0,e_1)=\left(\begin{array}{c}1\\ 0\end{array}\right),\left(\begin{array}{c}0\\ 0\end{array}\right),\left(\begin{array}{c}0\\ 0\end{array}\right)`$ if $`x=q^2\alpha _1`$ and
$`\xi ^1(e_1,e_0,e_1)=\left(\begin{array}{c}0\\ 0\end{array}\right),\left(\begin{array}{c}0\\ \left(q^2+1\right)\alpha _1\end{array}\right),\left(\begin{array}{c}1\\ \alpha _0\end{array}\right)`$ if $`x=q^2\alpha _1`$,
(b’) $`\tau (e_1,e_0,e_1)=\left(\begin{array}{cc}0& s\\ 0& 0\end{array}\right),\left(\begin{array}{cc}\alpha _0& 0\\ 0& \alpha _0\end{array}\right),\left(\begin{array}{cc}t& 1\\ 0& t\end{array}\right)`$ with $`st=0`$,
$`\xi ^1(e_1,e_0,e_1)=\left(\begin{array}{c}0\\ q^2s\alpha _0\end{array}\right),\left(\begin{array}{c}\left(q^2+1\right)s\\ 0\end{array}\right),\left(\begin{array}{c}0\\ \alpha _0\end{array}\right)`$, if $`\alpha _1=0`$ and $`t=\alpha _1`$,
$`\xi ^2(e_1,e_0,e_1)=\left(\begin{array}{c}\alpha _0\\ 0\end{array}\right),\left(\begin{array}{c}\left(q^2+1\right)\alpha _1\\ 0\end{array}\right),\left(\begin{array}{c}0\\ 0\end{array}\right)`$, if $`\alpha _1=0`$ and $`t=q^2\alpha _1`$,
(b”) $`\tau (e_1,e_0,e_1)=\left(\begin{array}{cc}s& 1\\ 0& s\end{array}\right),\left(\begin{array}{cc}\alpha _0& 0\\ 0& \alpha _0\end{array}\right),\left(\begin{array}{cc}0& 0\\ 0& 0\end{array}\right)`$,
$`\xi ^1(e_1,e_0,e_1)=\left(\begin{array}{c}0\\ 1\end{array}\right),\left(\begin{array}{c}0\\ 0\end{array}\right),\left(\begin{array}{c}0\\ 0\end{array}\right)`$, if $`\alpha _1=0`$ and $`s=\alpha _1`$,
$`\xi ^2(e_1,e_0,e_1)=\left(\begin{array}{c}0\\ 0\end{array}\right),\left(\begin{array}{c}\left(q^2+1\right)\alpha _1\\ 0\end{array}\right),\left(\begin{array}{c}\alpha _0\\ 0\end{array}\right)`$, if $`\alpha _1=0`$ and $`s=q^2\alpha _1`$,
(c’) $`\tau (e_1,e_0,e_1)=\left(\begin{array}{cc}s& 0\\ 0& u\end{array}\right),\left(\begin{array}{cc}\alpha _0& 0\\ 0& \alpha _0\end{array}\right),\left(\begin{array}{cc}t& 0\\ 0& v\end{array}\right)`$ with $`st=uv=0`$,
$`\xi ^1(e_1,e_0,e_1)=\left(\begin{array}{c}0\\ 0\end{array}\right),\left(\begin{array}{c}\left(q^2+1\right)\alpha _1\\ 0\end{array}\right),\left(\begin{array}{c}\alpha _0\\ 0\end{array}\right)`$, if $`s=q^2\alpha _1`$ and $`t=\alpha _1`$,
$`\xi ^2(e_1,e_0,e_1)=\left(\begin{array}{c}\alpha _0\\ 0\end{array}\right),\left(\begin{array}{c}\left(q^2+1\right)\alpha _1\\ 0\end{array}\right),\left(\begin{array}{c}0\\ 0\end{array}\right)`$, if $`s=\alpha _1`$ and $`t=q^2\alpha _1`$,
$`\xi ^3(e_1,e_0,e_1)=\left(\begin{array}{c}0\\ 0\end{array}\right),\left(\begin{array}{c}0\\ \left(q^2+1\right)\alpha _1\end{array}\right),\left(\begin{array}{c}0\\ \alpha _0\end{array}\right)`$, if $`u=q^2\alpha _1`$ and $`v=\alpha _1`$,
$`\xi ^4(e_1,e_0,e_1)=\left(\begin{array}{c}0\\ \alpha _0\end{array}\right),\left(\begin{array}{c}0\\ \left(q^2+1\right)\alpha _1\end{array}\right),\left(\begin{array}{c}0\\ 0\end{array}\right)`$, if $`u=\alpha _1`$ and $`v=q^2\alpha _1`$,
where $`s,t,u,v𝐂`$. The solutions (b’), (c’) contain (b), (c).
The condition $`\overline{\mathrm{\Delta }}()\overline{𝒜}`$ implies that $`\chi (b_{(1)})f(b_{(2)})=0`$ for all $`b`$ and $`𝐂`$-linear functionals $`f`$ on $`𝒜`$ with $`f(𝒜_c^+)=\{0\}`$, therefore $`\chi _1f,\chi _2f\mathrm{Lin}_𝐂\{\chi _1,\chi _2\}`$. Such a functional (taken from ) is
$$f=\alpha _0(l^2\epsilon )\alpha _1l^{(+)}l+\alpha _1l^{()}l$$
with $`l=l_{11}^{(+)}`$, $`l^{(+)}=l_{12}^{(+)}`$ and $`l^{()}=l_{21}^{()}`$ as specified for $`𝒜=𝒪(\mathrm{SL}_q(2))`$ in . The conditions $`\chi _1f,\chi _2f\mathrm{Lin}_𝐂\{\chi _1,\chi _2\}`$, checked by evaluation of $`\chi _1`$, $`\chi _2`$, $`\chi _1f`$ and $`\chi _2f`$ on $`e_i`$, $`e_ie_j`$ and $`e_ie_je_k`$ for $`i,j,k=1,\mathrm{\hspace{0.17em}0},\mathrm{\hspace{0.17em}1}`$, and $`dim_𝐂\mathrm{Lin}_𝐂\{\chi _1,\chi _2\}=2`$ reduce the total number of solutions to seven:
1.+2. $`c=q/(\pm q+1)^2`$, (a) with $`x=\pm q^1\alpha _1`$ and $`\frac{\beta _1}{\beta _2}=\frac{\alpha _1^2}{\alpha _0x(x\alpha _1)}`$.
3.+4. $`c=c(2)`$, (f) with $`\frac{\beta _1}{\beta _2}\{\frac{(q^4+1)\alpha _1}{q^4\alpha _0},\frac{(q^4+1)\alpha _1}{\alpha _0}\}`$.
5.–7. $`c=0`$, $`\alpha _1=0=\alpha _1`$,
(b”) with $`s=0`$ and $`\beta _1^{}0=\beta _2^{}`$,
(b’) with $`s=t=0`$ and $`\beta _1^{}0=\beta _2^{}`$,
(c’) with $`s=t=u=v=0`$ and $`\beta _1^{}\beta _4^{}\beta _2^{}\beta _3^{}`$.
5.–7. $`c=0`$, $`\alpha _10=\alpha _1`$,
(c’) with $`s=q^2\alpha _1`$, $`u=q^4\alpha _1`$, $`t=v=0`$ and $`\beta _10\beta _2`$,
(a’) with $`x=q^2\alpha _1`$ and $`\frac{\beta _1}{\beta _2}=\frac{q^2}{(q^41)\alpha _0}`$, $`\frac{\beta _2}{\beta _1^{}}=\frac{\alpha _1}{(q^21)\alpha _0}`$,
(c’) with $`s=q^2\alpha _1`$, $`u=q^2\alpha _1`$, $`t=v=0`$ and $`\frac{\beta _1}{\beta _1^{}}=\frac{\alpha _1}{(q^21)\alpha _0}`$, $`\beta _20`$.
5.–7. $`c=0`$, $`\alpha _1=0\alpha _1`$,
(a) with $`x=q^2\alpha _1`$ and $`\frac{\beta _1}{\beta _2}=\frac{q^2}{(q^41)\alpha _0}`$, $`\frac{\beta _2}{\beta _1^{}}=\frac{\alpha _1}{(q^21)\alpha _0}`$,
(c’) with $`s=u=0`$, $`t=q^2\alpha _1`$, $`v=q^4\alpha _1`$ and $`\beta _10\beta _2`$,
(c’) with $`s=u=0`$, $`t=q^2\alpha _1`$, $`v=q^2\alpha _1`$ and $`\frac{\beta _1}{\beta _2^{}}=\frac{\alpha _1}{(q^21)\alpha _0}`$, $`\beta _20`$.
We have carried out the calculations for all embeddings of $`𝒪(\mathrm{S}_{qc}^2)`$ in $`𝒜`$. For each embedding, $`\epsilon `$ determines a classical point of $`𝒪(\mathrm{S}_{qc}^2)`$, i.e. an algebra homomorphism $`𝒪(\mathrm{S}_{qc}^2)𝐂`$; if $`cc(1)`$, this is a one-to-one correspondence. The equivariant derivations of $`_c`$ corresponding to the solutions 1–7 are given by $`\mathrm{d}ba=\chi (b_{(1)})b_{(2)}a`$. They are independent of the embedding and satisfy $`dim_{\epsilon ,\mathrm{r}}\mathrm{d}=2`$ for each classical point $`\epsilon `$, in the case of the solutions 1, 2 and 5–7 also $`dim_{\epsilon ,\mathrm{l}}\mathrm{d}=2`$; all this is proved in Section 4. The solutions for $`c=0`$, not being coboundaries if $`\alpha _1=0=\alpha _1`$, do not correspond to derivations with an $`\omega \mathrm{\Gamma }(\mathrm{d})`$ for which $`\mathrm{d}a=\omega aa\omega `$ for all $`a_c`$. If $`\mathrm{d}a=_i\mathrm{d}e_ia_i`$, then $`\chi (a)=\chi (_i\epsilon (a_i)e_i)`$. Since $`dim_𝐂\chi (\mathrm{Lin}_𝐂\{e_1,e_0,e_1\})=1`$ for the solutions 3, 5 and 6, but $`dim_𝐂\chi (_c)=2`$, these do not correspond to derivations for which $`\{\mathrm{d}e_1,\mathrm{d}e_0,\mathrm{d}e_1\}`$ generates $`\mathrm{\Gamma }(\mathrm{d})`$ as a right $`_c`$-module.
## 4. CONSTRUCTIONS
We retain the notations of Section 3. In the case of a quantum group, $`𝒜`$ is coquasitriangular, i.e. equipped with a $`𝐂`$-linear map $`r:𝒜𝒜𝐂`$ which satisfies
$$\begin{array}{c}r(a_{(1)}b_{(1)})a_{(2)}b_{(2)}=b_{(1)}a_{(1)}r(a_{(2)}b_{(2)}),\hfill \\ r(abc)=r(ac_{(1)})r(bc_{(2)}),r(1c)=\epsilon (c),\hfill \\ r(abc)=r(a_{(1)}c)r(a_{(2)}b),r(a1)=\epsilon (a)\hfill \end{array}$$
for all $`a,b,c𝒜`$. We use a construction method introduced in .
###### Lemma.
Let $`𝒜`$ be a coquasitriangular Hopf algebra and $``$ a right $`𝒜`$-comodule algebra. Let $`b_1,\mathrm{},b_N`$ be $`𝐂`$-linearly independent elements, $`\mathrm{\Delta }(b_i)=_jb_j\psi _i^j`$ and $`\nu `$ a comodule algebra endomorphism of $``$. Let $`\mathrm{\Gamma }`$ be the $``$-bimodule generated by the symbols $`\gamma ^1,\mathrm{},\gamma ^N`$ with the relations $`a\gamma ^j=_i\gamma ^i\nu (a_{(1)})r(\psi _i^ja_{(2)})`$, $`a`$, $`j=1,\mathrm{},N`$. Then $`\gamma ^1,\mathrm{},\gamma ^N`$ is a basis of $`\mathrm{\Gamma }`$ as a right $``$-module, and $`\mathrm{d}a:=\omega aa\omega `$ with $`\omega =_i\gamma ^ib_i`$ defines an equivariant derivation of $``$. Moreover, if $`\nu `$ is bijective, then $`\gamma ^ia=_j\nu ^1(a_{(1)})r(\psi _j^iS(a_{(2)}))\gamma ^j`$, and $`\gamma ^1,\mathrm{},\gamma ^N`$ is a basis of $`\mathrm{\Gamma }`$ as a left $``$-module, too. ∎
We set $`𝒜=𝒪(\mathrm{SL}_q(2))`$ and denote by $`u_j^i`$, $`i,j=1,\mathrm{\hspace{0.17em}2}`$, the canonical generators of $`𝒜`$. Then $`r`$ is defined by $`r(u_k^iu_l^j)=q^{1/2}R_{kl}^{ij}`$, where $`R`$ is the R-matrix of $`\mathrm{SL}_q(2)`$ specified in .
If $`c=q/(q+1)^2`$ and $`x_i\mathrm{Lin}_𝐂\{u_i^1,u_i^2\}`$, $`i=1,\mathrm{\hspace{0.17em}2}`$, with
$$\frac{\epsilon (x_1)}{\epsilon (x_2)}=\frac{q+1}{q}\frac{\alpha _1}{\alpha _0}\text{and}\epsilon (x_1)^2=\frac{q+1}{q(q1)}\frac{\alpha _1}{\alpha _0},$$
then the subalgebra $``$ of $`𝒜`$ generated by $`x_1`$, $`x_2`$ is generated by $`x_1`$, $`x_2`$ with the relation $`x_1x_2qx_2x_1=1`$, and it contains $`_c`$ as the subalgebra of the elements of even degree ($`x_1`$, $`x_2`$ being of degree 1), cf. . The lemma with $`N=2`$, $`b_i=x_i`$ and $`\nu =\mathrm{id}`$ yields an equivariant derivation $`\mathrm{d}`$ of $``$, and $`\mathrm{d}|__c`$ is an equivariant derivation of $`_c`$, which is by construction independent of the embedding of $`_c`$ in $`𝒜`$. We obtain
$$\mathrm{d}ab=_i\gamma ^i\left(x_ia_ja_{(1)}x_jr(u_i^ja_{(2)})\right)b$$
for all $`a,b`$. In particular, $`\omega =\gamma ^1x_1+\gamma ^2x_2`$,
$$\begin{array}{cc}\mathrm{d}e_1\hfill & =(q1)(\omega e_1+q^2\alpha _0\gamma ^1x_2),\hfill \\ \mathrm{d}e_0\hfill & =(q1)(\omega e_0q^2\alpha _0(\gamma ^1x_1q^2\gamma ^2x_2)),\hfill \\ \mathrm{d}e_1\hfill & =(q1)(\omega e_1+\alpha _0\gamma ^2x_1)\hfill \end{array}$$
and $`\omega =\frac{q^2(q+1)^2}{(q1)(q^31)\alpha _0^2}\left(\mathrm{d}e_1e_1+(q^2+1)^1\mathrm{d}e_0e_0+q^2\mathrm{d}e_1e_1\right)`$, thus $`\gamma ^ix_j\mathrm{\Gamma }(\mathrm{d}|__c)`$ for all $`i,j=1,\mathrm{\hspace{0.17em}2}`$. We set $`\gamma ^{in}=\gamma ^i(q)^{\delta _{n1}}x_{3n}`$ and calculate that
$$\mathrm{d}ab=_{in}\gamma ^{in}x_n\left(x_ia_ja_{(1)}x_jr(u_i^ja_{(2)})\right)b$$
for all $`a,b`$, thus $`\overline{\mathrm{d}ab}=\overline{_{in}\gamma ^{in}\epsilon (x_n)\chi _i(a)\epsilon (b)}`$ for all $`a,b_c`$, where $`\overline{\text{}}:\mathrm{\Gamma }(\mathrm{d}|__c)\mathrm{\Gamma }(\mathrm{d}|__c)/(\mathrm{\Gamma }(\mathrm{d}|__c)_c^+)`$ is the canonical projection and $`\chi _i(a)=\epsilon (x_i)\epsilon (a)r(x_ia)`$. Since $`dim_𝐂\mathrm{Lin}_𝐂\{\chi _1|__c,\chi _2|__c\}=2`$ and $`dim_𝐂\mathrm{Lin}_𝐂\{\overline{_n\gamma ^{in}\epsilon (x_n)}|i=1,\mathrm{\hspace{0.17em}2}\}=2`$, this implies $`dim_𝐂\overline{\mathrm{\Gamma }(\mathrm{d}|__c)}=2`$, that is, $`dim_{\epsilon ,\mathrm{r}}(\mathrm{d}|__c)=2`$. Our proof also shows that $`\{\mathrm{d}e_1,\mathrm{d}e_0,\mathrm{d}e_1\}`$ generates $`\mathrm{\Gamma }(\mathrm{d}|__c)`$ as a right $`_c`$-module. A similar argument starting with
$$\begin{array}{cc}\hfill a\mathrm{d}b& =_{ijkl}a\left(x_lb_{(1)}r(u_i^kS(b_{(2)}))bx_l\delta _{ik}\right)r(u_k^jS(u_j^l))\gamma ^i\hfill \\ & =_{ijk}a\left(_lx_lb_{(1)}r(S(u_j^l)b_{(2)})bx_j\right)r(u_i^kS(u_k^j))\gamma ^i\hfill \end{array}$$
shows $`dim_{\epsilon ,\mathrm{l}}(\mathrm{d}|__c)=2`$. For the definition of $`𝒪(\mathrm{S}_{qc}^2)`$ as a right $`𝒪(\mathrm{SO}_{q^2}(3))`$-comodule algebra only the square of $`q`$ is needed, therefore we can replace $`q`$ by $`q`$ and obtain the corresponding result for $`c=q/(q+1)^2`$. This establishes the claims about the solutions 1 and 2 in Section 3.
If $`c=c(2)`$ and $`x_i,y_i\mathrm{Lin}_𝐂\{u_i^1,u_i^2\}`$, $`i=1,\mathrm{\hspace{0.17em}2}`$, with either
$$\begin{array}{c}\text{(i)}\frac{\epsilon (x_1)}{\epsilon (x_2)}=\frac{(q^4+1)\alpha _1}{q^4\alpha _0}\text{and}\frac{\epsilon (y_1)}{\epsilon (y_2)}=\frac{(q^4+1)\alpha _1}{q\alpha _0}\text{or}\hfill \\ \text{(ii)}\frac{\epsilon (x_1)}{\epsilon (x_2)}=\frac{(q^4+1)\alpha _1}{\alpha _0}\text{and}\frac{\epsilon (y_1)}{\epsilon (y_2)}=\frac{(q^4+1)\alpha _1}{q^5\alpha _0},\hfill \end{array}$$
then $`y_ix_j_c`$ for all $`i,j=1,\mathrm{\hspace{0.17em}2}`$, and $`x_1y_2qx_2y_1=\zeta \mathrm{\hspace{0.17em}1}`$ with $`\zeta 𝐂\{0\}`$. Moreover, $`_mu_m^i\tau _{jm}(e_k)=_{mn}\tau _{mi}(e_n)u_j^m\pi _k^n`$ for the solution (f) in Section 3. We consider the $`_c`$-$`𝒜`$-bimodule with the right $`𝒜`$-module basis $`\gamma ^1`$, $`\gamma ^2`$ and the left $`_c`$-module operation $`a\gamma ^j=_i\tau _{ij}(a)\gamma ^i`$. Then $`\mathrm{d}a:=\omega aa\omega `$ with $`\omega =_i\gamma ^ix_i`$ defines an equivariant derivation of $`_c`$:
$$\mathrm{d}ab=_i\gamma ^i\left(x_ia_j\tau _{ij}(a)x_j\right)b$$
for all $`a,b_c`$. In particular, $`\omega =\gamma ^1x_1+\gamma ^2x_2`$,
$$\begin{array}{cc}\mathrm{d}e_1\hfill & =\omega e_1\frac{q^21}{q^4+1}\alpha _0\gamma ^1x_2,\hfill \\ \mathrm{d}e_0\hfill & =\omega e_0+\frac{q^21}{q^4+1}\alpha _0(\gamma ^1x_1q^2\gamma ^2x_2),\hfill \\ \mathrm{d}e_1\hfill & =\omega e_1q^2\frac{q^21}{q^4+1}\alpha _0\gamma ^2x_1.\hfill \end{array}$$
In addition, $`\omega =\frac{(q^4+1)^2}{(q^21)^2(q^2+1)\alpha _0^2}\left(\mathrm{d}e_1e_1+(q^2+1)^1\mathrm{d}e_0e_0+q^2\mathrm{d}e_1e_1\right)`$ in the case (ii), in which we see that $`\{\mathrm{d}e_1,\mathrm{d}e_0,\mathrm{d}e_1\}`$ generates $`\mathrm{\Gamma }(\mathrm{d})`$ as a right $`_c`$-module. Using $`\gamma ^{in}=\gamma ^i\zeta ^1(q)^{\delta _{n1}}x_{3n}`$, we can proceed like before to show $`dim_{\epsilon ,\mathrm{r}}\mathrm{d}=2`$. In the case (i), since the left action of $`\mathrm{Lin}_𝐂\{1,e_1,e_0,e_1\}`$ on $`\mathrm{Lin}_𝐂\{\gamma ^1,\gamma ^2\}`$ consists of all $`𝐂`$-linear endomorphisms, $`\{\mathrm{d}e_j,e_i\mathrm{d}e_j|i,j=1,\mathrm{\hspace{0.17em}0},\mathrm{\hspace{0.17em}1}\}`$ generates $`\mathrm{\Gamma }(\mathrm{d})`$ as a right $`_c`$-module. Exploiting e.g. the relations
$$\begin{array}{c}q^2e_1\mathrm{d}e_0=\left(e_0q^2\frac{q^21}{q^4+1}\alpha _0\right)\mathrm{d}e_1,q^2e_1\mathrm{d}e_0=\left(e_0+\frac{q^21}{q^4+1}\alpha _0\right)\mathrm{d}e_1,\\ \mathrm{d}e_1\left(q^2e_0\frac{q^61}{q^4+1}\alpha _0\right)=\mathrm{d}e_0e_1,\mathrm{d}e_1\left(e_0+\frac{q^61}{q^4+1}\alpha _0\right)=q^2\mathrm{d}e_0e_1,\\ (q^2+1)\mathrm{d}e_1e_1+\mathrm{d}e_0e_0+(q^2+1)\mathrm{d}e_1e_1=0\end{array}$$
we get $`dim_{\epsilon ,\mathrm{r}}\mathrm{d}2`$, while $`dim_{\epsilon ,\mathrm{r}}\mathrm{d}2`$ follows from Theorem 1 (iii). However, one checks that in both cases $`dim_{\epsilon ,\mathrm{l}}\mathrm{d}=0`$. We identify the cases (i) and (ii) with the solutions 3 and 4 in Section 3.
We take the basis $`\stackrel{~}{e}_2,\mathrm{},\stackrel{~}{e}_2`$ in of the spin 2 subcomodule of $`_c`$ for $`b_1,\mathrm{},b_N`$ and set $`\nu =\mathrm{id}`$ in the above lemma. This yields the equivariant derivation $`\mathrm{d}`$ of $`_c`$ which is given by $`\mathrm{d}ba=\stackrel{~}{\chi }(b_{(1)})b_{(2)}a`$ with $`\stackrel{~}{\chi }_i(b)=r(\stackrel{~}{e}_ib)\epsilon (\stackrel{~}{e}_i)\epsilon (b)`$ for all $`a,b_c`$, $`i=2,\mathrm{},\mathrm{\hspace{0.17em}2}`$. If $`c=0`$, then it corresponds to each case of the solution 5 in Section 3. To compute this, we write $`\stackrel{~}{\chi }_i`$ in terms of $`l(a)=r(u_1^1a)`$, $`l^1(a)=r(u_2^2a)`$ and $`l^{()}(a)=qr(u_1^2a)`$, while $`r(u_2^1a)=0`$, and use the relations $`ll^1=\epsilon =l^1l`$, $`ll^{()}=ql^{()}l`$ and $`(fg)|__c=0`$, where $`f`$ is defined in Section 3 and $`g`$ is any $`𝐂`$-linear functional on $`𝒜`$. Hence, we can restrict ourselves to $`\alpha _1=0=\alpha _1`$ to calculate
$$\mathrm{d}(e_1^3)=q^6(q^4+q^2+1)(q^2\mathrm{d}\stackrel{~}{e}_2e_1\mathrm{d}e_1\stackrel{~}{e}_2).$$
Since $`\mathrm{d}`$ is equivariant, this implies that for each $`a`$ in the spin 3 subcomodule of $`_c`$, $`\mathrm{d}a`$, and thus $`\mathrm{\Gamma }(\mathrm{d})`$, is contained in the right $`_c`$-module generated by $`\{\mathrm{d}e_1,\mathrm{d}e_0,\mathrm{d}e_1,\mathrm{d}\stackrel{~}{e}_2,\mathrm{},\mathrm{d}\stackrel{~}{e}_2\}`$. Exploiting e.g. the relations
$$\begin{array}{c}(q^2+1)\mathrm{d}e_1e_1+\mathrm{d}e_0e_0+(q^2+1)\mathrm{d}e_1e_1=0,\\ \mathrm{d}e_1(q^2e_0+\alpha _0)=\mathrm{d}e_0e_1,\\ (q^4+1)\mathrm{d}\stackrel{~}{e}_2(q^2e_0+\alpha _0)\mathrm{d}\stackrel{~}{e}_1e_1=q^4(q^2+1)\alpha _0\mathrm{d}e_1e_1,\\ \alpha _0^2\mathrm{d}e_1=\frac{q^4+1}{q^2}\mathrm{d}\stackrel{~}{e}_2e_1+\frac{1}{q^2(q^2+1)}\mathrm{d}\stackrel{~}{e}_1e_0\frac{1}{q^2(q^4+q^2+1)}\mathrm{d}\stackrel{~}{e}_0e_1,\\ \alpha _0^2\mathrm{d}e_0=\mathrm{d}\stackrel{~}{e}_1e_1+\frac{q^41}{q^61}\mathrm{d}\stackrel{~}{e}_0e_0+q^4\mathrm{d}\stackrel{~}{e}_1e_1,\\ \alpha _0^2\mathrm{d}e_1=\frac{q^4}{q^4+q^2+1}\mathrm{d}\stackrel{~}{e}_0e_1+\frac{1}{q^2+1}\mathrm{d}\stackrel{~}{e}_1e_0+\frac{q^4+1}{q^4}\mathrm{d}\stackrel{~}{e}_2e_1\end{array}$$
we get $`dim_{\epsilon ,\mathrm{r}}\mathrm{d}2`$, while $`dim_{\epsilon ,\mathrm{r}}\mathrm{d}2`$ follows from Theorem 1 (iii). Similarly, the derivation associated to the solution 6, say, $`\mathrm{d}^{}`$, arises if $`r^{}`$ with $`r^{}(ab)=r(S(b)a)`$, $`a,b𝒜`$, is used instead of $`r`$. With regard to the left $`_c`$-module structure of $`\mathrm{\Gamma }`$ the derivation $`\mathrm{d}`$ is given by $`a\mathrm{d}b=\stackrel{~}{\chi }^{}(b_{(1)})ab_{(2)}`$ with $`\stackrel{~}{\chi }_i^{}(b)=r(S(\stackrel{~}{e}_i)b)\epsilon (\stackrel{~}{e}_i)\epsilon (b)`$ for all $`a,b_c`$, $`i=2,\mathrm{},\mathrm{\hspace{0.17em}2}`$, thus $`\mathrm{d}^{}ba=\stackrel{~}{\chi }^{}(b_{(1)}^{})(b_{(2)}a)^{}`$. One checks that $`\mathrm{d}^{}=\mathrm{d}^{}`$, if $`e_i^{}=e_i`$, and therefore $`dim_{\epsilon ,\mathrm{l}}\mathrm{d}=dim_{\epsilon ,\mathrm{r}}\mathrm{d}^{}=2`$. Finally, the 2-dimensional covariant differential calculus described in corresponds to the solution 7. It is equal to $`\mathrm{d}_{(_c^+)^2}`$, if $`\alpha _1=0=\alpha _1`$.
The classification problem in Section 3 with $`dim_{\epsilon ,\mathrm{l}}\mathrm{d}=2`$ instead of $`dim_{\epsilon ,\mathrm{r}}\mathrm{d}=2`$ is equivalent, because the Hopf algebra $`𝒪(\mathrm{SL}_q(2))`$ with the opposite multiplication is isomorphic to $`𝒪(\mathrm{SL}_{q^1}(2))`$ and correspondingly the right comodule algebra $`𝒪(\mathrm{S}_{qc}^2)`$ with the opposite multiplication to $`𝒪(\mathrm{S}_{q^1c}^2)`$. It is still open for $`c=c(0),c(1),\mathrm{}`$ and for $`\rho =\lambda =0`$. Since in the one-to-one correspondence in Theorem 1 (iii) only the trivial derivation $`\mathrm{d}=\mathrm{d}_^+`$ satisfies $`dim_{\epsilon ,\mathrm{l}}\mathrm{d}=0`$, the equivalents with the opposite comultiplication of the equivariant derivations for $`c=c(2)`$ constructed above do not occur, i.e. they do not arise from right ideals.
## ACKNOWLEDGMENTS
I am grateful to Prof. K. Schmüdgen for showing his interest in my work and to S. Kolb for detailed discussions. This work was supported by the Deutsche Forschungsgemeinschaft within the scope of the postgraduate scholarship programme “Graduiertenkolleg Quantenfeldtheorie” at the University of Leipzig.
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# Probing the gas content of the dwarf galaxy NGC 3109 with background X-ray sources
## 1 Introduction
NGC 3109 (DDO 236), a late-type dwarf spiral galaxy, has been classified as Sm IV (Sandage & Tammann 1981) and is seen almost edge-on with an inclination close to 80 (Carignan 1985). Since this galaxy has no nucleus it has been assigned the morphological type Ir (van den Bergh 1999). It is a southern Magellanic dwarf galaxy well resolved into stars and is one of the largest Magellanic dwarfs close to the Local Group at a distance of 1.36$`\pm `$0.10 Mpc (Musella et al. 1997) with an apparent major-axis diameter of 30′ (12 kpc) and minor-axis diameter of 6′ (Demers et al. 1985).
NGC 3109 has a similar dimension as the LMC. With an absolute magnitude comparable to the SMC of $`M_\mathrm{B}=15.7`$ mag it is rather underluminous. Estimates of the total mass run from 0.6 to $`\mathrm{1.6\; 10}^{10}\mathrm{M}_{}`$. The galaxy is surrounded by a huge Hi envelope, quite larger than its optical size (Materne 1980; Huchtmeier et al. 1980; Jobin & Carignan 1990). Jobin & Carignan (1990) derive an Hi mass of $`5\pm 1\times 10^8\mathrm{M}_{}`$. The rotation curve requires a large dark matter halo.
On deep exposures NGC 3109 exhibits spiral structure. Associations of knots of stars are visible along spiral arms. A globular cluster search revealed ten candidates. Bright blue stars are found evenly distributed over the face of the galaxy but particularly along the spiral arms, indicating that star formation is taking place on a galaxy-wide scale (Demers et al. 1985). The metallicity of NGC 3109 has been found to be low, similar to the SMC (Richer & McCall 1995). Minniti et al. (1999) establish the existence of an extended halo of old and metal-poor stars.
NGC 3109 belongs as most luminous member to a subgroup of the Local Group dwarfs which is relatively isolated (Mateo 1998). Another member of this group is the dwarf spheroidal galaxy Antila at a distance of only 1.2 from NGC 3109. From this apparent separation on the sky the lower limit to the distance between both galaxies is $``$26 kpc (Whiting et al. 1997).
In this article we derive the X-ray population in the field of NGC 3109 from ROSAT PSPC observations. We classify a few of the X-ray sources which coincide with the Hi extent of NGC 3109 from their X-ray spectral properties. We find two candidate AGN which we use to probe the gas content of NGC 3109.
## 2 Observations
The observations used were carried out with the PSPC detector of the ROSAT observatory during two pointed observations in May-June and November 1992. The satellite, X-ray telescope (XRT) and the focal plane detector (PSPC) are described in detail in Trümper (1983) and Pfeffermann et al. (1986). We note that no observations have been performed in the direction of NGC 3109 with the HRI detector of ROSAT.
The data of the two pointings given in Table 1 have been retrieved from the public ROSAT archive. The data sets from the two observations have been merged to one data set using standard EXSAS procedures (Zimmermann et al. 1994).
## 3 The catalog of X-ray sources
Three detection procedures (local, map, and maximum likelihood) were applied to the merged pointings using EXSAS commands (Zimmermann et al. 1994). In the local and map detection a square shaped detection cell is slid over the image. Source counts are determined within the cell, background counts either from an area surrounding the cell (local background) or from a (smoothed) background image within the same cell. In the maximum likelihood detection the distribution of the detected photons above the background is compared in a maximum likelihood ratio test with the analytical point-spread function.
The analysis was performed in the five energy channel ranges Soft = (channel 11-41, 0.1-0.4 keV), Hard = (channel 52-201, 0.5-2.1 keV), Hard1 = (channel 52-90, 0.5-0.9 keV) and Hard2 = (channel 91-201, 0.9-2.0 keV) and broad (0.1-2.4 keV). The five source lists were merged to one final source list taking detections at off-axis angles $``$50′ into account. The maximum likelihood algorithm was used to determine the final source position, the counts in five energy bands and the source extent. A one-dimensional energy and position dependent Gaussian distribution was applied in order to obtain the source extent. The source extent ($`Ext`$) is given as the Gaussian $`\sigma _{\mathrm{Gauss}}`$
$$Ext=\sigma _{\mathrm{Gauss}}=FWHM_{\mathrm{Gauss}}/2.35$$
(1)
Hardness ratios $`HR1`$ and $`HR2`$ were calculated from the counts in the bands as $`HR1=(HS)/(H+S)`$ and $`HR2=(H2H1)/(H1+H2)`$. The existence likelihood ratio and the extent likelihood ratio was calculated according to Cash (1979) and Cruddace et al. (1988). We selected for our final source catalog only detections with an existence likelihood ratio $`LH_{\mathrm{exist}}10`$, which is equal to a probability of existence $`P(1\mathrm{exp}(LH_{\mathrm{exist}}))(14.5\times 10^5)`$. We give the value for the extent only in case the extent likelihood ratio is $`LH_{\mathrm{ext}}20`$. A $`90\%`$ source error radius was calculated, adding quadratically a systematic error of $`5`$″ (cf. Kürster 1993).
$$P_e=2.1\times \sqrt{x_{\mathrm{err}}^2+y_{\mathrm{err}}^2+(5\mathrm{})^2}$$
(2)
The positional error derived for large off-axis angles $`\mathrm{\Delta }{}_{}{}^{}{}_{}{}^{>}30^{}`$ may be somewhat underestimated due to the asymmetry of the point-spread-function. But the positional error should not be larger than $``$1′.
We finally screened the catalog of ROSAT sources by displaying the positions of these sources on the hard, soft, and broad band ROSAT PSPC image. We could confirm 91 of the detected sources. This screened catalog of pointlike and moderately extended sources is given in Table 6. <sup>1</sup><sup>1</sup>1Table 6 will be made available on-line with the electronically published version. We give in the first column of the catalog for each confirmed source the sequence number and in the second column the source number (from the catalog of unscreened detections). We always refer to the source number in the text.
In Fig. 1 we show the ROSAT PSPC image of the NGC 3109 field in the hard band (0.5 – 2.0 keV). We note that the dwarf spheroidal galaxy Antila which is $``$1.2 to the south of NGC 3109 is outside of the observed field of view. In Fig. 2 we mark on the ROSAT PSPC image of the central 20′ of NGC 3109 the sources for with accurate hardness ratios $`\delta HR2`$0.2 have been determined.
## 4 Classification of the X-ray sources
The X-ray colors (hardness ratios $`HR1`$ and $`HR2`$) as given in Table 6 can be used for a source classification. Kahabka et al. (1999, hereafter KPFH99) have made a classification of the ROSAT PSPC X-ray sources in the field of the Small Magellanic Cloud (SMC). X-ray binaries have on average harder X-ray spectra than supernova remnants and background AGN. The range of hardness ratios $`HR1`$ and $`HR2`$ is then different. In addition absorption due to intervening cold gas affects the X-ray colors. KPFH99 have derived the dependence of $`HR1`$ and $`HR2`$ on the absorbing column density for a low metallicity galaxy like the SMC. The hydrogen column density in the direction of the X-ray source has been derived from the high-resolution 21-cm map of Stanimirovic et al. (1999). Then from the measured X-ray color $`HR2`$ a classification could be obtained. The strength of this method is that X-ray binaries and AGN cover different parts of the diagram for intervening hydrogen columns $`{}_{}{}^{}{}_{}{}^{<}6\times 10^{21}\mathrm{cm}^2`$. AGN are located in a band in this diagram with radio loud AGN populating the upper regime of the band while radio quiet AGN occupy the lower part of the band (cf. Laor et al. 1997; Brinkmann et al. 1997).
For radio loud AGN spectra (powerlaw photon index $`2.0`$) the dependence of the X-ray colors $`HR1`$ and $`HR2`$ on the hydrogen column density $`N_\mathrm{H}`$ ($`10^{21}\mathrm{cm}^2`$) for abundances $``$0.2 solar (SMC abundances) has been derived from simulations as
$$HR1=1.000.82\left(\frac{N_\mathrm{H}}{0.25}\right)^{1.35}(N_\mathrm{H}0.3)$$
(3)
$$HR2=0.145+0.65\left(\frac{N_\mathrm{H}}{15}\right)^{0.68}(N_\mathrm{H}0.3)$$
(4)
and for radio quiet AGN spectra (powerlaw photon index $`2.6`$)
$$HR1=1.001.6\left(\frac{N_\mathrm{H}}{0.25}\right)^{1.52}(N_\mathrm{H}0.3)$$
(5)
$$HR2=0.01+0.84\left(\frac{N_\mathrm{H}}{15}\right)^{0.82}(N_\mathrm{H}0.3)$$
(6)
For X-ray binaries with powerlaw photon index $`0.8`$ spectra the dependence of the X-ray colors $`HR1`$ and $`HR2`$ on the hydrogen column density is
$$HR1=1.000.20\left(\frac{N_\mathrm{H}}{0.25}\right)^{0.94}(N_\mathrm{H}0.3)$$
(7)
$$HR2=0.415+0.48\left(\frac{N_\mathrm{H}}{15}\right)^{0.72}(N_\mathrm{H}0.3)$$
(8)
Equations 3 to 8 can be solved for the hydrogen column density $`N_\mathrm{H}`$. In combining the equations for $`HR1`$ and $`HR2`$ one can derive analytical solutions of X-ray binary and AGN tracks in the $`HR1`$$`HR2`$ plane.
For radio loud AGN (powerlaw photon index $`2.0`$) one derives the track
$$HR2=0.145+0.0402\left(\frac{1.0HR1}{0.82}\right)^{0.504}(N_\mathrm{H}0.3)$$
(9)
and for radio quiet AGN (powerlaw photon index $`2.6`$) the track
$$HR2=0.01+0.0292\left(\frac{1.0HR1}{1.6}\right)^{0.539}(N_\mathrm{H}0.3)$$
(10)
For X-ray binaries (powerlaw photon index $`0.8`$) we derive the track
$$HR2=0.415+0.0252\left(\frac{1.0HR1}{0.16}\right)^{0.766}(N_\mathrm{H}0.3)$$
(11)
These tracks allow a classification of a ROSAT source as an AGN (or an X-ray binary) without the requirement that the intervening hydrogen column density is known. It even is possible to constrain in the $`HR1`$$`HR2`$ plane the intervening hydrogen column density for the source. We note that X-ray binaries are not necessarily seen through the total gas column of the galaxy disk. With a galactic foreground column of $`4\times 10^{20}\mathrm{cm}^2`$ towards NGC 3109 we derive with Equation 11 for X-ray binaries a lower bound on $`HR2`$ of $`HR20.45`$.
The tracks given above have been derived for a metallicity $``$0.2 solar which is consistent with the metallicity derived for Hii regions in NGC 3109 (cf. Minniti et al. 1999). Higher metallicity tracks deviate somewhat from these tracks especially if one approaches the high column density regime $`HR11.0`$.
## 5 The Hi extent of NGC 3109
Huchtmeier et al. (1980, hereafter HSM80) measured the Hi distribution across NGC 3109 with the Effelsberg radio telescope and a beam size of 9′. We apply the procedure used for the X-ray detections in the field of the SMC to the X-ray detections in the field of NGC 3109. First we make use of the Hi map of HSM80 which has an extent of 60′ $`\times `$ 30′. We convert the Hi intensity into hydrogen column densities with the equation given in Dickey & Lockman (1990), cf. Kahabka (1999) and we derive a peak of $`N_\mathrm{H}1.3\times 10^{21}\mathrm{cm}^2`$. HSM80 have found that NGC 3109 has a large extent in Hi, with an extension (distortion) of the Hi in the SW. With our X-ray catalog we cover the whole extent of the Hi of NGC 3109. If we restrict the analysis to the inner 20′ of this field and to sources with well constrained hardness ratios $`\delta HR2`$0.2 then we can classify 7 sources as AGN.
A much higher resolution Hi map of NGC 3109 with a beam size of 40″ has been derived by Jobin & Carignan (1990, hereafter JC90) with the VLA. We now make use of this high-resolution Hi image which we take from plate 67 of JC90. The Hi distribution of NGC 3109 has an extent of 40′ $`\times `$ 12′. The peak hydrogen column density is $`2.3\times 10^{21}\mathrm{cm}^2`$, and the lowest column density is $`10^{19}\mathrm{cm}^2`$.
In Fig. 3 we show the positions of the cataloged ROSAT PSPC sources overlaid on the gray scale Hi image of NGC 3109 and taken from JC90. We find 26 ROSAT PSPC sources within the Hi contours of JC90. If we restrict the analysis to the inner 20′ of this field and to sources with well constrained hardness ratios $`\delta HR2`$0.2, then we can classify 3 sources (number 36, 41 and 63) as AGN and 2 sources (number 53 and 73) as X-ray binaries. Source 53 may be seen through higher gas columns while source 73 is seen through lower gas columns (see Fig. 4 and Equation 11). One source (number 82) can be either class. The two sources 52 and 92 cannot be classified as AGN or X-ray binaries. Source 92 may be a foreground object (cf. Fig. 4, Fig. 5 and Table 2). Source 52 has similar hardness ratios as the LMC SNR 0548-70.4 (Haberl & Pietsch 1999). It could be a (young) SNR in NGC 3109.
NGC 3109 has a mass smaller than, or comparable to, the LMC and 2 X-ray binaries with luminosities above a few times $`10^{36}\mathrm{erg}\mathrm{s}^1`$ would be in agreement with extrapolations from ROSAT findings for the LMC (cf. Haberl & Pietsch 1999).
There are 14 ROSAT PSPC sources which are projected onto NGC 3109 intrinsic hydrogen columns of $`N_\mathrm{H}10^{21}\mathrm{cm}^2`$ (cf. Fig. 3). The largest column is derived for source 86 ($`N_\mathrm{H}=\mathrm{2.3\; 10}^{21}\mathrm{cm}^2`$). This source could be associated with a spiral arm or an Hii region of NGC 3109. Two further sources, 53 and 59, are close to another region of large column density and source 59 is also close to the optical center of NGC 3109.
We further test the correctness of the classification by constructing the distribution of the number of detected X-ray sources $`N`$ with fluxes in excess of a given flux $`S`$, the $`\mathrm{log}N\mathrm{log}S`$ (cf. Hasinger et al. 1993). We restrict the analysis to sources within the central 20′ of NGC 3109 which have been classified as AGN or X-ray binaries. We correct the X-ray fluxes for the intervening hydrogen columns by using for the galactic contribution a value of $`\mathrm{4.3\; 10}^{20}\mathrm{cm}^2`$ (Dickey & Lockman 1990) and for the contribution due to NGC 3109 the Hi model of JC90. We find that the observed $`\mathrm{log}N\mathrm{log}S`$ is in excess of the $`\mathrm{log}N\mathrm{log}S`$ of the soft extragalactic X-ray background (Hasinger et al. 1993). This can be accounted for if an additional less steep component (e.g. due to X-ray binaries) is added. If we only construct the $`\mathrm{log}N\mathrm{log}S`$ of the candidate AGN, then the $`\mathrm{log}N\mathrm{log}S`$ of the soft extragalactic X-ray background is reproduced. This would mean that no significant additional hydrogen is needed to explain the observed sources as background AGN and X-ray binaries although only a few candidate AGN are within the Hi extent of NGC 3109. The two brightest candidate AGN (source 36 and 63) could give an excess in the $`\mathrm{log}N\mathrm{log}S`$. We cannot exclude that they are at least in part related to NGC 3109.
## 6 Optical counterparts of the X-ray sources
We generated finding charts in the B-band COSMOS blue plates and in the I-band (NTT EMMI, cf. Fig. 6) to search for optical counterparts of the ROSAT sources. The B-band plates were used to classify the source using the flux ratio $`f_x/f_{\mathrm{opt}}`$ while the I-band image was used to determine the source morphologies from an analysis of their point spread function (see below).
### 6.1 Matches in the B-band
We produced finding charts of the ROSAT sources given in Table 6 using the COSMOS blue plates. We especially investigated the classified ROSAT sources given in Table 2.
We searched for the optical counterparts in the 90% confidence circle of the ROSAT source. We list the parameters of the matches in Table 3. For the two sources 28 and 92 we find bright (B=11.7 mag and B=12.1 mag) optical counterparts in the 90% confidence circle of these ROSAT sources. This indicates a galactic foreground nature of these sources. For the sources 36, 41, 52, 68 and 73 we find weak (B=20.4 mag to B=22.3 mag) unresolved optical counterparts. For source 73 there is also a galaxy (B=22.7 mag) in the error circle. For source 53 we cannot find an optical counterpart as it is in the crowded field of NGC 3109 and for 63 no optical counterpart exists within the ROSAT 90% confidence circle.
In addition we calculated the flux ratio $`f_\mathrm{x}/f_{\mathrm{opt}}`$ with the equation given in Haberl & Pietsch (1999). We make use of the B-magnitude given in Table 3 and the PSPC countrate given in Table 6. We use the abbreviation $`k_{\mathrm{PSPC}}=\mathrm{log}(10^{11}\mathrm{𝑃𝑆𝑃𝐶}\mathrm{countrate})`$ and obtain the flux ratio from the equation:
$$\mathrm{log}(f_\mathrm{x}/f_{\mathrm{opt}})=k_{\mathrm{PSPC}}+0.4m_\mathrm{B}+5.37$$
(12)
Using the scheme given in Table 3 of Haberl & Pietsch (1999) we classify the sources 28 and 92 as candidate foreground stars and the sources 36 and 41 as candidate AGN (cf. Table 3).
### 6.2 Matches in the I-band
For the first run of optical matching of X-ray sources in the central part of NGC 3109 we obtained three optical I-band images (60.E-0818) from the ESO archive.<sup>2</sup><sup>2</sup>2Based on observations made with ESO Telescopes at the La Silla or Paranal Observatories under programme ID 60.E-0818(A). All images were taken on 1998 February 2 with the red arm of the ESO Multi-Mode Instrument (EMMI) at the New Technology Telescope (NTT) with a total exposure time of 2700 s. The images were taken with the Tektronix 2048$`\times `$2048 pix<sup>2</sup> chip with a scale of 0.268″/pix. The resulting field of view covered by the stacked images is 8.8′ $`\times `$ 8.8′ and the final image is centered on the galaxy nucleus. The seeing was measured to be $`0.75`$″ throughout all three exposures.
The data are used to identify optical couterparts within the region of highest Hi column density, which resides at the center of NGC 3109 (Jobin & Carignan 1990) and is best covered by the field of view of this particular dataset, and to support or reject the classification provided by the hardness-ratio estimate (see Sect. 5). Further photometric studies in optical passbands of ROSAT sources for most objects in the catalog in Table 6 with accurate positions will be given in a subsequent paper.
In order to obtain bona-fide optical counterparts we transformed the pixel coordinates of the I-band to new equatorial coordinates. We used positions of 43 stars out of the USNO Astrometric Catalog (Monet et al. 1998) which was obtained from Centre de Données de Astronomiques de Strasbourg (CDS). Using the task ccmap within the IRAF<sup>3</sup><sup>3</sup>3IRAF is distributed by the National Optical Astronomy Observatories, which are operated by the Association of Universities for Research in Astronomy, Inc., under cooperative agreement with the National Science Foundation. environment (Tody et al. 1993) we find the mean rms uncertainty for all given optical coordinates being $`\mathrm{\Delta }`$RA$`0.62`$″ and $`\mathrm{\Delta }`$Dec$`0.48`$″. Applying the plate solution to coordinates provided by the ROSAT PSPC catalog (see Table 6) we find 7 optical matches within the field of view. The finding charts for all these regions are given in Fig. 6 with the appropriate scaling of the 90%-confidence circle as given in Table 6. The five brightest objects are encircled and labeled according to their luminosity (label 1 marks the brightest object).
Subsequently, we perform an analysis of the shape parameters of the optical point spread function (PSF) for the 5 brightest objects within each of the 7 object’s finding circles. We used the source-extraction software SExtractor v2.1.4 (Bertin & Arnouts 1996) for determination of ellipticity $`\epsilon `$, full width at half maximum (FWHM), and the star-galaxy-classification parameter CL (using a neural-network algorithm which was extensively trained, see Bertin & Arnouts for details). We compare our findings to all remaining optical objects most of which are assumed to be point-sources. The code yields $``$7500 detections for which we plot the three PSF-shape parameters ($`\epsilon `$, FWHM, and CL) as a function of instrumental I-band magnitude. Fig. 7 shows the parameters of all detections together with the parameters for the brightest source within the finding circle (see Table 4). Note that the brightest object is not necessarily the optical counterpart of the ROSAT source.
The brightest sources 52.1, 62.1 and 65.1 (the first number being the source number and the decimal the source label) in Fig. 6 are at least one magnitude brighter than the next fainter object within the finding circle (see Table 4). Object 63.1 has about the same magnitude as object 63.2. Therefore no clear preference can be given just from considering the I-band flux. The PSF-shape parameters of these 4 objects agree with those found for stars (see Fig. 7). The CL classification of SExtractor puts all 4 sources in the point-source regime. Yet, the optical findings are prone to mis-classification since not every brightest object is most central within the finding circle.
Although they are not the brightest, the sources 62.5 and 63.2 are the most likely candidate optical counterparts (infering from their most central position only). These were classified as slightly extended and slightly elliptical, respectively.
As can be seen from Fig. 6 there are 2 finding charts (i.e. charts 62 and 65) from the edge of the CCD image. Thus, we cannot exclude brighter objects within the area not covered by the exposure. Chart 59 is too crowded to give a reliable identification only on the basis of any photometric PSF-shape parameter. However, we note that the center of the ROSAT error circle coincides with the galaxy’s center. Thus, it is very likely that the ROSAT source is connected with a X-ray source which is located near the center of the galaxy.
## 7 Candidate AGN
In Sect.5 we have used the hardness ratios $`HR1`$ and $`HR2`$ to classify 3 ROSAT sources (with numbers 36, 41 and 63) and projected onto the Hi extent of NGC 3109 as candidate AGN.
We have searched in the optical B-band COSMOS finding charts of these sources for optical counterparts. In addition we have generated deep I-band finding charts for source 63 and other ROSAT sources and we have looked for faint unresolved and extended objects within the positional error circles of these ROSAT sources (cf. Fig. 6). We have found several optical candidates in the positional error circle of each ROSAT source. Therefore optical spectroscopy of all candidates is required to firmly identify the optical counterpart of these ROSAT sources.
If the background nature of candidate AGN “close to” NGC 3109 is established then these AGN can be used to probe the gas column density of NGC 3109 in the direction of these AGN (taking the additional galactic gas column into account). For two sufficiently X-ray bright candidate AGN RX J1003.2-2607 (source 63) and RX J1002.5-2602 (source 36) we perform X-ray spectral fitting with the EXSAS spectral analysis task (Zimmermann et al. 1994). We assume a galactic foreground column density of $`\mathrm{4.3\; 10}^{20}\mathrm{cm}^2`$ and determine the additional hydrogen column density (due to NGC 3109) assuming reduced metallicities ($``$0.2 solar). For both sources we find that the hydrogen column deduced from the X-ray spectral fit is consistent with the hydrogen column deduced from the 21-cm survey of HSM80 and JC90 (cf. Table 5 and Fig. 8).
This result is also consistent with the location of these two sources in the $`HR1`$$`HR2`$ plane (cf. Fig. 5). We do not include in the spectral analysis source 53 which has been classified as a candidate X-ray binary in Sect.5.
## 8 Molecular content of NGC 3109
Assuming that both AGN have canonical powerlaw photon indices $`\alpha `$ = $``$(1.75 to 2.25), cf. Laor et al. (1997), we further tighten the constraints on the hydrogen column density. For RX J1003.2-2607 we find from the X-ray spectral fit $`N_\mathrm{H}=11\pm _5^7\times 10^{20}\mathrm{cm}^2`$ and from the 21-map of JC90 $`N_\mathrm{H}=8\pm 2\times 10^{20}\mathrm{cm}^2`$. For RX J1002.5-2602 we constrain the hydrogen column density from the X-ray spectral fit to $`N_\mathrm{H}=4\pm _4^6\times 10^{20}\mathrm{cm}^2`$ and from the 21-map of JC90 to $`N_\mathrm{H}1\times 10^{20}\mathrm{cm}^2`$. Both values are in agreement within the uncertainties.
The column densities of absorbing material from the X-ray data are slightly larger than those from the 21-cm line, but the uncertainty of the X-ray spectral fit leaves room for some additional hydrogen (e.g. molecular hydrogen) in the line of sight of the AGN. We assume that the photoionisation cross section is $`2.8`$ times larger for molecular hydrogen than for atomic hydrogen (cf. Cruddace et al. 1974; Yan et al. 1998). Then we determine the column density due to molecular hydrogen $`N_{\mathrm{H}_2}`$ from the total hydrogen column density $`N_\mathrm{H}^{\mathrm{tot}}`$ derived from the X-ray spectral fit and the atomic hydrogen column density $`N_{\mathrm{H}\mathrm{i}}`$ derived from the 21-cm observations as
$$N_{\mathrm{H}_2}=\frac{1}{2.8}\left(N_\mathrm{H}^{\mathrm{tot}}N_{\mathrm{H}\mathrm{i}}\right)$$
(13)
We so constrain the gas column due to molecular hydrogen in NGC 3109 from the X-ray spectral fit of RX J1003.2-2607 to $`{}_{}{}^{}{}_{}{}^{<}4\times 10^{20}\mathrm{cm}^2`$.
Towards RX J1003.2-2607 the amount of $`\mathrm{H}_2`$ is $``$$`10^{20}\mathrm{cm}^2`$. Along this line of sight the molecular mass fraction, $`f_\mathrm{m}=(N_\mathrm{H}^{\mathrm{tot}}N_{\mathrm{H}\mathrm{i}})/(N_\mathrm{H}^{\mathrm{tot}}+0.4N_{\mathrm{H}\mathrm{i}})=0.21\pm 0.37`$. This means $`{}_{}{}^{}{}_{}{}^{<}60`$% of the mass of the total gas is in molecular form. This result can be compared with CO observations of NGC 3109, from which the mass of molecular hydrogen has been determined to $`{}_{}{}^{}{}_{}{}^{>}`$$`\mathrm{4\; 10}^7M_{}`$ (Rowan-Robinson, Philips & White 1980). With a Hi mass of $`\mathrm{5\; 10}^8M_{}`$ a molecular mass fraction of $`{}_{}{}^{}{}_{}{}^{>}`$10% is obtained.
Extraplanar absorbing clouds have been found in a high-resolution survey of a sample of 12 edge-on galaxies by Howk & Savage (1999) at distances of $`z=0.5`$ to $`1.5`$ kpc from the galaxy plane. For an inclination of 80 a radius of the galaxy disk of 12 kpc and a distance to the galaxy of 1.4 Mpc similar clouds would be projected in NGC 3109 1′ to 4′ from the galaxy plane. From these clouds dust absorption has been observed but the dust should coexist with gas in the molecular and atomic phase. Similar clouds may exist in NGC 3109 and the candidate AGN RX J1003.2-2607 would be seen through gas which is at a height $`z`$1 kpc above the galaxy plane while the candidate AGN RX J1002.5-2602 is seen through gas with $`z`$2.4 kpc.
## 9 Summary and conclusions
From ROSAT PSPC observations of the dwarf galaxy NGC 3109 we derive 10 X-ray sources which are contained within the Hi extent of this galaxy and which have accurate hardness ratios. We classify 2 of these sources as foreground stars, 2 as candidate X-ray binaries and 3 as candidate background AGN. From X-ray spectral fitting we derive for 2 of the AGN total hydrogen column densities which we compare with the Hi column densities inferred from 21-cm line measurements. We estimate that the molecular mass fraction of the gas is not larger than $``$60%. Upcoming spectroscopy of candidate optical counterparts will help to constrain the nature of these objects and help to understand the hydrogen content of NGC 3109. Furthermore, spectroscopy of globular clusters around NGC 3109 will help to trace the enrichment history (i.e. the major star formation episodes) and the kinematics (i.e. total mass) of NGC 3109.
###### Acknowledgements.
The ROSAT project is supported by the Max-Planck-Gesellschaft and the Bundesministerium für Forschung und Technologie (BMFT). This research made use of the COSMOS digitized optical survey of the southern sky, operated by the Royal Observatory Edinburgh and the Naval Research Laboratory, with support from NASA. This research has made use of the SIMBAD data base operated at CDS, Strasbourg, France. We thank K.S. de Boer and U. Klein for critically reading the manuscript. We thank an anonymous referee for useful comments.
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# References
ITP-UH-06/00
hep-th/0005044
New N=2 Supersymmetric Vector-Tensor Interaction
Ulrich Theis
Institut für Theoretische Physik, Universität Hannover,
Appelstraße 2, 30167 Hannover, Germany
utheis@itp.uni-hannover.de
Abstract
An $`N=2`$ supersymmetric self-interaction of the vector-tensor multiplet is presented, in which the vector provides the gauge field for local central charge transformations. The dual description in terms of a vector multiplet and an $`N=1`$ superspace formulation are given.
Since the discovery of the relevance of the $`N=2`$ vector-tensor multiplet to certain heterotic string theory compactifications , its possible interactions have been investigated in a number of publications . In the central charge of the multiplet was gauged by a coupling to an abelian vector multiplet that provides the gauge field for local central charge transformations, while in a self-interaction of the multiplet was introduced, giving rise to a Chern-Simons coupling of the vector and tensor. These two kinds of interactions have subsequently been rederived in harmonic superspace from deformations of the superfield constraints that determine the multiplet .
In the present letter it is shown how the central charge of the vector-tensor multiplet can be gauged by a self-interaction, without a coupling to an additional vector multiplet. The model is obtained from a deformation of the free component action and its corresponding supersymmetry and gauge transformations by means of the Noether method. This results in a realization of the supersymmetry algebra which holds only on-shell, i.e. on fields satisfying their equations of motion. We shall address the problem of constructing a manifestly $`N=2`$ supersymmetric off-shell formulation in a future publication and content ourselves with an off-shell version in terms of $`N=1`$ superfields for the time being.
Gauge field theories of the kind considered here have been investigated in detail in in the framework of $`N=1`$ supersymmetry. They are actually members of a family of (bosonic) models discovered by Henneaux and Knaepen in , that in four spacetime dimensions involve non-polynomial interactions of 1- and 2-forms. In it was shown that every four-dimensional Henneaux-Knaepen model admits an $`N=1`$ supersymmetric generalization. The vector-tensor multiplet with gauged central charge as found in was so far the only known model with two supersymmetries. In the following we present the second example.
The vector-tensor multiplet contains a 1-form $`A_\mu `$, a 2-form $`B_{\mu \nu }`$, a real scalar $`\varphi `$, and an SU(2) doublet of Weyl spinors $`\psi _\alpha ^i`$. The bosonic and fermionic off-shell degrees of freedom can be matched by inclusion of a real auxiliary scalar $`U`$.
The supersymmetry algebra of the free multiplet involves a central charge that acts non-trivially on the gauge fields,
$$\mathrm{\Delta }_cA_\mu =cH_\mu ,\mathrm{\Delta }_cB_{\mu \nu }=cF_{\mu \nu },c.$$
(1)
Here $`F^{\mu \nu }=\epsilon ^{\mu \nu \rho \sigma }_\rho A_\sigma `$ and $`H^\mu =\frac{1}{2}\epsilon ^{\mu \nu \rho \sigma }_\nu B_{\rho \sigma }`$ are the Hodge-dual field strengths of $`A_\mu `$ and $`B_{\mu \nu }`$ respectively. The conserved current that corresponds to this global symmetry of the free action is simply
$$J^\mu =F^{\mu \nu }H_\nu .$$
(2)
A coupling of $`A_\mu `$ to this current yields a consistent interaction vertex $`gA_\mu J^\mu `$ of dimension five (where accordingly $`g`$ is a coupling constant of dimension $`1`$). It requires to modify the gauge transformations of $`A_\mu `$ and $`B_{\mu \nu }`$ by first-order terms that are precisely the central charge transformations given above,
$`\delta A_\mu `$ $`=_\mu C(x)+gC(x)H_\mu +O(g^2)`$
$`\delta B_{\mu \nu }`$ $`=_\mu C_\nu (x)_\nu C_\mu (x)+gC(x)F_{\mu \nu }+O(g^2),`$ (3)
which amounts to gauging the global symmetry. In order to restore supersymmetry to order $`g`$, the vertex $`A_\mu J^\mu `$ is to be accompanied by additional terms of dimension five, which however depend on the gauge potentials only via the field strengths and therefore do not introduce further first-order deformations of the gauge transformations. The question now is whether the action and the transformations can be extended such that both gauge invariance and $`N=2`$ supersymmetry are realized to all orders in the coupling constant.
Since the interaction falls into the family of Henneaux-Knaepen models, it is to be expected that the resulting action and transformations are non-polynomial in the fields, so a Noether construction order by order seems impractical. As it turns out, however, the introduction of an auxiliary vector $`V^\mu `$, which to lowest order in $`g`$ equals $`H^\mu `$ on-shell, greatly simplifies the construction: it essentially suffices to determine the deformation up to second order, all higher-order contributions are due to prefactors that are functions of the scalar field $`\varphi `$ only, which are easily completed to all orders.
The full Lagrangian constructed in this way reads
$``$ $`={\displaystyle \frac{1}{2}}\mathrm{cos}(2g\varphi )V^\mu V_\mu V_\mu H^\mu {\displaystyle \frac{1}{4}}\mathrm{cos}(2g\varphi )^{\mu \nu }_{\mu \nu }{\displaystyle \frac{1}{4}}\mathrm{sin}(2g\varphi )^{\mu \nu }_{\mu \nu }`$
$`+{\displaystyle \frac{1}{2}}\mathrm{cos}(2g\varphi )^\mu \varphi _\mu \varphi \mathrm{i}\mathrm{cos}(2g\varphi )\left(\psi ^i\sigma ^\mu \stackrel{}{_\mu }\overline{\psi }_i\right)+2g\mathrm{sin}(2g\varphi )V_\mu \psi ^i\sigma ^\mu \overline{\psi }_i+{\displaystyle \frac{1}{2}}U^2`$
$`\mathrm{i}g_{\mu \nu }\left(\mathrm{e}^{\mathrm{i}g\varphi }\psi ^i\sigma ^{\mu \nu }\psi _i\mathrm{e}^{\mathrm{i}g\varphi }\overline{\psi }^i\overline{\sigma }^{\mu \nu }\overline{\psi }_i\right){\displaystyle \frac{g^2}{\mathrm{cos}(2g\varphi )}}\left(\mathrm{e}^{\mathrm{i}g\varphi }\psi ^i\psi ^j+\mathrm{e}^{\mathrm{i}g\varphi }\overline{\psi }^i\overline{\psi }^j\right)^2.`$ (4)
Here $`_{\mu \nu }`$ is an extended field strength,
$$_{\mu \nu }=(_\mu +gV_\mu )A_\nu (_\nu +gV_\nu )A_\mu .$$
(5)
The Lagrangian is invariant, up to a total derivative, under gauge transformations
$`\delta A_\mu `$ $`=(_\mu +gV_\mu )C`$
$`\delta B_{\mu \nu }`$ $`=gC\left[\mathrm{cos}(2g\varphi )_{\mu \nu }\mathrm{sin}(2g\varphi )_{\mu \nu }2g\left(\mathrm{e}^{\mathrm{i}g\varphi }\psi ^i\sigma _{\mu \nu }\psi _i+\mathrm{e}^{\mathrm{i}g\varphi }\overline{\psi }^i\overline{\sigma }_{\mu \nu }\overline{\psi }_i\right)\right]`$
$`+_\mu C_\nu _\nu C_\mu `$
$`\delta \varphi `$ $`=\delta \psi _\alpha ^i=\delta V_\mu =\delta U=0,`$ (6)
and rigid supersymmetry transformations generated by
$`𝒟_\alpha ^i\varphi `$ $`=\psi _\alpha ^i`$
$`𝒟_\alpha ^iV_\mu `$ $`=\mathrm{i}_\mu \psi _\alpha ^i`$
$`𝒟_\alpha ^iA_\mu `$ $`=\mathrm{i}g\psi _\alpha ^iA_\mu +\mathrm{i}\mathrm{e}^{\mathrm{i}g\varphi }(\sigma _\mu \overline{\psi }^i)_\alpha `$
$`𝒟_\alpha ^iB_{\mu \nu }`$ $`=2\mathrm{cos}(2g\varphi )(\sigma _{\mu \nu }\psi ^i)_\alpha 2g\mathrm{e}^{\mathrm{i}g\varphi }A_{[\mu }(\sigma _{\nu ]}\overline{\psi }^i)_\alpha `$
$`𝒟_\alpha ^i\overline{\psi }_{\dot{\alpha }}^j`$ $`=\frac{1}{2}\epsilon ^{ij}\sigma _{\alpha \dot{\alpha }}^\mu (V_\mu +\mathrm{i}_\mu \varphi )`$
$`𝒟_\alpha ^i\psi ^{\beta j}`$ $`=\frac{1}{2}\epsilon ^{ij}\mathrm{e}^{\mathrm{i}g\varphi }_{\mu \nu }\sigma ^{\mu \nu }{}_{\alpha }{}^{}{}_{}{}^{\beta }+\frac{\mathrm{i}}{2}\epsilon ^{ij}\delta _\alpha ^\beta U\mathrm{i}g\epsilon ^{ij}\psi _\alpha ^k\psi _k^\beta g\mathrm{tan}(2g\varphi )\delta _\alpha ^\beta \psi ^i\psi ^j`$
$`{\displaystyle \frac{\mathrm{i}g}{\mathrm{cos}(2g\varphi )}}\delta _\alpha ^\beta \overline{\psi }^i\overline{\psi }^j`$
$`𝒟_\alpha ^iU`$ $`=\mathrm{cos}(2g\varphi )(\sigma ^\mu _\mu \overline{\psi }^i)_\alpha +g\mathrm{sin}(2g\varphi )(_\mu \varphi \mathrm{i}V_\mu )(\sigma ^\mu \overline{\psi }^i)_\alpha `$
$`g\mathrm{e}^{\mathrm{i}g\varphi }_{\mu \nu }(\sigma ^{\mu \nu }\psi ^i)_\alpha {\displaystyle \frac{2\mathrm{i}g^2}{\mathrm{cos}(2g\varphi )}}\left(\mathrm{e}^{2\mathrm{i}g\varphi }\psi ^i\psi ^j+\overline{\psi }^i\overline{\psi }^j\right)\psi _{\alpha j}.`$ (7)
Using the transformation properties of $`_{\mu \nu }`$,
$`\delta _{\mu \nu }`$ $`=gC(_\mu V_\nu _\nu V_\mu )`$
$`𝒟_\alpha ^i_{\mu \nu }`$ $`=\mathrm{i}g\psi _\alpha ^i_{\mu \nu }+2\mathrm{i}(_{[\mu }+gV_{[\mu })\left(\mathrm{e}^{\mathrm{i}g\varphi }\sigma _{\nu ]}\overline{\psi }^i\right)_\alpha ,`$ (8)
gauge invariance is easily verified, while supersymmetry requires some effort.
Note that although the combination $`_\mu +gV_\mu `$ resembles a covariant derivative, there is no gauge transformation associated with $`V_\mu `$, it is merely an auxiliary field. In order to eliminate it, we collect all terms involving $`V_\mu `$. They can be written as
$$=\frac{1}{2}V_\mu K^{\mu \nu }V_\nu V_\mu ^\mu +\mathrm{},$$
(9)
where we have employed the abbreviations
$`K^{\mu \nu }`$ $`=\mathrm{cos}(2g\varphi )\left[\eta ^{\mu \nu }(1g^2AA)+g^2A^\mu A^\nu \right]`$
$`^\mu `$ $`=\frac{1}{2}\epsilon ^{\mu \nu \rho \sigma }_\nu B_{\rho \sigma }+g\mathrm{cos}(2g\varphi )F^{\mu \nu }A_\nu +g\mathrm{sin}(2g\varphi )F^{\mu \nu }A_\nu `$
$`2g\mathrm{sin}(2g\varphi )\psi ^i\sigma ^\mu \overline{\psi }_i+2\mathrm{i}g^2\left(\mathrm{e}^{\mathrm{i}g\varphi }\psi ^i\sigma ^{\mu \nu }\psi _i\mathrm{e}^{\mathrm{i}g\varphi }\overline{\psi }^i\overline{\sigma }^{\mu \nu }\overline{\psi }_i\right)A_\nu .`$ (10)
The equation of motion of $`V_\mu `$ then yields ($``$ denotes on-shell equality)
$$V_\mu (K^1)_{\mu \nu }^\nu ,(K^1)_{\mu \nu }=\frac{\eta _{\mu \nu }g^2A_\mu A_\nu }{\mathrm{cos}(2g\varphi )(1g^2AA)},$$
(11)
which upon insertion into the Lagrangian and the transformations results in expressions non-polynomial in both $`gA_\mu `$ and $`g\varphi `$. If one expands the Lagrangian in powers of $`g`$, the coupling $`gA_\mu J^\mu `$ is recovered to first order, which was the goal of the construction.
The supersymmetry algebra closes on-shell; if the equations of motion hold, the commutator of two supersymmetry transformations $`\mathrm{\Delta }_\xi =\xi _i^\alpha 𝒟_\alpha ^i+\overline{\xi }_{\dot{\alpha }}^i\overline{𝒟}_i^{\dot{\alpha }}`$ yields a translation and a gauge transformation,
$$[\mathrm{\Delta }_\xi ,\mathrm{\Delta }_\zeta ]a^\mu _\mu \delta ,$$
(12)
with parameters
$`a^\mu `$ $`=\mathrm{i}\left(\zeta _i\sigma ^\mu \overline{\xi }^i\xi _i\sigma ^\mu \overline{\zeta }^i\right)`$
$`C`$ $`=a^\mu A_\mu +{\displaystyle \frac{\mathrm{i}}{g}}\left(\mathrm{e}^{\mathrm{i}g\varphi }\xi _i\zeta ^i+\mathrm{e}^{\mathrm{i}g\varphi }\overline{\xi }^i\overline{\zeta }_i\right)`$
$`C_\mu `$ $`=a^\nu B_{\nu \mu }+{\displaystyle \frac{1}{2g}}\mathrm{sin}(2g\varphi )a_\mu +\left(\mathrm{e}^{\mathrm{i}g\varphi }\overline{\xi }^i\overline{\zeta }_i\mathrm{e}^{\mathrm{i}g\varphi }\xi _i\zeta ^i\right)A_\mu .`$ (13)
Note that while the parameter $`C`$ is singular for $`g=0`$, the transformation $`\delta `$ is not, since $`C`$ occurs either under a derivative or with a factor $`g`$.
It is quite remarkable that the gauged central charge transformations do not commute with supersymmetry transformations, not even on-shell. Rather, one has
$$[\mathrm{\Delta }_\xi ,\delta ]\delta ^{},$$
(14)
where now
$`C^{}`$ $`=\mathrm{i}gC\left(\xi _i\psi ^i\overline{\xi }^i\overline{\psi }_i\right)`$
$`C_\mu ^{}`$ $`=gC\left(\mathrm{e}^{\mathrm{i}g\varphi }\xi ^i\sigma _\mu \overline{\psi }_i\mathrm{e}^{\mathrm{i}g\varphi }\psi ^i\sigma _\mu \overline{\xi }_i\right).`$ (15)
A free vector-tensor multiplet is dual to an abelian vector multiplet. Since in our interacting model the 2-form $`B_{\mu \nu }`$ occurs only via its field strength, it too can be dualized into a pseudo-scalar. In fact, the use of an auxiliary vector makes the dualization particularly simple. Considering $`H^\mu `$ as an independent field and implementing the Bianchi identity $`_\mu H^\mu =0`$ by means of a Lagrange multiplier field $`\phi `$, we find that according to the equation of motion of $`H^\mu `$, we can replace $`V_\mu `$ with $`_\mu \phi `$,
$$V_\mu _\mu \phi ,𝒟_\alpha ^i\phi =\mathrm{i}\psi _\alpha ^i.$$
(16)
$`\varphi `$ and $`\phi `$ can be combined into a complex scalar $`X`$ that is chiral, $`\overline{𝒟}_{\dot{\alpha }i}X=0`$. With suitable field redefintions, we obtain the (on-shell) field content of an $`N=2`$ vector multiplet with the standard supersymmetry transformations<sup>1</sup><sup>1</sup>1Every function of $`(\phi \mathrm{i}\varphi )`$ is chiral, but only $`X`$ as in (17) satisfies the additional constraint $`𝒟^i𝒟^jX+\overline{𝒟}^i\overline{𝒟}^j\overline{X}=0`$ that gives rise to a vector multiplet.,
$$X=\frac{1}{2g}\mathrm{e}^{g(\phi \mathrm{i}\varphi )},\widehat{A}_\mu =\mathrm{e}^{g\phi }A_\mu ,\lambda _\alpha ^i=\mathrm{e}^{g(\phi \mathrm{i}\varphi )}\psi _\alpha ^i.$$
(17)
Up to a factor, $`_{\mu \nu }`$ turns out to be nothing but the field strength of the new vector $`\widehat{A}_\mu `$,
$$_{\mu \nu }=\mathrm{e}^{g\phi }\widehat{F}_{\mu \nu }.$$
(18)
The holomorphic prepotential $`f(X)`$ that determines the action of the vector multiplet can be obtained from the field strength terms in the Lagrangian (4),
$$\frac{1}{4}\mathrm{cos}(2g\varphi )^{\mu \nu }_{\mu \nu }\frac{1}{4}\mathrm{sin}(2g\varphi )^{\mu \nu }_{\mu \nu }=\frac{1}{4}\mathrm{Im}\left[\frac{\mathrm{i}}{(2gX)^2}\widehat{F}_{\mu \nu }(\widehat{F}^{\mu \nu }\mathrm{i}\widehat{F}^{\mu \nu })\right].$$
Comparing with the general expression common to every vector multiplet,
$$\frac{1}{4}\mathrm{Im}\left[f^{\prime \prime }(X)\widehat{F}_{\mu \nu }(\widehat{F}^{\mu \nu }\mathrm{i}\widehat{F}^{\mu \nu })\right],$$
we read off the second derivative of the prepotential, and conclude that
$$f(X)=\frac{\mathrm{i}}{4g^2}\mathrm{ln}(gX).$$
(19)
Finally, we briefly demonstrate how the present model can be derived from an $`N=1`$ superspace integral, drawing on results obtained in :
Let us embed $`B_{\mu \nu }`$, $`\varphi `$ and $`\psi _\alpha ^1`$ in a chiral spinor superfield $`\mathrm{\Psi }_\alpha (x,\theta ,\overline{\theta })`$, $`A_\mu `$ and $`\psi _\alpha ^2`$ in a real superfield $`A(x,\theta ,\overline{\theta })`$, and $`V_\mu `$ in a real superfield $`V(x,\theta ,\overline{\theta })`$ with a mass dimension shifted by $`+1`$. From these, we construct field strength superfields $`Y_\alpha (x,\theta ,\overline{\theta })`$ and $`W_\alpha (x,\theta ,\overline{\theta })`$,
$$Y_\alpha =\frac{\mathrm{i}}{4}\overline{D}^2\left[\mathrm{e}^{2\mathrm{i}gV}D_\alpha (\mathrm{e}^{\mathrm{i}gV}A)\right],W_\alpha =\frac{1}{2}\overline{D}^2D_\alpha V,$$
(20)
where $`D_\alpha `$ and $`\overline{D}_{\dot{\alpha }}`$ are the usual $`N=1`$ supercovariant derivatives. $`\mathrm{\Psi }_\alpha `$, $`Y_\alpha `$ and $`W_\alpha `$ are chiral,
$$\overline{D}_{\dot{\alpha }}\mathrm{\Psi }_\alpha =\overline{D}_{\dot{\alpha }}Y_\alpha =\overline{D}_{\dot{\alpha }}W_\alpha =0,$$
(21)
so the action
$$S=d^4xd^2\theta \left(W^\alpha \mathrm{\Psi }_\alpha +Y^\alpha Y_\alpha \frac{1}{4g^2}d^2\overline{\theta }\mathrm{e}^{2\mathrm{i}gV}\right)+\text{c.c.}$$
(22)
is manifestly $`N=1`$ supersymmetric. The component gauge transformations (6) follow from the superfield transformations
$$\delta \mathrm{\Psi }_\alpha =\mathrm{i}\overline{D}^2D_\alpha R2\mathrm{i}g\mathrm{\Lambda }Y_\alpha ,\delta A=\mathrm{i}\left(\mathrm{e}^{\mathrm{i}gV}\mathrm{\Lambda }\mathrm{e}^{\mathrm{i}gV}\overline{\mathrm{\Lambda }}\right),\delta V=0,$$
(23)
where $`R(x,\theta ,\overline{\theta })`$ is a real and $`\mathrm{\Lambda }(x,\theta ,\overline{\theta })`$ a chiral superfield. Gauge invariance of the action is then due to the relations
$$\delta Y_\alpha =\mathrm{i}g\mathrm{\Lambda }W_\alpha ,D^\alpha W_\alpha \overline{D}_{\dot{\alpha }}\overline{W}^{\dot{\alpha }}=0,$$
(24)
the former being the superfield analog of the gauge transformation (8) of $`_{\mu \nu }`$. Passing to Wess-Zumino gauge for $`A`$ and $`\mathrm{\Psi }_\alpha `$ and eliminating the auxiliary fields (except $`V_\mu `$) results in the $`N=2`$ supersymmetric Lagrangian (4).
Evidently, there is an infinite number of $`N=1`$ supersymmetric Henneaux-Knaepen models that describe the field content of a vector-tensor multiplet, for one may specify different functions of $`V`$ (the $`d^2\theta d^2\overline{\theta }`$-part in the action) that allow to eliminate the auxiliary fields (another example was given in ), but in general these will not possess a second supersymmetry.
The previously known interactions of the vector-tensor multiplet can all be derived by dimensional reduction from interactions of the $`N=(1,0)`$ tensor multiplet in six dimensions . In these cases the gauged central charge is a remnant of translations in the additional spacelike directions. It would be interesting to know whether this applies to the present model as well, or whether the fact that the gauge field for the local central charge transformations resides in the vector-tensor multiplet itself requires a different mechanism in six dimensions. Models that might be relevant in this context can be found in .
The new model introduced here indicates that our current knowledge about the possible interactions of the vector-tensor multiplet is far from being exhaustive. In particular, the problem of finding interactions between several such multiplets has not been attacked successfully as yet. A possible starting point would be to consider the general $`N=1`$ supersymmetric Henneaux-Knaepen models of and to single out those that exhibit an SU(2) $`R`$-symmetry. Work in this direction is under way.
*Remark:* In the course of this work, a further $`N=2`$ supersymmetric Henneaux-Knaepen model was found, involving two vector multiplets and a double-tensor multiplet .
Acknowledgements
I would like to thank Friedemann Brandt and Sergei Kuzenko for helpful discussions and suggestions. Thanks also to Augusto Sagnotti for bringing the work on six-dimensional tensor multiplet couplings to my attention.
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# Light Element Abundance Patterns in the Orion Association: I) HST Observations of Boron in G-dwarfs
## 1 Introduction
The first results for boron abundances in stars of the Orion association were provided by a study of the B ii 1362 Å line in four B-type stars (Cunha et al. 1997). The non-LTE abundances derived for these Orion stars were between log $`ϵ`$(B) = 2.5 - 2.9, i.e. close to the boron abundance measured in solar system meteorites (log $`ϵ`$(B) = 2.78) or the solar photosphere (log $`ϵ`$(B) = 2.7, where log $`ϵ`$(X)=log N(X)/N(H) + 12). A second study of boron in the Orion association using the G-dwarf BD-5 1317 found a boron abundance lower by about 0.7 dex than the solar/meteoritic value (Cunha, Smith & Lambert 1999). This particular Orion member was picked because i) it had an undepleted lithium abundance and ii) it exhibited one of the largest oxygen abundances measured in the association (log $`ϵ`$(O)=9.11; Cunha, Smith & Lambert 1998). The absence of Li depletion guarantees that the boron was not destroyed in the stellar interior. As a consequence, the observed boron abundance is expected to reflect the original abundance at the time of star formation. As for the high oxygen abundance, it is attributed to the self-enrichment of the Orion cloud (Cunha & Lambert 1994), caused by the explosion of massive stars as type II supernovae (SN II). According to this model, the O-rich SN II ejecta ‘pollute’ certain pockets of gas around the site of the explosion, from which new stars form with enhanced oxygen abundances. This idea is corroborated by the fact that the most O-rich stars in the Orion association are also the youngest members and generally co-located (Cunha et al. 1998).
In the present study, two additional stars are examined. They consist of G-dwarfs whose radial velocities and proper motions indicate that they are definite members of the Orion association (Cunha, Smith & Lambert 1995). These two stars also show different degrees of oxygen enrichment and fulfill the requirement that they do not show significant Li depletion. They can thus be used to study the self-enrichment process within the Orion clouds. This is an important goal from the theoretical point of view, since SN II’s are known to be the source not only of O, but also indirectly of LiBeB. According to current knowledge (e.g. Vangioni-Flam et al. 2000), beryllium and boron can be synthesized only by spallative processes in which heavier nuclei (mainly C, N and O) are spallated in reactions induced by energetic particles (EPs). Two distinct mechanisms have been invoked, depending on whether the EPs are nuclei or neutrinos. In the first case, B (and Li and Be) is produced by the spallation of C, N, and O by protons and $`\alpha `$’s. The other mechanism, known as the $`\nu `$-process (or neutrino-spallation), consists in breaking nuclei with energetic neutrinos released in enormous numbers during the explosion of Type II Supernovae. The carbon nuclei spalled in this process to make <sup>11</sup>B are those synthesized by the massive SN II progenitor. As a consequence, the $`\nu `$-process boron is ‘ready-mixed’ with the CNO in the SN II ejecta and all these nuclei are released together into the ambient medium. The $`\nu `$-process has been invoked to cope with a classical problem resulting from the confrontation of the observed (solar system) values of the B/Be and <sup>11</sup>B/<sup>10</sup>B ratios with those predicted by the nucleo-spallation process alone that produces too little <sup>7</sup>Li and <sup>11</sup>B. The $`\nu `$-process is predicted to produce <sup>7</sup>Li and <sup>11</sup>B but virtually no <sup>6</sup>Li, no <sup>9</sup>Be and no <sup>10</sup>B. Therefore, it offers a way to resolve the classical problem.
From the detailed analyses of Be and B abundances in metal-poor halo stars, it has been recently proposed that the EPs responsible for most of the light element production are probably accelerated inside superbubbles created by repeated SN II’s in OB associations (Parizot & Drury, 1999, 2000; Parizot, 2000; Ramaty et al. 2000). Since it is known that the Orion association blew just such a superbubble (SB) from the winds and explosions of massive stars in subgroup Ia (e.g. Brown et al. 1995), it may be expected that significant Li, Be and B production occurred recently in or close to the Orion molecular cloud. The self-enrichment in O of the Orion stellar association should therefore be accompanied by a significant enrichment in light elements as well, and the relation between the two processes is worth studying.
In the last decade, B abundances have been measured in disk and halo stars (Duncan, Lambert & Lemke 1992; Duncan et al. 1997; García-López et al. 1998; Primas et al. 1999 and Cunha et al. 2000) showing a global positive correlation of boron with oxygen which is understood as the progressive enrichment of the interstellar medium (ISM) in both metals and light elements. However, the situation in local environments has never been studied, and one may ask whether the same kind of correlation also holds on smaller scales, within individual stellar associations. Our study intends to answer the question in the case of Orion.
## 2 STIS Observations and Spectra
The two Orion targets, BD -6 1250 and HD 294297, were observed with the STIS spectrograph on the Hubble Space Telescope (HST), over 10 and 11 orbits each, which were needed in order to obtain combined spectra with a S/N $``$50 for the analysis of B I at 2496.771 Å. (The measured S/N across the spectral region synthesized ranged from 40-60 in both stars). The observations were obtained with the MAMA detector in the ACCUM mode with the first order grating G230M, plus 52X0.1 slit to record the spectral region between 2454 and 2545 Å with a resolution R$``$14,000.
The observed STIS spectra were processed by the pipeline calibration from STScI with the IRAF package ‘calstis’ which was used in order to subtract the bias, divide by the flatfield and wavelength calibrate the data. These pipeline calibrated spectra were then combined with the IRAF task ’mscombine’ and extracted one-dimensional spectra were obtained with the task ’apsum’. Unlike the spectra obtained with the echelle grating, the levels of scattered light in the first-order grating spectra are insignificant and the step of subtraction of scattered light was not needed. The final spectra of the two Orion targets are shown in the two panels of Figure 1.
## 3 Analysis
The stellar parameters adopted in the calculation of the model atmospheres for the two studied stars (plus a sample of Orion members) are presented in Table 1. The effective temperatures from Cunha et al. (1995) were obtained from a calibration of the Strömgren $`\beta `$ indices with spectroscopic effective temperatures defined by the Fe I lines present in the optical spectra of slowly rotating stars in the direction of Orion. The surface gravities for the rapid rotators were assumed in Cunha et al. (1995) to be 4.0. Here we have adjusted the log g for HD 294297 to a higher value (log g=4.4) because this produced an overall better fit in the spectral region between 2494.5 and 2499.5 Å. The sensitivity of the derived boron abundance to changes in log g will be discussed in Section 4.2.
Boron abundances were derived from spectrum synthesis of the region around the B i resonance line at 2496.771 Å, which is the less blended boron resonance line available. In an effort to place all derived boron abundances for near-solar metallicity stars on the same absolute scale, the present analysis of the two Orion targets is entirely consistent with the previous analysis of a sample of near-solar metallicity solar type dwarfs (with \[Fe/H\] ranging from -0.75 to +0.15) from HST archival data (Cunha et al. 2000). The analysis of the B I spectral region in near solar-temperature and solar-metallicity stars is complicated by the definition of the continuum level; the spectra are crowded with strong absorption lines and no regions are free from absorption. As discussed in previous studies of B I (Cunha & Smith 1999; Cunha et al. 1999; 2000), the continuum level in this region was set for the Sun, where observed specific intensities are available, using the combination of synthesis code, line list, and solar model atmosphere. The relative position of the continuum to spectral line depths does not vary much over the limited range of T<sub>eff</sub> spanned by the Orion G stars. The continuum level is allowed to vary slightly in order to provide the best fit to the absorption lines.
This analysis technique was applied to 14 field F and G dwarfs (Cunha et al. 2000), for which reasonably accurate distances were available, and a consistency check was carried out between predicted and observed continuum fluxes. A model surface flux at 2500 Å, F<sub>2500</sub>, comes from the model atmosphere, and the predicted continuum flux observed at the earth, f<sub>2500</sub>, is f<sub>2500</sub>=F<sub>2500</sub>(R/D)<sup>2</sup>, where $`R`$ is the stellar radius and $`D`$ is the distance; to compute the stellar radius requires distance, apparent magnitude, reddening, a bolometric correction, and T<sub>eff</sub>. We call this the calculated continuum flux, f(calc). This flux can be compared to the empirical continuum flux at 2500 Å set by the spectrum synthesis, which we call the observed continuum flux, f(obs).
The comparison of f(calc)/f(obs) for the field F and G dwarfs from Cunha et al. (2000) was quite good (their Figure 3): with an average f(calc)/f(obs) being 0.91$`\pm `$0.19 and having no trend in this ratio with T<sub>eff</sub> over the range 5600-6700 K. This same technique can be applied to the Orion stars as the distance to the Orion association is reasonably well-known, e.g. Warren & Hesser (1978); de Zeeuw et al. (1999). The Orion association subgroups Ib and Ic extend over a distance of $``$40 pc (de Zeuw et al. 1999) and we adopt a single distance of 490 pc for the 3 G-dwarfs observed to date by HST. A reddening of E(B-V)=0.05 is assumed for all Orion stars and the (small) bolometric corrections are from Böhm-Vitense (1989). Results for f(cal)/f(obs) in the Orion members are shown versus T<sub>eff</sub> in the top panel of Figure 2, along with the previously derived results for the field F and G dwarfs from Cunha et al. (2000). The Orion members agree well with the field stars and indicate that the analysis technique used to define the continuum is consistent. The mean value of f(calc)/f(obs) for the Orion members is 0.89$`\pm `$0.18, very similar to the field stars. The comparison between the field-star sample and the Orion stars indicates that there are not significant model-induced systematic differences between the two sets of analyses.
A quantitative estimate of the uncertainties associated with empirically fitting the continuum levels can be made by investigating line depths as functions of T<sub>eff</sub> and \[Fe/H\] for the Orion targets, as well as the larger sample of field dwarfs studied in Cunha et al. (2000). Line depths measured relative to a continuum level should depend primarily on the effective temperature and metallicity. The flux at a local high point in the spectra of these stars, near 2500 Å (whose level is controlled by line blanketing), was then compared to the continuum flux derived empirically from the spectrum synthesis (f(obs)) at this wavelength. It is found that this ratio of the observed line depth to the empirically derived continuum is a well-defined, smooth curve as a function of T<sub>eff</sub>, for near solar-metallicity stars (with \[Fe/H\]$``$-0.2). This is shown in the bottom panel of Figure 2 where we plot this ratio versus the effective temperature for the field and Orion dwarfs. Inspection of the figure shows that this ratio varies from $``$0.25 at T<sub>eff</sub>=5600K, to $``$0.65 at T$`{}_{eff}{}^{}=`$6700K. Note that the line-depth to continuum values increase for the metal-poor stars, as expected due to reduced metal line blanketing. The observed scatter of the ratios about the mean curve defined for the near-solar metallicity stars is $`\pm `$0.04. The flux ratio of the local high point (at 2500 Å) to the continuum varies from 0.3 to 0.4 for the three studied Orion stars and these three stars follow the relation of flux ratio versus T<sub>eff</sub> defined by the near-solar metallicity field stars. If the continuum flux level for the Orion G-dwarfs is thus allowed to vary by $`\pm `$0.04, the resulting derived B abundance varies by $`\pm `$0.11 dex: this is a fair measurement of the abundance uncertainty introduced from the uncertainties in the continuum level.
Our synthetic spectra were computed using the program LINFOR (originally developed at Kiel University by H. Holweger, M. Steffen & W. Steenbock) and model atmospheres generated with the ATLAS9 code (R. L. Kurucz, 1993 - private communication) for the stellar parameters in Table 1. The calculations of the synthetic spectra included a more modern value for the bound-free cross-section of Mg i 3p$`{}_{}{}^{3}P_{}^{o}`$, which is the dominant source of continuous opacity in the 2500 Å region. The importance of this choice is pointed out in Cunha & Smith (1999). In the calculation of the synthetic spectra the adopted Mg abundance was scaled with the oxygen abundance for each star. The adopted line list was compiled first in order to fit the disk-center spectrum of the Sun, later this line list was fine-tuned in order to fit the spectra of the sample of near-solar metallicity dwarfs (mentioned above) with effective temperatures in the range between 5650K and 6700K. This line list can be found in Cunha et al. (2000). In Figure 3 we illustrate synthetic spectra for the studied stars. Using a simple $`\chi ^2`$ minimization we derive an LTE boron abundance of log $`ϵ`$(B)=2.60$`\pm `$0.20 for HD 294297 and log $`ϵ`$(B)=2.30$`\pm `$0.20 for BD -6 1250 with the uncertainties set by the sharpness of the $`\chi ^2`$ minima. Corrections for non-LTE effects in stars of solar temperatures and metallicities are small: Kiselman & Carlsson’s (1996) non-LTE calculations for the B i lines indicate revised non-LTE abundances of log $`ϵ`$(B)=2.34 and log $`ϵ`$(B)=2.65, respectively, for HD 294297 and BD -6 1250. The derived boron abundances for the target stars plus their derived Li, Fe and O abundances are assembled in Table 1. We also added to this table other Orion members that have been studied in previous papers (Cunha et al. 1995, 1998) and that will be brought into the discussion of the abundance results and uncertainties that follow. We note the addition of one star (P1374) that has not been published previously.
## 4 Results and Discussion
### 4.1 The Abundances
A summary of boron and oxygen abundances in Orion is presented in Figure 4. Combined with the non-LTE B-star boron abundances<sup>1</sup><sup>1</sup>1A large correction for non-LTE effects on the B ii 1362 Å line was included (Cunha et al. 1997). The B iii 2066 Å line was predicted to be minimally affected by departures from LTE. Proffitt et al. (1999) observed and analyzed the 2066 Å line in one of the four stars. Our reanalysis (Lambert et al. 2000) of their spectrum with our model atmosphere gives a non-LTE abundance in fair agreement with the non-LTE value from the B ii 1362 Å line, bolstering our confidence in the B ii non-LTE corrections., the three points defined by the Orion solar-type stars indicate clearly that there is no positive trend of boron with oxygen in Orion. In fact, if there is any discernable trend of B and O, it is an anticorrelation. A possible anticorrelation is made all the more convincing by the combination of two sets of results from very different types of stars: B stars with T<sub>eff</sub> $``$ 18000-22000K (analyzing B ii) and G stars with T<sub>eff</sub> $``$ 5850-6150K (analyzing B i). Taken together, the B and G stars seem to define a single relation of decreasing boron with increasing oxygen among the stellar members of Orion. The most oxygen-rich and boron-poor Orion star in Figure 4 is the G-star BD -5 1317, whose B I spectrum was analyzed in Cunha et al. (1999). Much of the weight of a significant anticorrelation between B and O falls on this star. Cunha et al. discussed whether the boron abundance derived from B I could be influenced by some unkonwn effect, most probably a chromosphere, but were unable to find an obvious explanation for the weakened B I line in BD -5 1317 (including the addition of a chromosphere to the model atmosphere), other than a low B abundance. This star was studied for Li, Fe, and O in Cunha et al. (1995, 1998) and no obvious spectral peculiarities were noted; however, future spectroscopic studies of this star are encouraged.
If a linear least-squares fit is performed on the log $`ϵ`$(B) versus log $`ϵ`$(O) data from Orion, a slope of -1.1$`\pm `$0.2 is found (with a correlation coeficient r= -0.94), indicating a significant decrease of B with O.
To probe the robustness of this apparent decrease of B as O increases, tests on the boron and oxygen abundance dataset were conducted. As a first step, the data point for BD -5 1317 was excluded from another linear least-squares fit on the remaining 6 points; in this case a significant anticorrelation (r=-0.90) is still found between B and O with a slope of -1.3$`\pm `$0.3 (within the errors, the same slope found with the inclusion of BD -5 1317). Reduced values of $`\chi ^2`$ can also be computed from the expression
$$\chi _\mathrm{r}^2=\frac{1}{\nu 1}\mathrm{\Sigma }\frac{(observed_ifitted_i)^2}{\sigma _i^2},$$
(1)
where ‘observed’ and ‘fitted’ refer to the observed values and linear least-squares computed values, respectively, $`\sigma `$ is the associated error, and $`\nu `$ is the number of degrees of freedom. Using $`\sigma `$=0.20 dex, the values of $`\chi _\mathrm{r}^2`$ are 1.65 (with the inclusion of BD -5 1317) and 1.53 (with the exclusion of BD -5 1317): both values of $`\chi _\mathrm{r}^2`$ indicate a good fit with a linear trend of log $`ϵ`$(B) versus log $`ϵ`$(O). If a zero slope between B and O is assumed (i.e. no trend), the associated values of $`\chi _\mathrm{r}^2`$ are 15.40 (with the inclusion of BD -5 1317) and 7.82 (with the exclusion of BD -5 1317): both of these values indicate extremely poor fits to the data (well past the 3$`\sigma `$ level of confidence). Because the observed trend is significant and opposite to the global correlation between B and O in the Galaxy, over three orders of magnitude in metallicity, an investigation into whether errors can produce such an anticorrelation is carried out.
### 4.2 Elemental Correlations and Stellar Parameter Errors
As this paper deals with an analysis of G-type stars, the discussion of possible errors in boron and oxygen abundances is confined to stars of this spectral type for which the B i 2497 Å line and the O i 7770 Å triplet provide the B and O abundances, respectively. The ground-based and HST spectra analyzed for both oxygen and boron are of relatively high S/N, and errors in the derived abundances are due primarily to uncertainties in the crucial stellar parameters: effective temperature (T<sub>eff</sub>), surface gravity (log g), and microturbulent velocity ($`\xi `$). Typical O i equivalent widths for the three G stars were used as input in order to compute the effects of the parameter changes on derived O abundances. For boron uncertainties, the discussions from Primas et al. (1999) and Boesgaard et al. (1998) were used. The results are that the oxygen abundances change by -0.08 dex for $`\mathrm{\Delta }`$T = +100K, +0.06 dex for $`\mathrm{\Delta }`$log g = +0.3 dex, and -0.03 dex for $`\mathrm{\Delta }`$$`\xi `$ = +0.2 km s<sup>-1</sup>: the corresponding numbers for boron are +0.10 (for $`\mathrm{\Delta }`$T), -0.02 (for $`\mathrm{\Delta }`$log g), and -0.06 (for $`\mathrm{\Delta }`$$`\xi `$) dex, respectively. Of special note are the anticorrelated O and B abundance errors for changes in temperature and surface gravity; for example, if the effective temperature of a star were overestimated by 100K, the derived O abundance would be too low by 0.08 dex, while the derived B abundance would be too high by 0.10 dex. This effect results in spurious anticorrelations between B and O for random temperature and gravity errors; however, we note that the observed scale of the anticorrelation over the range of derived B and O abundances would require errors much larger than estimated, for example, $`\mathrm{\Delta }`$T $``$ 600K.
Further constraints on the reality of a B-O anticorrelation are provided by the Fe abundances from Cunha et al. (1995; 1998) for 9 Orion members. These stars span a T<sub>eff</sub> range of 5600-6150K and a tight distribution of Fe abundances is obtained with a spread of $`\pm `$0.13 dex: this scatter is completely accounted for by the observational errors, as argued by Cunha et al. Therefore, the Orion F-G stars show a single value of the Fe abundance, but not a single O abundance. For the Fe I lines used by Cunha et al., it is found that Fe abundances change by +0.08 dex for $`\mathrm{\Delta }`$T= +100K, -0.02 dex for $`\mathrm{\Delta }`$log g=+0.3 dex, and -0.03 dex for $`\mathrm{\Delta }`$$`\xi `$= +0.2 km s<sup>-1</sup>. Note that B I and Fe I have a similar behavior. A 600K error in T<sub>eff</sub>’s would clearly lead to a noticeable spread in the Fe abundances but that is not seen.
In order to test the magnitudes of spurious anticorrelations produced by random T<sub>eff</sub>, log g, and $`\xi `$ errors, a program was used to generate gaussian distributed noise of specified means and standard deviations in these parameters, which were then used to compute errors in boron, oxygen, and iron abundances. Different starting model abundances were used in order to understand under what conditions a spurious anticorrelation, of the magnitude observed, between B and O could be generated. The fact that two very different types of stars (spectral types B and G) seem to fall along a single relationship is ignored here, and only uncertainties relevant to the G stars are considered. The internal abundance uncertainties from Cunha et al. (1995; 1998) and this study, for the Orion members, are dominated primarily by uncertainties in the stellar parameters T<sub>eff</sub>, log g, and $`\xi `$. Estimates of these uncertainties are $`\pm `$150K in T<sub>eff</sub>, $`\pm `$0.3 dex in log g, and $`\pm `$0.2 km s<sup>-1</sup> in $`\xi `$. The estimated errors in stellar parameters are taken from discussions in Cunha et al. (1995, 1998) for the Orion F/G dwarfs. In these previous studies, based upon a comparison of photometric and spectroscopic effective temperature scales, average differences of $`\pm `$70K were found. We adopt a conservative approach here and double these errors to $`\mathrm{\Delta }`$T=150K. The uncertainties in surface gravity and microturbulence are those from Cunha et al. (1995; 1998) and are not as critical, as O I and B I are most sensitive to effective temperature.
A discussion of the abundance results begins with Fe and O in the Orion G-stars from Cunha et al. (1998) who found that the observed Fe abundances could be modelled as a single value, while the values of the derived O abundances were too scattered to be explained by a single value. This contention is re-examined here using a different analysis. The error simulation technique employed here consists of a beginning set of model data points that represent the observed sample: in the case of Fe and O, this consists of 9 points, while in the case of B and O it is 3 points. An input abundance distribution is assumed; for example, in the initial modelling of Fe and O, a constant abundance value for each element in the Orion members is tried. Random T<sub>eff</sub>, log g, and $`\xi `$ errors are then generated for each input data point, resulting in changes to the input Fe and O abundances ($`\delta `$Fe and $`\delta `$O) which are then added to the input model abundances. The input abundances have now been perturbed by random stellar parameter errors. Because this investigation is to probe possible correlations (or anticorrelations) between pairs of elements, a linear least-squares fit is then performed on the perturbed model points and a slope derived: this slope can be compared to the corresponding slope derived from the observed abundances. The above procedure can be performed an arbitrary number of times (each time is labelled as a “realization”), with a slope computed for each realization. A distribution of slopes can then be constructed and compared to the observationally derived slope. The impact that stellar parameter uncertainties can have on any underlying abundance correlations or anticorrealtions can then be investigated. In addition, standard deviations in the abundances from the perturbed model points can be compared to the observed values. This exercise can thus test whether the estimates of stellar parameter uncertainties are reasonable, as well as the reality of possible correlations (or anticorrelations or no correlations) between pairs of elements.
The top panel of Figure 5 shows the Fe versus O abundances as taken from Cunha et al. (1998). No obvious trend exists and a linear regression finds no correlation, with an insignificant positive slope of 0.1$`\pm `$0.2 derived. Note that in the previous discussion of the various elemental abundance sensitivities to stellar parameters, the Fe I and O I lines have both temperature and gravity sensitivties that are of nearly equal magnitudes but in opposite senses, thus, substantial, random stellar parameter errors would produce anticorrelated Fe and O abundances. This is clearly not observed and suggests already that random errors are unlikely to be responsible for the observed B-O anticorrelation. However, an error simulation of the Fe and O abundances can provide clues as to whether our error estimates are reasonable, and how these errors translate into possible effects on the B versus O trend.
The bottom panel of Figure 5 illustrates simulated data given uncertainties in the stellar parameters: distributions of fitted linear slopes to log $`ϵ`$(Fe) versus log $`ϵ`$(O) abundances are shown (for 1000 realizations) using two different underlying model abundances; a model with a single Fe and O abundance and a model with single Fe and varying O. The slopes were derived from 9 input points (as in the observed sample). For the input model consisting of constant Fe and O abundances we took the average log $`ϵ`$(Fe)= 7.35 and log$`ϵ`$(O)= 8.70, to which random abundance errors were added given errors in T<sub>eff</sub>, log g, and $`\xi `$. Because the Fe I and O I lines are most sensitive to temperature, and in opposite senses, a false anticorrelation is derived for the model points, which is manifested in the bottom panel of Figure 4 as the distribution of negative slopes centered on d(Fe)/d(O)= -0.8, reflecting the dominant sensitivites to temperature of d(Fe)= +0.08 dex/100K and d(O)= -0.08 dex/100K. The average slope of -0.8, instead of -1.0 which would be derived from temperature errors only, results from the lower sensitivity of Fe I to gravity errors when compared to O I.
The observed slope of +0.1 (insignificant from zero slope) is far from the derived average slope in the simulations of -0.80. In addition, using $`\mathrm{\Delta }`$T<sub>eff</sub> =150K, $`\mathrm{\Delta }`$(log g)= 0.3 dex and $`\mathrm{\Delta }`$$`\xi `$= 0.2 km s<sup>-1</sup>, the Fe scatter is found to be $`\pm `$0.12 dex: in excellent agreement with the observed 0.13 dex. Oxygen, on the other hand, is found to scatter in the simulations by 0.13 dex, far less than the observed 0.24 dex. As found by Cunha et al. (1998), an intrinsic abundance scatter within the Orion members of $``$0.5 dex is required. If this oxygen spread is then included in the input model data (while retaining a constant Fe abundance) the second distribution of simulated slopes is shown in the bottom panel of Figure 5 (again for 9 model points and 1000 realizations). Here the derived slopes from the simulated data are close to zero, as observed in the real Orion members. The difference between the slope distributions in the two models is caused by the oxygen abundance spread, as opposed to a single value. With a spread in O, the leverage to create a slope, produced by random errors, is lessened, as well as the scatter in the derived slopes. These simulations verify the conclusions from Cunha et al. (1998) that the Orion F- and G-stars examined had a uniform Fe abundance, but significantly varying O abundances (presumed to be due to selective enrichment of the Orion Association from very massive SN II).
The trend of boron versus oxygen is now examined in light of the results obtained for iron and oxygen. Figure 6 illustrates the simulations conducted on the B and O data, with the top panel again showing the B versus O abundances derived for the Orion G-dwarf members. As was the case with Fe and O, if a single value for the O and B abundances are assumed for all Orion members, random stellar parameter errors will produce an apparent anticorrelation of B with O; however, for realistic values of T<sub>eff</sub>, log g, and microturbulence uncertainties, the model abundance scatter is much smaller than that observed. Inputting the O scatter of $``$0.5 dex, as suggested by the Fe-O analysis, two model simulations are shown in the bottom panel of Figure 6: each distribution of slopes is generated from 3 model points (as in the observed sample), each run through 1000 realizations. One input model had a constant B abundance, and the distribution of slopes is scattered about a mean slope of 0.0, very far from the observed slope of -1. The second input model assumed an anticorrelations of B with O, with log $`ϵ`$(B) proportional to log $`ϵ`$(O)<sup>-1</sup>. In this case, with the given uncertainties in stellar parameters, the slope distribution scatters about -1.0. Note that with only 3 points, the scatter in slopes is larger than for the 9 points used for the Fe and O model. These simulations suggest that the combination of Fe, O, and B abundances in the Orion members indicate a constant Fe abundance, variable O abundance, and a variable B abundance that is inversely proportional to oxygen.
Finally, lithium abundances for Orion members, when compared to oxygen abundances, can further constrain the interpretation of boron in Orion. The top panel of Figure 7 shows the Li and O abundances for 10 Orion members. Interpretation of Li must take into account the susceptibility of this element to destruction by warm protons. Two separate symbols are used in Figure 7: the filled symbols denote stars with lithium near the expected undepleted abundance of log $`ϵ`$(Li)= 3.3, while the open symbols are stars which have almost certainly depleted lithium somewhat. In Cunha et al. (1995), it was found that the lower Li abundances in Orion members were for the slowest rotators (Vsin$`\iota `$ $``$ 10 km s<sup>-1</sup>) and may be stars in which lithium has been destroyed from mixing induced by rotational spindown. If the low Li stars are set aside, the very flat values of undepleted Li versus O in the top panel of Figure 7 suggest lithium is independent of oxygen. Because the Li i line is very temperature sensitive (and in the opposite sense to O i lines), the lack of a large scatter in the Li abundances and the absence of an apparent anticorrelation of Li and O, is an additional argument that large errors in the stellar parameters do not exist for the Orion sample of solar-type stars. The bottom panel of Figure 7 shows distributions of slopes in simulated data for two different models: constant Li with varying O (resulting in model slope distributions near 0.0, as observed in the real data), and Li abundances which decrease as O<sup>-1.0</sup>. The observed Li abundances suggest no significant trend of Li with O, and are most easily fit with the estimated errors of 150K in T<sub>eff</sub>, 0.3 dex in log g, and 0.2 km s<sup>-1</sup> in $`\xi `$. With these errors, the derived anticorrelation of boron with oxygen is real and must be explained.
Of interest to this investigation is a comparison with the interstellar abundances of boron and oxygen in the Association’s gas. Both abundances have been measured along sightlines in Orion but not to stars for which we have B and O abundances. Meyer, Jura, & Cardelli (1998) provided the “definitive” interstellar O abundance including measurements for 5 stars in Orion. From these and other sightlines, the O abundance measured from O i 1356 Å line was remarkably uniform with no evidence for direct condensation of oxygen atoms onto grains. When an estimate of oxygen atoms tied up in grains is added to the measured abundance, the total oxygen abundance was put at log $`ϵ`$(O) = 8.60 with an uncertainty of less than 0.1 dex.
In diffuse interstellar clouds, gaseous boron is expected to be predominantly present as B<sup>+</sup> ions. Detections of these ions through weak absorption at the 1362 Å B ii line are reported by Jura et al. (1996), Lambert et al. (1998), and Howk, Sembach, & Savage (2000) for various sightlines. Abundances, as summarized by Lambert et al. (1998), for sightlines to Orion are log $`ϵ`$(B) $``$ 2.0 when independent observations of the H column density (mostly H i) are considered. This is a lower limit to the interstellar B abundance because some boron atoms may have condensed onto interstellar grains. The combination of measured interstellar abundances in the direction of Orion of log $`ϵ`$(B) $``$ 2.0 and log $`ϵ`$(O) = 8.6 (Meyer et al. 1998) is shown in Figure 4, where it is seen to lie off the trend shown by the Orion stars. Recently, Howk, Sembach, & Savage (2000) have reported a higher interstellar boron abundance, log $`ϵ`$(B) $``$ 2.4. They attribute their higher abundance to observations of clouds of lower density in which depletion of boron onto grains is less severe. However, the newly observed lines of sight are to distant stars, primarily toward the Galactic center. If the diffuse clouds are at a similar distance, their boron abundance could differ from the local (Orion) value on account of Galactic abundance gradients. Therefore, we adopt the local interstellar B abundance but recognize that it is a lower limit.
## 5 Supernovae and the Self-Enrichment of Molecular Clouds
Assuming that the surface composition of a star reflects the composition of the gas from which it formed, the observed variations from star to star within an association can only be understood if the parent cloud was chemically inhomogeneous and/or its composition evolved and all the stars did not form at the same place and time. To interpret the various abundances of oxygen, iron and the light elements in the Orion stars, one therefore has to determine the history and geometry of the chemical enrichment, or if one prefers, the distribution in both space and time of the different nucleosynthetic episodes.
When a stellar association forms from the collapse of a chemically roughly homogeneous cloud, the first-generation stars have approximately the same composition. The most massive stars evolve quickly, on timescales of a few million years, and explode as SN II which release in the ambient medium several solar masses of enriched material, notably more than 1 $`M_{}`$ of pure oxygen per SN II. In addition to this direct contamination of the gas, the SN II’s have a strong dynamical influence on the ambient medium: they produce a shock wave which accelerates particles to relativistic energies and they compress the surrounding gas. Both theoretical models and direct observation indicate that the explosion of a SN II within, or close to, a molecular cloud can trigger the fragmentation of the gas and lead to further star formation. Depending on the mixing of the SN II ejecta with the ambient, chemically unperturbed, gas the new stars formed in the wake of previous SN II can show various O abundances, bounded from below by the initial ISM O abundance and from above by the O abundance in the ejecta. The overabundance of oxygen in some stars of an association can thus be attributed to the *self*-enrichment of the molecular cloud.
Note that in this model the O abundance varies in time, as more and more SN II’s explode and release O-rich material in the ambient medium, but also from one place to another as the hazards of mixing and gas fragmentation dictate. Of course, one cannot expect to model the Orion clouds in sufficient detail to determine the distribution of abundances, density and other physical parameters over the whole region, and we must limit ourselves to general trends and average numbers.
To describe the variation of O abundances in Orion, we divide the gas into two distinct components: the ejecta and the uncontaminated ISM, which we simply call here the *ambient medium*. This simple model considers addition of ejecta of mass $`M_{\mathrm{ej}}`$ to ambient material of mass $`M_{\mathrm{amb}}`$ to provide a total mass $`M_t=M_{\mathrm{amb}}+M_{\mathrm{ej}}`$ from which stars form. If $`\alpha (X)`$ denotes the mass fraction of element $`X`$ and $`f=M_{\mathrm{ej}}/(M_{\mathrm{ej}}+M_{\mathrm{amb}})`$, the composition of a star formed from the mixed gas of mass $`M_t`$ is given by
$$\alpha _{}(X)=(1f)\alpha _{\mathrm{amb}}(X)+f\alpha _{\mathrm{ej}}(X).$$
(2)
The relation between the B and O abundances in Orion can be investigated using the simple model described above, in which new stars form from various amounts of two chemically distinct gas components (the SN II ejecta and the ‘ambient’ medium) before they are fully mixed and their compositions get homogenized. Combining Eq. (2) for O and B, one can eliminate the mixing parameter, $`f`$, and obtain:
$$\alpha _{}(\mathrm{B})=\alpha _{\mathrm{ej}}(\mathrm{B})+K_{\mathrm{B}/\mathrm{O}}[\alpha _{}(O)\alpha _{\mathrm{ej}}(\mathrm{O})],$$
(3)
where
$$K_{\mathrm{B}/\mathrm{O}}=\frac{\alpha _{\mathrm{ej}}(\mathrm{B})\alpha _{\mathrm{amb}}(B)}{\alpha _{\mathrm{ej}}(\mathrm{O})\alpha _{\mathrm{amb}}(O)}.$$
(4)
Alternatively, by exchanging $`f`$ and $`1f`$ as well as $`\alpha _{\mathrm{ej}}(\mathrm{B})`$ and $`\alpha _{\mathrm{amb}}(\mathrm{B})`$ in Eqs. (2) and (3), we have:
$$\alpha _{}(\mathrm{B})=\alpha _{\mathrm{amb}}(\mathrm{B})+K_{\mathrm{B}/\mathrm{O}}[\alpha _{}(O)\alpha _{\mathrm{amb}}(\mathrm{O})],$$
(5)
From Eqs. (3) or (5), we see that the B abundance in Orion stars can be either correlated or anticorrelated with O, depending on whether the slope $`K_{\mathrm{B}/\mathrm{O}}`$ is positive or negative, respectively. The anticorrelation reported in this paper is thus compatible with our simple ‘mixing model’ if $`\alpha _{\mathrm{ej}}(\mathrm{B})<\alpha _{\mathrm{amb}}(\mathrm{B})`$ (since $`\alpha _{\mathrm{ej}}(\mathrm{O})>\alpha _{\mathrm{amb}}(\mathrm{O})`$). This is true if the $`\nu `$-process for <sup>11</sup>B production is negligible, since in that case $`\alpha _{\mathrm{ej}}(\mathrm{B})0`$. Such a conclusion would be very important for the light element nucleosynthesis, since the $`\nu `$-process is generally invoked to increase the B/Be and <sup>11</sup>B/<sup>10</sup>B ratios produced by standard (high energy) nucleo-spallation processes (cf. introduction). If the Orion observations can be used to rule out the $`\nu `$-process, then it seems inevitable that a low-energy cosmic ray component (LECR) exists in the ISM, whose nature and origin remains to be determined.
If we identify the most O-poor (log $`ϵ`$(O) = 8.3) star and the most O-rich (log $`ϵ`$(O) = 9.1) star as having formed from the (initial) ambient and the most severely contaminated mixed gas, respectively, $`\alpha _{}(\mathrm{O})/\alpha _{\mathrm{amb}}(\mathrm{O})`$ is given approximately by the ratio of the observed abundances of the most oxygen-rich to oxygen poor stars ($`10^{0.8}=6.3`$). (The qualification ‘approximately’ is necessary because the equation is framed in terms of mass fractions but the spectroscopic analyses provide the O/H ratio subject to assumptions about the chemical composition, especially about the He/H ratio, and, in the case of the B i line in F-G stars, about the Mg/H ratio (Cunha & Smith 1999).)
In order for B and O to be anticorrelated, the SN II ejecta must have a B/O ratio that is less than the ambient material. The limiting case obviously occurs when the ejecta are thoroughly depleted in B and rich in O, a condition that denies the $`\nu `$-process a significant role in the synthesis of B. In such a case,
$`\alpha _{}(B)=(1f)\alpha _{\mathrm{amb}}(B).`$ (6)
The maximum observed value of $`\alpha _{}(\mathrm{O})/\alpha _{\mathrm{amb}}(\mathrm{O})`$ $``$ 6.3 can be used to set a lower limit on the ratio $`\alpha _{\mathrm{ej}}(\mathrm{O})/\alpha _{\mathrm{amb}}(\mathrm{O})`$ (the lower limit would be this value if the star were composed purely of SN II ejecta). If a ratio $`\alpha _{\mathrm{ej}}(\mathrm{O})/\alpha _{\mathrm{amb}}(\mathrm{O})`$ is assumed, Eq. 2 gives an estimate of $`f`$ that may be used in Eq. 6 to calculate the reduction in the B abundance between the initial ambient gas and the most heavily contaminated star. Obviously, the greater the ratio $`\alpha _{\mathrm{ej}}(\mathrm{O})/\alpha _{\mathrm{amb}}(\mathrm{O})`$, the smaller the reduction of the B abundance. For an arbitrary value of $`\alpha _{\mathrm{ej}}(\mathrm{O})/\alpha _{\mathrm{amb}}(\mathrm{O})`$ = 7 (slightly larger than the lower limit of 6.3) and $`\alpha _{}(\mathrm{O})/\alpha _{\mathrm{amb}}(\mathrm{O})=6`$ (see above), $`f`$ = 0.83 and the observed B abundance in such a star would be $`\alpha _{}(B)=0.17\alpha _{\mathrm{amb}}(B)`$ (a reduction of about 0.8 dex, which is roughly that observed). A value of $`\alpha _{\mathrm{ej}}(\mathrm{O})/\alpha _{\mathrm{amb}}(\mathrm{O})7`$ is not an unreasonable value for a Type II supernova, where 1$`M_{}`$ of oxygen may be synthesized and ejecta may amount to 10$`M_{}`$. In this model, there is no constraint that the B-reduction must approximately equal the O-increase. For example, if $`\alpha _{\mathrm{ej}}(\mathrm{O})/\alpha _{\mathrm{amb}}(\mathrm{O})`$ is raised to 10, the B reduction is only 0.25 dex. This exercise does not address the feasibility of retention of an adequate mass of SN II ejecta by gas that will subsequently form the stars with oxygen abundances above ambient values. This interesting issue which was aired by Cunha & Lambert (1992) may be bypassed here; a more compelling challenge to this simple explanation of the O versus B anticorrelation must be faced.
In the simple mixing model presented here, if B is not synthesized in significant amounts by the $`\nu `$-process, then possibly <sup>7</sup>Li is not produced either. In such a case, the Li and B abundances should be correlated and decline in tandem with increasing O abundance. Lithium is not detectable in B-type stars. Lithium and oxygen abundances of the F-G stars were derived by Cunha et al. (1995, 1998). Figure 7 shows that the lithium abundances are nearly independent of the O abundances, with perhaps a slight decrease. A simple, but by no means unique, interpretation is that synthesis of Li accompanies the O synthesis. However, the interpretation of Li abundances is complicated by the possibility that Li is destroyed within stars. It is important to note that in Figure 7 no undepleted Li abundances are measured below log $`ϵ`$(O)=8.7. There is a tendency for the Orion stars with lower oxygen abundances to have lower Li abundances: these stars are older, tend to rotate more slowly, and have probably depleted their initial Li abundances (Cunha et al. 1995). In order to properly place Li within the context of the evolution of B and O in Orion, it will be necessary to measure the behavior of undepleted Li abundances in Orion members with lower oxygen abundances.
Although the absence of a significant $`\nu `$-process is a sufficient condition for a local anticorrelation between B and O, we draw attention to the fact that it is not a necessary condition. In a model advocated by Parizot (2000), energetic particles (SBEPs) accelerated inside a superbubble created by repeated SN II’s in an OB association may be responsible for synthesis of Li, Be, and B in a supershell. The geometry of the Orion-Eridanus superbubble (see e.g. Burrows et al. 1993) is such that the Orion clouds themselves can actually be considered as part of the supershell. Indeed, the history of star formation in Orion indicates that the molecular cloud is ‘eroded’ from the side which faces the Orion-Eridanus superbubble, and a star formation wave propagates deeper and deeper into the cloud. Now this star-forming side of the cloud is irradiated by SBEPs (Parizot 1998) and an overabundance of spallation products like Be and B is to be expected there. Therefore, the ‘ambient gas’ component in Eq. (2) could be very much enriched in Li, Be and B. This makes it possible that the resulting B abundance, $`\alpha _{\mathrm{amb}}(B)`$, is greater than $`\alpha _{\mathrm{ej}}(B)`$ even in the presence of a significant $`\nu `$-process. From the astrophysical point of view, it depends on the penetration of the SBEPs inside the Orion clouds. The more they penetrate, the less $`\alpha _{\mathrm{amb}}(B)`$, since the B production is then distributed over a larger volume.
## 6 Perspective
As recalled in the introduction, theoretical studies of light element production have up to now confined themselves to the interpretation of the *general* increase of LiBeB abundances as a function of metallicity. On local scales, complex behaviors may result from the fact that CNO and LiBeB are not necessarily produced at exactly the same place or at the same time. Indeed, before the local production of B and O, say, effectively results in a global increase of the B and O abundances in the ISM, the O-rich and B-rich material have to mix together and with the rest of the ISM. Now if new stars form before, or during this mixing episode, one should expect them to show quite unusual compositions. Therefore, it is not obvious that the *local* B-O anticorrelation found in Orion is contradictory with the results obtained on the Galactic scale. In a recent study, Parizot & Drury (2000) have addressed the question of a possible scatter in the Be/O and B/O ratios in the Galaxy, across a mean value accounted for by the superbubble model. The idea was very similar: because O and BeB are not produced (or released) together in superbubbles, various elemental ratios can actually result from inhomogeneous mixing of the O-rich gas with the BeB-rich material. Observational evidence of such a scatter in the Be data is available in Boesgaard et al. (1999).
Detailed calculations of the expected values for $`\alpha _{\mathrm{amb}}(B)`$ and $`\alpha _{\mathrm{ej}}(B)`$, as well as the implications of the Orion data for the SB model will be presented in Paper II. Here, we have noted that the observed B-O anticorrelation can be accounted for within a simple self-enrichment model. First, if the $`\nu `$-process is negligible, the anti-correlation is an inevitable prediction of the model. Second, if the $`\nu `$-process proves to be significant, the observed anticorrelation simply means that $`\alpha _{\mathrm{amb}}(B)>\alpha _{\mathrm{ej}}(B)`$, which is at least plausible considering the very high local production of B from SBEPs (see Paper II for more details and quantitative estimates). Finally, it is important to realize that the Orion observations can actually be used to constrain the light element production models. Although the B-O anticorrelation alone cannot be used to determine the weight of the $`\nu `$-process in the Galaxy, a joint study of the Be abundances should give decisive information. Since Be is not expected to be a $`\nu `$-process product and is then solely a fruit of spallation, it can be used to set the parameters of the SB and the mixing models. Any difference between Be and B can then be attributed to the $`\nu `$-process, and serve to quantify it. We predict that Be will also be found anticorrelated with O in the Orion association, and *even more* than B. This is because B is produced by two means: nucleo-spallation, which leads to the anticorrelation, and neutrino-spallation, which alone would lead to a B-O positive correlation, since the boron produced in this way is ‘ready-mixed’ with the oxygen in the SN II ejecta. Since the $`\nu `$-process is irrelevant for Be, the ‘slope’ of the Be-O anticorrelation must be at least equal, and possibly greater (in absolute value) than that of the B-O anticorrelation. Note that the slopes are here the ones described by Eq. (3), namely $`|K_{\mathrm{B}/\mathrm{O}}|`$ and $`|K_{\mathrm{Be}/\mathrm{O}}|`$ (which should actually be normalized by dividing them by $`(\mathrm{B}/\mathrm{O})_{}`$, say, and $`(\mathrm{Be}/\mathrm{O})_{}`$, respectively). This is different from the slope seen in Figure 4, where the abundances are plotted on logarithmic scales. It is remarkable, however, that the only anticorrelation in logarithmic variables which is compatible with a anticorrelation in linear variables (as predicted by the mixing model, Eq. (3), is one of slope $`1`$, which is what we observe for Orion.
Finally, the case of lithium is more complicated, since the Galactic evolution models indicate that the spallation processes may not be the main contributors, and AGB stars probably produced most of the Li currently observed in the Galaxy. The absence of a clear correlation or anticorrelation found in Figure 7 could thus result from the fact that the $`\nu `$-process is roughly balanced by the SBEP-induced production, or by the fact that none of these processes modifies significantly the Li abundance in the ejecta and the ambient medium.
We thank Ivo Busko for help in issues related to the STIS spectra. This research is supported in part by NASA through the contract NAG5-1616, and the grant GO-06520.01.95A from the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS5-26555. We also acknowledge support from the National Science Foundation through grant AST96-18459. EP was supported by the TMR programme of the European Union under contract FMRX-CT98-0168.
. TABLE 1 Stellar Parameters and Abundances Star T<sub>eff</sub>(K) Log g Log $`ϵ`$(Li)<sub>nlte</sub> Log $`ϵ`$(B)<sub>lte</sub> Log $`ϵ`$(B)<sub>nlte</sub> Log $`ϵ`$(O)<sub>nlte</sub> Log $`ϵ`$(Fe)<sub>lte</sub> P1179 6050 4.0 3.13 8.71 P1374 6390 4.0 3.24 9.00 7.46 P1455 5950 4.0 3.26 8.86 7.59 P1657 6100 3.8 2.56 8.72 7.20 P1789 6120 4.0 3.06 9.29 7.31 P1955 5890 4.0 3.18 9.13 7.53 P2374 5900 3.9 2.30 8.78 7.41 BD-5 1317 5850 4.0 3.11 2.0 2.06 9.11 7.34 BD-6 1250 5950 4.0 2.74 2.3 2.34 8.74 7.33 HD294297 6150 4.4 2.56 2.6 2.65 8.61 7.32
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# PARAMETRIZATION OF POLARIZED PARTON DISTRIBUTIONS
## 1 Introduction
From the measurements of the proton’s polarized structure function $`g_1`$, we learned that the proton spin cannot be understood in a simple quark model. It should be interpreted mainly by a gluon contribution and/or effects of angular momenta; however, the situation is not satisfactory for specifying the major carrier of the proton spin. Our study is intended to understand the current status of polarized parton distributions by analyzing inclusive spin asymmetry $`A_1`$ data prior to RHIC-Spin measurements. Semi-inclusive data became available, but they are not accurate enough to provide any significant constraint. Our group name AAC stands for Asymmetry Analysis Collaboration, which consists of theorists and experimentalists.$`^\mathrm{?}`$ In this paper, we report our analysis of the asymmetry $`A_1`$ for getting optimum polarized parton distributions.
In Sec.2, a parametrization form of the polarized parton distributions is discussed. The results are shown in Sec.3. According to our analysis, the polarized antiquark distribution cannot be well determined, and this issue is discussed in Sec.4. The results in Sec.3 and 4 are quoted from Ref.1. Finally, conclusions are given in Sec.5.
## 2 Parametrization
First, an $`x`$-dependent functional form of parametrized distributions is explained. Because the positivity condition is imposed not only in the leading order (LO) but also in the next-to-leading order (NLO), it is convenient to have an unpolarized distribution multiplied by an $`x`$-dependent function:
$$\mathrm{\Delta }f_i(x,Q_0^2)=A_ix^{\alpha _i}(1+\gamma _ix^{\lambda _i})f_i(x,Q_0^2),$$
(1)
at the initial $`Q^2`$ point ($`Q_0^2`$=1 GeV<sup>2</sup>). Since there is no accurate data to find the difference between $`\mathrm{\Delta }\overline{u}`$, $`\mathrm{\Delta }\overline{d}`$, and $`\mathrm{\Delta }\overline{s}`$, these distributions are assumed to be the same. Of course, they are not expected to be equal by considering the unpolarized situation.$`^\mathrm{?}`$ With this assumption and the flavor number $`N_f=3`$, we have four distributions ($`\mathrm{\Delta }u_v`$, $`\mathrm{\Delta }d_v`$, $`\mathrm{\Delta }\overline{q}`$, $`\mathrm{\Delta }g`$) to be determined by a $`\chi ^2`$ analysis. Because there are four parameters for each distribution, we have sixteen parameters in total. However, the first moments of $`\mathrm{\Delta }u_v`$ and $`\mathrm{\Delta }d_v`$ are fixed ($`\eta _{u_v}=0.926`$ and $`\eta _{d_v}=0.341`$) by using semi-leptonic decay constants. It means that the number of actual free parameters is fourteen.
The polarized distributions are evolved $`^\mathrm{?}`$ to experimental $`Q^2`$ points of the spin asymmetry $`A_1`$, which is expressed by the structure functions as
$$A_1(x,Q^2)\frac{2xg_1(x,Q^2)[1+R(x,Q^2)]}{F_2(x,Q^2)}.$$
(2)
The longitudinal-transverse structure function ratio $`R`$ is taken from the analysis of SLAC-1990, $`F_2`$ is calculated by the GRV98 distributions, and $`g_1`$ is calculated by our parametrized distributions. Then, $`\chi ^2`$ is evaluated in comparison with experimental data: $`\chi ^2=(A_1^{\mathrm{data}}(x,Q^2)A_1^{\mathrm{calc}}(x,Q^2))^2/(\mathrm{\Delta }A_1^{\mathrm{data}}(x,Q^2))^2`$. The optimum set of parameters is found by minimizing $`\chi ^2`$ by the subroutine MINUIT.
## 3 Results
In our analysis, we use the asymmetry data set of so called “large tables”. On the other hand, other analyses use small data tables in which $`Q^2`$ values are averaged at certain $`x`$ points. Although the present data are not accurate enough to provide $`Q^2`$-dependence information, we believe that the raw data should be used as much as we can. For 375 data points with $`Q^2>1`$ GeV<sup>2</sup>, we obtain the minimum $`\chi ^2`$ values $`\chi ^2`$/d.o.f.=322.6/360 and 300.4/360 for LO and NLO, respectively. There is a significant $`\chi ^2`$ reduction in NLO, so that it is important to analyze the data by the NLO expressions. It should be noted that all our NLO results are obtained in the $`\overline{\mathrm{MS}}`$ scheme. There are two major sources for the $`\chi ^2`$ reduction, and they are from the HERMES proton and E154 neutron data in Figs.2 and 2.
In these figures, our optimum LO and NLO asymmetry curves are shown at $`Q^2`$=5 GeV<sup>2</sup>. Because the data are taken at various $`Q^2`$ points depending on the $`x`$ region, it is not straightforward to compare them with the fixed $`Q^2`$ curves. However, the figures show that the agreement is satisfactory. The errors of deuteron data are so large that its $`A_1`$ figure is not shown here.
Obtained polarized parton distributions are shown in Fig.4. Because the first moments of $`\mathrm{\Delta }u_v`$ and $`\mathrm{\Delta }d_v`$ are fixed at positive and negative numbers, respectively, they are mainly positive and negative distributions. It is rather difficult to determine the gluon distribution, but a positive distribution is favored in both LO and NLO. The antiquark distribution becomes negative, and we discuss the issue of its determination in the next section.
Using the obtained distributions, we show $`Q^2`$ dependence of $`A_1`$ in comparison with data in Fig.4. The distributions are provided at $`Q^2`$=1 GeV<sup>2</sup> and they are evolved to larger $`Q^2`$ by the DGLAP equations. The curves show that there could be strong $`Q^2`$ dependence in the small-$`Q^2`$ region ($`Q^2<`$2 GeV<sup>2</sup>). Experimentalists sometimes assume that $`A_1`$ is $`Q^2`$ independent for evaluating $`g_1`$ from the $`A_1`$ data. Although it does not matter due to the present experimental accuracy, we should be careful about such an assumption for a precise analysis.
Using these distributions, we get the quark spin content
$$\mathrm{\Delta }\mathrm{\Sigma }=0.201(\mathrm{LO}),0.051(\mathrm{NLO}).$$
(3)
The NLO value seems to be significantly smaller than usually-quoted ones 10%$``$30%. The small spin content originates from the small-$`x`$ behavior of our antiquark distribution, so that we discuss this issue in the next section.
## 4 Polarized antiquark distribution
Because the small-$`x`$ extrapolation could affect the value of the spin content, we compare our antiquark distribution with some of the recent analyses. In Fig.5, our AAC curve is shown together with the LSS (Leader-Sidorov-Stamenov) and SMC (Spin Muon Collaboration) distributions. The antiquark distribution is not explicitly calculated in the SMC paper, so that the curve is obtained by transforming their distributions. The SMC curve deviates from others in the large-$`x`$ region. However, the large-$`x`$ difference does not matter at this stage because the antiquark distribution does not contribute to $`g_1`$ at large $`x`$. It should be clarified for example by the Drell-Yan measurements at RHIC in any case. Another important point is in the small-$`x`$ region. Our distribution falls off rather slowly in comparison with others, which results in the small quark spin content. Unfortunately, available data are taken in the region $`x>0.004`$, and $`x`$ is not small enough to determine $`\mathrm{\Delta }\overline{q}`$ uniquely. It is also still too far away from meaningful flavor decomposition for $`\mathrm{\Delta }\overline{q}`$.$`^\mathrm{?}`$
Due to the lack of small-$`x`$ data, we had better fix the small-$`x`$ behavior of $`\mathrm{\Delta }\overline{q}`$. As theoretical guidelines, the Regge theory and perturbative QCD could be used. First, the Regge model predicts $`g_1(x)x^\alpha `$ as $`x0`$ with the intercepts $`\alpha `$ of $`a_1(1260)`$, $`f_1(1285)`$, and $`f_1(1420)`$ trajectories. However, the intercepts are not well known and it is usually assumed as $`\alpha _{a_1}=0.50`$. The small-$`x`$ functional form of the GRV98 is $`x\overline{q}x^{0.14}`$ at $`Q^2=1`$ GeV<sup>2</sup> according to our numerical estimate, so that the Regge prediction is $`\mathrm{\Delta }\overline{q}/\overline{q}x^{1.11.6}`$. Second, the perturbative QCD could also suggest the small-$`x`$ behavior. However, we have to assume that singlet and gluon distributions are constants as $`x0`$ at certain $`Q^2(Q_1^2)`$. Although it is not obvious whether such $`Q_1^2`$ exists, let us assume $`Q_1^2=0.30.5`$ GeV<sup>2</sup>. Then, the singlet distribution becomes $`x^{0.120.09}`$ as $`x0`$, hence $`\mathrm{\Delta }\overline{q}/\overline{q}x^{1.0}`$. The Regge and “pQCD” distributions fall off much faster than ours at small $`x`$.
From these theoretical suggestions, we analyzed the data again by fixing the parameter $`\alpha _{\overline{q}}`$ ($`\mathrm{\Delta }\overline{q}/\overline{q}x^{\alpha _{\overline{q}}}`$). First, we fixed the parameter at 1.0 which is suggested by pQCD and it is also about the lower bound of the Regge model. Second, it is taken as 1.6 which is the upper bound of the Regge. Then, we obtain the NLO results:
$`\chi ^2=`$ $`305.8,`$ $`\mathrm{\Delta }\mathrm{\Sigma }=`$ $`0.241,`$ $`\mathrm{for}\alpha _{\overline{q}}=`$ $`1.0,`$
$`323.5,`$ $`0.276,`$ $`1.6.`$ (4)
The $`\chi ^2`$ change for $`\alpha _{\overline{q}}`$=1.0 from the previous NLO value (300.4) is merely 1.8%, but we notice the large change in the spin content (5%$``$24%). Because of the small $`\chi ^2`$ change, the $`\alpha _{\overline{q}}`$=1.0 fit could be equally taken as a good solution. On the other hand, the $`\chi ^2`$ change is rather large for $`\alpha _{\overline{q}}`$=1.6. According to these results, the spin content is 24%$``$28% which is in the range of the widely-quoted values. In this way, we find that the spin content is very sensitive to the small-$`x`$ behavior of the polarized antiquark distribution and that it cannot be fixed by the present $`g_1`$ data.
## 5 Conclusions
It is clarified in our $`\chi ^2`$ analysis of the $`A_1`$ data that the polarized antiquark distribution is not well determined in the small- and large-$`x`$ regions. This fact leads to the conclusion that the quark spin content is not well constrained only by the present $`A_1`$ data. We need small- and large-$`x`$ measurements in future. From our studies, we have proposed three sets of distributions: LO, NLO-1 with free $`\alpha _{\overline{q}}`$, and NLO-2 with fixed $`\alpha _{\overline{q}}=1.0`$.$`^\mathrm{?}`$
## Acknowledgments
S.K. was partly supported by the Grant-in-Aid for Scientific Research from the Japanese Ministry of Education, Science, and Culture under the contract number 10640277. He would like to thank the theory division of CERN, where this manuscript is partially written, for supporting his stay.
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# On the asymptotic geometry of abelian-by-cyclic groupsTo appear in Acta Mathematica
## 1 Introduction
Gromov’s Polynomial Growth Theorem \[Gro81\] states that the property of having polynomial growth characterizes virtually nilpotent groups among all finitely generated groups.
Gromov’s theorem inspired the more general problem (see, e.g. \[GdlH91\]) of understanding to what extent the asymptotic geometry of a finitely-generated solvable group determines its algebraic structure—in short, are solvable groups quasi-isometrically rigid? In general they aren’t: very recently A. Dioubina \[Dio99\] has found a solvable group which is quasi-isometric to a group which is not virtually solvable; these groups are finitely generated but not finitely presentable. In the opposite direction, first steps in identifying quasi-isometrically rigid solvable groups which are not virtually nilpotent were taken for a special class of examples, the solvable Baumslag-Solitar groups, in \[FM98\] and \[FM99b\].
The goal of the present paper is to show that a much broader class of solvable groups, the class of finitely-presented, nonpolycyclic, abelian-by-cyclic groups, is characterized among all finitely-generated groups by its quasi-isometry type. We also give a complete quasi-isometry classification of the groups in this class; such a classification for nilpotent groups remains a major open question. Motivated by these results, we offer a conjectural picture of quasi-isometric classification and rigidity for polycyclic abelian-by-cyclic groups in §10.1.
The proofs of these results lead one naturally from a geometry of groups problem to the theory of dynamical systems via the asymptotic behavior of certain flows and their associated foliations.
### 1.1 The abelian-by-cyclic group $`\mathrm{\Gamma }_M`$
A group $`\mathrm{\Gamma }`$ is abelian-by-cyclic if there is an exact sequence
$$1A\mathrm{\Gamma }Z1$$
where $`A`$ is an abelian group and $`Z`$ is an infinite cyclic group. If $`\mathrm{\Gamma }`$ is finitely generated, then $`A`$ is a finitely generated module over the group ring $`𝐙[Z]`$, although $`A`$ may not be finitely generated as a group.
By a result of Bieri and Strebel \[BS78\], the class of finitely presented, torsion-free, abelian-by-cyclic groups may be described in another way. Consider an $`n\times n`$ matrix $`M`$ with integral entries and $`detM0`$. Let $`\mathrm{\Gamma }_M`$ be the ascending HNN extension of $`𝐙^n`$ given by the monomorphism $`\varphi _M`$ with matrix $`M`$. Then $`\mathrm{\Gamma }_M`$ has a finite presentation
$$t,a_1,\mathrm{},a_n|[a_i,a_j]=1,ta_it^1=\varphi _M(a_i),i,j=1,\mathrm{},n$$
where $`\varphi _M(a_i)`$ is the word $`a_1^{m_1}\mathrm{}a_n^{m_n}`$ and the vector $`(m_1,\mathrm{},m_n)`$ is the $`i^{\text{th}}`$ column of the matrix $`M`$. Such groups $`\mathrm{\Gamma }_M`$ are precisely the class of finitely presented, torsion-free, abelian-by-cyclic groups (see \[BS78\] for a proof involving a precursor of the Bieri-Neumann-Strebel invariant, or \[FM99b\] for a proof using trees). The group $`\mathrm{\Gamma }_M`$ is polycyclic if and only if $`\left|detM\right|=1`$; this is easy to see directly, and also follows from \[BS80\].
### 1.2 Statement of results
The first main theorem in this paper gives a classification of all finitely-presented, nonpolycyclic, abelian-by-cyclic groups up to quasi-isometry. It is easy to see that any such group has a torsion-free subgroup of finite index, and so is commensurable (hence quasi-isometric) to some $`\mathrm{\Gamma }_M`$. The classification of these groups is actually quite delicate—the standard quasi-isometry invariants (ends, growth, isoperimetric inequalities, etc.) do not distinguish any of these groups from each other, except that the size of the matrix $`M`$ can be detected by large scale cohomological invariants of $`\mathrm{\Gamma }_M`$.
Given $`M\mathrm{GL}(n,𝐑)`$, the *absolute Jordan form* of $`M`$ is the matrix obtained from the Jordan form for $`M`$ over $`𝐂`$ by replacing each diagonal entry with its absolute value, and rearranging the Jordan blocks in some canonical order.
###### Theorem 1.1 (Classification theorem).
Let $`M_1`$ and $`M_2`$ be integral matrices with $`\left|detM_i\right|>1`$ for $`i=1,2`$. Then $`\mathrm{\Gamma }_{M_1}`$ is quasi-isometric to $`\mathrm{\Gamma }_{M_2}`$ if and only if there are positive integers $`r_1,r_2`$ such that $`M_1^{r_1}`$ and $`M_2^{r_2}`$ have the same absolute Jordan form.
###### Remark.
Theorem 1.1 generalizes the main result of \[FM98\], which is the case when $`M_1,M_2`$ are positive $`1\times 1`$ matrices; in that case the result of \[FM98\] says even more, namely that $`\mathrm{\Gamma }_{M_1}`$ and $`\mathrm{\Gamma }_{M_2}`$ are quasi-isometric if and only if they are commensurable. When $`n2`$, however, it’s not hard to find $`n\times n`$ matrices $`M_1,M_2`$ such that $`\mathrm{\Gamma }_{M_1},\mathrm{\Gamma }_{M_2}`$ are quasi-isometric but not commensurable. Polycyclic examples are given in \[BG96\], and the same ideas may be used to produce nonpolycyclic examples.
The following theorem shows that the algebraic property of being a finitely-presented, nonpolycyclic, abelian-by-cyclic group is in fact a large-scale geometric property.
###### Theorem 1.2 (Quasi-isometric rigidity).
Let $`\mathrm{\Gamma }=\mathrm{\Gamma }_M`$ be a finitely presented abelian-by-cyclic group, determined by an integer $`(n\times n)`$ matrix $`M`$ with $`\left|detM\right|>1`$. Let $`G`$ be any finitely generated group quasi-isometric to $`\mathrm{\Gamma }`$. Then there is a finite normal subgroup $`KG`$ such that $`G/K`$ is abstractly commensurable to $`\mathrm{\Gamma }_N`$, for some integer $`(n\times n)`$ matrix $`N`$ with $`\left|detN\right|>1`$.
###### Remark.
Theorem 1.2 generalizes the main result of \[FM99b\], which covers the case when $`M`$ is a positive $`1\times 1`$ matrix. The latter result was given a new proof in \[MSW\], and in §9 we follow the methods of \[MSW\] in proving Theorem 1.2.
###### Remark.
The “finitely presented” hypothesis in Theorem 1.2 cannot be weakened to “finitely generated”. Dioubina shows \[Dio99\] that the wreath product $`𝐙\text{w}r𝐙`$, an abelian-by-cyclic group of the form $`𝐙[𝐙]`$-by-$`𝐙`$, is quasi-isometric to the wreath product $`(𝐙F)\text{w}r𝐙`$ whenever $`F`$ is a finite group. But $`(𝐙F)\text{w}r𝐙`$ has no nontrivial finite normal subgroups, and when $`F`$ is nonabelian it is not abstractly commensurable to an abelian-by-cyclic group.
One of the key technical results used to prove Theorem 1.1 is the following theorem, which we believe is of independent interest. It describes a rigidity phenomenon for 1-parameter subgroups of $`\mathrm{GL}(n,𝐑)`$ which generalizes work of Benardete \[Ben88\] (see also \[Wit90\]).
A 1-parameter subgroup $`M^t`$ of $`\mathrm{GL}(n,𝐑)`$ determines a 1-parameter family of quadratic forms $`Q_M(t)=(M^t)^T(M^t)`$ on $`𝐑^n`$, where the superscript <sup>T</sup> denotes transpose. Each $`Q_M(t)`$ determines a norm $`_{M,t}`$ and a distance function $`d_{M,t}`$ on $`𝐑^n`$.
Theorem 5.11 (One-parameter subgroup rigidity) Let $`M^t,N^t`$ be 1-parameter subgroups of $`\mathrm{GL}(n,𝐑)`$, such that $`M=M^1`$ and $`N=N^1`$ have no eigenvalues on the unit circle. If there exists a bijection $`f:𝐑^n𝐑^n`$ and constants $`K1,C0`$ such that for each $`t𝐑`$ and $`p,q𝐑^n`$
$$C+\frac{1}{K}d_{M,t}(p,q)d_{N,t}(f(p),f(q))Kd_{M,t}(p,q)+C$$
then $`M`$ and $`N`$ have the same absolute Jordan form.
The proof of Theorem 5.11 is given in §6, and shows that in fact $`f`$ is a homeomorphism with a reasonably high degree of regularity; see Proposition 6.3.
###### Remark.
The case of Theorem 5.11 when $`f`$ is the identity map follows from a theorem of D. Benardete \[Ben88\]. See also D. Witte \[Wit90\]. Benardete’s theorem determines precisely when two one-parameter subgroups of $`\mathrm{GL}(n,𝐑)`$ diverge, and it applies as well to matrices with eigenvalues on the unit circle.
### 1.3 Homogeneous spaces
Using coarse topological and geometrical methods, we reduce the study of quasi-isometries of $`\mathrm{\Gamma }_M`$ to that of a certain Lie group $`G_M`$.
After squaring $`M`$ if necessary, we can assume that $`detM>0`$ and that $`M`$ lies on a 1-parameter subgroup $`M^t`$ of $`\mathrm{GL}(n,𝐑)`$. The group $`\mathrm{\Gamma }_M`$ is a cocompact subgroup of the solvable Lie group $`G_M=𝐑^n_M𝐑`$, where $`𝐑`$ acts on $`𝐑^n`$ by the $`1`$-parameter subgroup $`M^t`$. The group $`\mathrm{\Gamma }_M`$ is discrete in $`G_M`$ if and only if $`detM=1`$. See section 4 for details.
The groups $`G_M`$, with their left invariant metrics, give a rich and familiar collection of examples, including: all real hyperbolic spaces, when $`M`$ is a constant times the identity; many negatively curved homogeneous spaces, when $`M`$ has all eigenvalues $`>1`$ in absolute value; and 3-dimensional solv geometry, when $`M`$ is a $`2\times 2`$ hyperbolic matrix of determinant $`1`$. The negatively curved examples associated to a real diagonal matrix with all eigenvalues $`>1`$ were studied by Pansu \[Pan89a\] (and later Gromov \[Gro93\]), who analyzed their quasi-isometric geometry using the idea of “conformal dimension”.
We should mention also the result of Heintze \[Hei74\] that the class of connected, negatively curved homogeneous spaces consists precisely of those spaces of the form $`N𝐑`$ where $`N`$ is a nilpotent Lie group, and the action of $`𝐑`$ on the Lie algebra has all eigenvalues strictly outside the unit circle.
### 1.4 Outline of proofs
After preliminary sections, §3 on Linear Algebra, and §4 on The Solvable Lie Group $`G_M`$, the proof of Theorem 1.1 can be divided into 3 main parts: §5,6 on the Dynamics of $`G_M`$; §7 on Quasi-Isometries of $`\mathrm{\Gamma }_M`$ via Coarse Topology; and §8 on Finding the Integers, where the pieces of the proof are put together. The proof of Theorem 1.2 is contained in §9 on Quasi-Isometric Rigidity. Finally we pose some conjectures and problems in §10.
### §5,6: Dynamics of $`G_M`$
In these two sections we classify the Lie groups $`G_M`$ up to horizontal-respecting quasi-isometry, that is up to quasi-isometries $`\varphi :G_MG_N`$ which take each set of the form $`𝐑^m\times \{t\}`$ to a set of the form $`𝐑^n\times \{h(t)\}`$ for some function $`h`$ called the *induced time change*.
Theorem 5.2 (Horizontal respecting quasi-isometries: special case) Let $`M,N`$ lie on 1-parameter subgroups $`M^t,N^t`$ of $`\mathrm{GL}(n,𝐑)`$, and suppose that $`detM,detN>1`$. If there exists a horizontal respecting quasi-isometry $`\varphi :G_MG_N`$, then there exist real numbers $`r,s>0`$ so that $`M^r`$ and $`N^s`$ have the same absolute Jordan form.
###### Remark.
In the special case where $`M,N`$ are diagonalizable with all eigenvalues $`>1`$, this can be extracted from work of Pansu \[Pan89a\] *without* the assumption that $`\varphi `$ is horizontal respecting. This special case was later reconsidered by Gromov (see \[Gro93\] Section 7.C), as an application of his “inf$`\delta `$im” invariant. Our statement and proof of Theorem 5.2 is inspired in part by the ideas of exponential growth rates built into the inf$`\delta `$im invariant (see also comments after Proposition 5.8).
In §5 we give a slightly more general version of this statement, Theorem 5.2.
The proof of Theorem 5.2 uses a certain dynamical system on $`G_M`$, the “vertical flow” which flows upward at unit speed along flow lines of the form (point)$`\times 𝐑𝐑^m_M𝐑`$. When $`M`$ has no eigenvalues on the unit circle this is a hyperbolic or Anosov flow, and in general it is a partially hyperbolic flow. We prove Theorem 5.2 in several steps, using stronger and stronger dynamical properties of flows in $`G_M`$.
##### Step 1 (Foliations Rigidity, Proposition 5.4)
Using the Shadowing Lemma from hyperbolic dynamics we show that $`\varphi `$ coarsely respects three dynamically defined foliations of $`G_M`$ and $`G_N`$: the weak stable, weak unstable, and center foliations. This, together with a result of Bridson-Gersten that depends in turn on work of Pansu (see Corollary 5.6), allows reduction to the case where $`M,N`$ have no eigenvalues on the unit circle.
##### Step 2 (Time rigidity, Proposition 5.8)
We show that the induced time change map of $`\varphi `$ is actually an *affine* map between the time parameters of $`G_M`$ and $`G_N`$. After taking a real power of $`N`$ and composing with a vertical translation, we can assume that $`\varphi `$ preserves the time parameter, that is $`h(t)=t`$.
##### Step 3 (One-parameter subgroup rigidity, Theorem 5.11)
From Step 2, $`\varphi `$ induces a quasi-isometry between corresponding level sets of the time parameter on $`G_M,G_N`$, which reduces the proof to Theorem 5.11, One-Parameter Subgroup Rigidity. The latter theorem is proved in §6, by studying rigidity properties of certain flags of foliations of $`𝐑^n`$ associated to the absolute Jordan form of $`M\mathrm{GL}(n,𝐑)`$.
### §7: Quasi-Isometries of $`\mathrm{\Gamma }_M`$ via Coarse Topology
Given an integer matrix $`M\mathrm{GL}(n,𝐑)`$ with $`detM>1`$, we study the geometry of $`\mathrm{\Gamma }_M`$ by constructing a contractible metric cell complex $`X_M`$ on which $`\mathrm{\Gamma }_M`$ acts freely, properly discontinuously and cocompactly by isometries, so that $`\mathrm{\Gamma }_M`$ is quasi-isometric to $`X_M`$. Topologically, $`X_M`$ is a product of $`𝐑^m`$ with the homogeneous directed tree $`T_M`$ with one edge entering and $`d`$ edges leaving each vertex. Here $`d=detM`$. Metrically, for every coherently oriented line $`\mathrm{}`$ in $`T_M`$, the metric on $`X_M`$ is such that $`𝐑^m\times \mathrm{}`$ is isometric to $`G_M`$.
The main result of this section, Proposition 7.1, says that a quasi-isometry $`f:X_MX_N`$ induces a quasi-isometry $`\varphi :G_MG_N`$ which respects horizontal foliations. This is proved using coarse geometric and topological methods. This is precisely where the condition $`detM,detN>1`$ is essential for our proof, since it gives that the trees $`T_M,T_N`$ have nontrivial branching, and this branching allows us to show that $`f`$ “remembers” the branch points (see Step 2 of §7.2).
While this proof is in the spirit of \[FM98\], further complications arise in this more general case (see §7.2). Also, for other applications (e.g. \[FM99a\], \[MSW\]), we shall derive Proposition 7.1 from a still more general result, Theorem 7.7, which applies to many graphs of groups whose vertex and edge groups are fundamental groups of aspherical manifolds of fixed dimension.
### §8: Finding the integers
Given integer matrices $`M,N\mathrm{GL}(n,𝐑)`$ with $`\left|detM\right|,\left|detN\right|>1`$ such that $`\mathrm{\Gamma }_M`$ and $`\mathrm{\Gamma }_N`$ are quasi-isometric, a simple argument allows us to reduce to the case of positive determinant, and then the results of §57 combine to show that there are positive real numbers $`r,s`$ so that $`M^r`$ and $`N^s`$ have the same absolute Jordan form. We need to show that integral $`r,s`$ exist. This is done by showing that a quasi-isometry $`X_MX_N`$ induces a bilipschitz homeomorphism between certain self-similar Cantor sets attached to $`X_M`$ and $`X_N`$. Applying a theorem of Cooper on bilipschitz types of these Cantor sets allows us to conclude that $`(detM)^p=(detN)^q`$ for some integers $`p,q1`$, from which the desired conclusion follows.
### §9: Quasi-Isometric Rigidity
To prove Theorem 1.2, we use the coarse topology results from §7 to show that a group quasi-isometric to some $`\mathrm{\Gamma }_M`$ admits a quasi-action on a tree of $`n`$-dimensional Euclidean spaces. We then use the results of \[MSW\] to convert this quasi-action into a true action on a tree, whose edge and vertex stabilizers are finitely generated groups quasi-isometric to $`𝐙^n`$. The proof is completed by invoking well-known quasi-isometry invariants, combined with a brief study of injective endomorphisms of virtually abelian groups.
### §10: Questions
We make some conjectures concerning possible extensions of this work to the polycyclic case. Also, we state some problems on the quasi-isometry group of $`\mathrm{\Gamma }_M`$.
##### Acknowledgements
We thank Kevin Whyte and Amie Wilkinson for all their help. We are also grateful to the IHES, where much of this work was done.
## 2 Preliminaries
This brief section reviews some basic material; see for example \[GdlH91\].
Given $`K1,C0`$, a $`(K,C)`$ *quasi-isometry* between metric spaces is a map $`f:XY`$ such that:
1. For all $`x,x^{}X`$ we have
$$\frac{1}{K}d_X(x,x^{})Cd_Y(f(x),f(x^{}))Kd_X(x,x^{})+C$$
2. For all $`yY`$ we have $`d_Y(y,f(X))C`$.
If $`f`$ satisfies (1) but not necessarily (2) then it is called a $`(K,C)`$ *quasi-isometric embedding*. If $`f`$ satisfies only the right hand inequality of (1) then $`f`$ is $`(K,C)`$ *coarsely lipschitz*, and if in addition $`C=0`$ then $`f`$ is $`K`$-*lipschitz*.
A *coarse inverse* of a quasi-isometry $`f:XY`$ is a quasi-isometry $`g:YX`$ such that, for some constant $`C^{}>0`$, we have $`d(gf(x),x)<C^{}`$ and $`d(fg(y),y)<C^{}`$ for all $`xX`$ and $`yY`$. Every $`(K,C)`$ quasi-isometry $`f:XY`$ has a $`K,C^{}`$ coarse inverse $`g:YX`$, where $`C^{}`$ depends only on $`K,C`$: for each $`yY`$ define $`g(y)`$ to be any point $`xX`$ such that $`d(f(x),y)C`$.
A fundamental fact observed by Efremovich, by Milnor \[Mil68\], and by Švarc, which we use repeatedly without mentioning, states that if a group $`G`$ acts properly discontinuously and cocompactly by isometries on a proper geodesic metric space $`X`$, then $`G`$ is finitely generated, and $`X`$ is quasi-isometric to $`G`$ equipped with the word metric.
Given a metric space $`X`$ and $`A,BX`$, we denote the *Hausdorff distance* by
$$d_{}(A,B)=inf\{r[0,\mathrm{}]|AN_r(B)\text{and}BN_r(A)\}$$
The following lemma says that an ambient quasi-isometry induces a quasi-isometry between subspaces of a certain type. A map $`\sigma :SX`$ between geodesic metric spaces is *uniformly proper* if there is function $`\rho :[0,\mathrm{})[0,\mathrm{})`$ with $`\underset{t\mathrm{}}{lim}\rho (t)=+\mathrm{}`$, and constants $`K1,C0`$, such that for all $`x,yS`$ we have:
$$\rho \left(d_S(x,y)\right)d_X(\sigma (x),\sigma (y))Kd_S(x,y)+C$$
The function $`\rho `$ and the constants $`K,C`$ are called *uniformity data* for $`\sigma `$.
###### Lemma 2.1.
Given geodesic metric spaces $`X,Y,S,T`$, a quasi-isometry $`f:XY`$, and uniformly proper maps $`\sigma :SX`$ and $`\tau :TY`$, suppose that $`d_{}(f\sigma (S),\tau (T))<\mathrm{}`$. Then $`S,T`$ are quasi-isometric. To be explicit, any function $`g:ST`$ such that $`d_Y(f\sigma (x),\tau g(x))`$ is uniformly bounded is a quasi-isometry; the quasi-isometry constants for $`g`$ depend only on those for $`f`$, the uniformity data for $`\sigma `$ and $`\tau `$, and the bound for $`d_Y(f\sigma (x),g\tau (x))`$.
###### Proof.
Pick $`K1`$, $`C0`$ and $`\rho :[0,\mathrm{})[0,\mathrm{})`$ such that $`f`$ is a $`(K,C)`$ quasi-isometry, $`d_Y(f\sigma (x),g\tau (x))C`$, and $`\rho ,K,C`$ are uniformity data for $`\sigma ,\tau `$.
Consider $`x,yS`$ such that $`d_S(x,y)1`$. We have $`d_Y(f\sigma (x),f\sigma (y))K^2+KC+C`$, and so $`d_Y(\tau g(x),\tau g(y))K^2+KC+3C`$, from which it follows that $`\rho \left(d_T(g(x),g(y))\right)K^2+KC+3C`$. Since $`\underset{t\mathrm{}}{lim}\rho (t)=\mathrm{}`$ we obtain a bound $`d_T(g(x),g(y))A`$ depending only on $`K,C,\rho `$. The usual “rubber band” argument, using geodesics in $`S`$ divided into subsegments of length $`1`$ with a terminal subsegment of length $`1`$, suffices to prove that $`g`$ is $`(K^{},C^{})`$ coarsely lipschitz, with $`K^{},C^{}`$ depending only on $`K,C,\rho `$.
For any $`\xi T`$ there is a point $`\overline{g}(\xi )S`$ such that $`d_Y(f\sigma \overline{g}(\xi ),\tau (\xi ))C`$. For any $`\xi ,\eta T`$ with $`d(\xi ,\eta )1`$ we have
$`d_Y(f\sigma \overline{g}(\xi ),f\sigma \overline{g}(\eta ))`$ $`d_Y(f\sigma \overline{g}(\xi ),\tau (\xi ))+d_Y(\tau (\xi ),\tau (\eta ))+d_Y(f\sigma \overline{g}(\eta ),\tau (\eta ))`$
$`K+3C`$
and so $`\rho \left(d_S(\overline{g}(\xi ),\overline{g}(\eta ))\right)d_X(\sigma \overline{g}(\xi ),\sigma \overline{g}(\eta ))K^2+4KC`$. As above we obtain an upper bound for $`d_S(\overline{g}(\xi ),\overline{g}(\eta ))`$ and the rubber band argument shows that $`\overline{g}`$ is coarsely lipschitz.
For any $`xS`$, setting $`\xi =g(x)T`$, we have
$`d_Y(f\sigma (x),f\sigma \overline{g}(\xi ))`$ $`d_Y(f\sigma (x),\tau g(x))+d_Y(\tau (\xi ),f\sigma \overline{g}(\xi ))`$
$`2C`$
It follows that $`d_X(\sigma (x),\sigma \overline{g}(\xi ))3KC`$ and so $`\rho \left(d_S(x,\overline{g}g(x))\right)=\rho \left(d_S(x,\overline{g}(\xi ))\right)3KC`$, yielding an upper bound for $`d_S(x,\overline{g}g(x))`$. Similarly, $`d_Y(\xi ,g\overline{g}(\xi ))`$ is bounded for all $`\xi T`$.
Knowing that $`g:ST`$ and $`\overline{g}:TS`$ are coarse lipschitz maps which are coarse inverses of each other, it easily follows that $`g`$ is a quasi-isometry, with quasi-isometry constants depending only on the coarse lipschitz constants for $`g`$ and $`\overline{g}`$, and on the coarse inverse constants for $`g,\overline{g}`$. ∎
## 3 Linear Algebra
In this section we collect some basic results about canonical forms of matrices, and growth of vectors under the action of a matrix.
Let $`(n,F)`$ denote all $`n\times n`$ matrices over a field $`F`$, and let $`\mathrm{GL}(n,F)`$ be the group of invertible matrices. Let $`\mathrm{GL}_0(n,𝐑)`$ be the identity component of $`\mathrm{GL}(n,𝐑)`$, consisting of all matrices of positive determinant.
### 3.1 Jordan Forms
A matrix $`J(k,𝐂)`$ is a *Jordan block* it it has the form $`J=J(k,\lambda )=\lambda \text{Id}+N`$ where $`\lambda 𝐂`$ and $`N_{ij}=\delta (i+1,j)`$, so $`N`$ is the $`k\times k`$ matrix with $`1`$’s on the superdiagonal and $`0`$’s elsewhere.
A matrix $`M(n,𝐂)`$ is in *Jordan form* if it is in block diagonal form
$$M=\left(\begin{array}{cccc}J_1& 0& \mathrm{}& 0\\ 0& J_2& \mathrm{}& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& \mathrm{}& J_I\end{array}\right)$$
where each $`J_i`$ is a Jordan block. Every matrix in $`(n,𝐂)`$ is conjugate, via an invertible complex matrix, to a matrix in Jordan form, unique up to permutation of the Jordan blocks. When all eigenvalues are real, say $`J_i`$ has eigenvalue $`\mathrm{}_i`$, we resolve the nonuniqueness by requiring $`\mathrm{}_1\mathrm{}_2\mathrm{}\mathrm{}_I`$, and for each $`i=1,\mathrm{},I1`$, if $`\mathrm{}_i=\mathrm{}_{i+1}`$ then $`\mathrm{rk}(J_i)\mathrm{rk}(J_{i+1})`$.
A matrix $`J(k,𝐑)`$ is a *real Jordan block* if it has one of the following two forms. The first form is an ordinary Jordan block $`J(k,\mathrm{})`$ where $`\mathrm{}𝐑`$. The second form, which requires $`k`$ to be even, has a $`2\times 2`$ block decomposition of the form
$$J=J(k,a,b)=\left(\begin{array}{ccccc}Q(a,b)& \text{Id}& \mathrm{}& 0& 0\\ 0& Q(a,b)& \mathrm{}& 0& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& \mathrm{}& Q(a,b)& \text{Id}\\ 0& 0& \mathrm{}& 0& Q(a,b)\end{array}\right)$$
where Id is the identity, $`0`$ is the zero matrix, $`Q(a,b)=\left(\begin{array}{cc}a& b\\ b& a\end{array}\right)`$, and $`b0`$.
A matrix $`M(n,𝐑)`$ is in *real Jordan form* if it is in block diagonal form as above where each block $`J_i`$ is a real Jordan block. Every matrix in $`(n,𝐑)`$ is conjugate, via an invertible real matrix, to a matrix in real Jordan form, unique up to permutation of blocks.
The *absolute Jordan form* of $`M(n,𝐑)`$ is the matrix obtained from the Jordan form of $`M`$ by replacing each diagonal entry $`\lambda `$ by $`\mathrm{}=\left|\lambda \right|`$, and permuting the blocks to resolve the nonuniqueness. If $`M`$ is invertible then the absolute Jordan form of $`M`$ can be written in block diagonal form
$$\left(\begin{array}{ccc}J_M^+& 0& 0\\ 0& J_M^0& 0\\ 0& 0& J_M^{}\end{array}\right)$$
where the diagonal entries of $`J_M^+`$ are $`>1`$, of $`J_M^0`$ are $`=1`$, and of $`J_M^{}`$ are $`<1`$. We call $`J_M^+`$ the *expanding part* of the absolute Jordan form, $`J_M^0`$ the *unipotent part*, and $`J_M^{}`$ the *contracting part*, and the block matrix $`\left(\begin{array}{cc}J_M^+& 0\\ 0& J_M^{}\end{array}\right)`$ is called the *nonunipotent part*. Of course, one or more of these parts might be empty.
Note that the Jordan form of the real matrix $`J(k,a,b)`$ is
$$\left(\begin{array}{cc}J(k/2,a+bi)& 0\\ 0& J(k/2,abi)\end{array}\right)$$
and so the absolute Jordan form of $`J(k,a,b)`$ is
$$\left(\begin{array}{cc}J(k/2,\sqrt{a^2+b^2})& 0\\ 0& J(k/2,\sqrt{a^2+b^2})\end{array}\right)$$
Given $`M(n,𝐑)`$, this process may be applied block-by-block to the real Jordan form of $`M`$, and the blocks then permuted, to obtain the absolute Jordan form of $`M`$.
Let $`\mathrm{GL}_\times (n,𝐑)`$ denote the set of all matrices in $`\mathrm{GL}(n,𝐑)`$ lying on a 1-parameter subgroup of $`\mathrm{GL}(n,𝐑)`$, so $`\mathrm{GL}_\times (n,𝐑)\mathrm{GL}_0(n,𝐑)`$. It is well-known and easy to see, given a matrix $`M\mathrm{GL}(n,𝐑)`$, that $`M\mathrm{GL}_\times (n,𝐑)`$ if and only if the negative eigenvalue Jordan blocks of $`M`$ may be paired up so that the two blocks occuring in each pair are identical to each other, and this occurs if and only if $`M`$ has a square root in $`\mathrm{GL}(n,𝐑)`$. Thus, if $`M`$ does not already lie on a 1-parameter subgroup then $`M^2`$ does. We are therefore free to replace a matrix by its square in order to land on a 1-parameter subgroup.
Given a 1-parameter subgroup $`\rho (t)`$ of $`\mathrm{GL}(n,𝐑)`$, if $`M=\rho (1)`$ then we will often abuse notation and write $`\rho (t)=M^t`$, despite the fact that $`M`$ may not lie on a unique 1-parameter subgroup.
Given $`A(n,𝐑)`$ in Jordan form—no $`J(k,a,b)`$ blocks—we say that $`\rho (t)=e^{At}`$ is a 1-parameter *Jordan subgroup*. Notice that the matrices $`e^{At}`$ are *not* themselves in Jordan form. For example when $`A=J(n,\mathrm{})=\mathrm{}\text{Id}+N`$ is a single $`n\times n`$ Jordan block then $`e^{At}`$ is obtained by multiplying the scalar $`e^\mathrm{}t`$ with the matrix
$$e^{Nt}=\underset{i=0}{\overset{n}{}}\frac{1}{i!}N^it^i=\left(\begin{array}{cccccc}1& t& \frac{t^2}{2!}& \mathrm{}& \frac{t^{n1}}{(n1)!}& \frac{t^n}{n!}\\ & 1& t& \mathrm{}& \frac{t^{n2}}{(n2)!}& \frac{t^{n1}}{(n1)!}\\ & & 1& \mathrm{}& \frac{t^{n3}}{(n3)!}& \frac{t^{n2}}{(n2)!}\\ & & & \mathrm{}& \mathrm{}& \mathrm{}\\ & & & & 1& t\\ & & & & & 1\end{array}\right)$$
(3.1)
Nevertheless, for any Jordan form matrix $`J=\mathrm{}\text{Id}+N`$ with $`\mathrm{}𝐑`$, the Jordan form of $`e^J`$ is $`e^{\mathrm{}}\text{Id}+N`$.
Given a general 1-parameter subgroup $`e^{\mu t}`$ in $`\mathrm{GL}(n,𝐑)`$, choose $`A`$ so that $`A^1\mu A`$ is in real Jordan form, and so $`A^1\mu A=\delta +\nu +\eta `$ where $`\delta `$ is diagonal, $`\nu `$ is superdiagonal, and $`\eta `$ is skew-symmetric. We then have
$$e^{\mu t}=(Ae^{(\delta +\nu )t}A^1)(Ae^{\eta t}A^1)$$
Since $`\eta `$ is skew symmetric it follows that $`e^{\eta t}`$ is in the orthogonal group $`\mathrm{O}(n,𝐑)`$. We have therefore proved (see \[Wit90\] for this particular formulation):
###### Proposition 3.1 (1-parameter Real Jordan Form).
Let $`M^t`$ be a 1-parameter subgroup of $`\mathrm{GL}(n,𝐑)`$. There exists a 1-parameter Jordan subgroup $`e^{Jt}`$, a matrix $`A\mathrm{GL}(n,𝐑)`$, and a bounded 1-parameter subgroup $`P^t`$ conjugate into the orthogonal group $`\mathrm{O}(n,𝐑)`$, such that $`e^J`$ is the absolute Jordan form of $`M`$, and letting $`\overline{M}^t=A^1e^{Jt}A`$ we have
$$M^t=\overline{M}^tP^t=P^t\overline{M}^t$$
###### Remark.
In \[Wit90\] the subgroup $`\overline{M}^t`$ is called the *nonelliptic part* of $`M^t`$, and $`P^t`$ is called the *elliptic part*. These two 1-parameter subgroups, which commute with each other, are uniquely determined by $`M^t`$.
### 3.2 Growth of vectors under a linear transformation
Consider a 1-parameter subgroup $`M^t`$ of $`\mathrm{GL}(n,𝐑)`$, with real Jordan form $`M^t=(A^1e^{Jt}A)P^t=\overline{M}^tP^t`$. Let
$$0<\lambda _1<\mathrm{}<\lambda _L$$
be the eigenvalues of $`\overline{M}`$. Let $`V_l=\mathrm{ker}((\lambda _l\text{Id}\overline{M})^m)`$ be the *root space* of the eigenvalue $`\lambda _l`$, where $`m`$ is the multiplicity of $`\lambda _l`$. Let $`n_l`$ be the *index of nilpotency* of $`\overline{M}|V_l`$, the smallest integer such that $`V_l=\mathrm{ker}((\lambda _l\text{Id}\overline{M})^{n_l})`$. For $`i=0,\mathrm{},n_l1`$ let $`V_{l,i}=\mathrm{ker}((\lambda _l\text{Id}\overline{M})^{i+1})`$, so $`V_{l,0}`$ is the eigenspace of $`\lambda _l`$ and $`V_{l,n_l1}=V_l`$. We thus have the *Jordan decomposition* of $`\overline{M}`$, which consists of the direct sum of root spaces
$$𝐑^n=V_1\mathrm{}V_L$$
together with the *Jordan filtrations*
$$V_{l,0}\mathrm{}V_{l,n_l1}=V_l,l=1,\mathrm{},L$$
This decomposition is uniquely determined by $`\overline{M}`$, and hence by $`M`$.
###### Proposition 3.2 (Growth of vectors).
With the above notation, there exist constants $`A,B>0`$ with the following properties. Given $`l=1,\mathrm{},L`$ with $`\lambda _l1`$, we have:
If $`vV_l`$ and $`t0`$ then
$$M^tvA\lambda _l^tv$$
In fact the same lower bound holds if $`vV_lV_{l+1}\mathrm{}V_L`$.
Given $`i=0,\mathrm{},n_l1`$, if $`vV_{l,i}`$ and $`t1`$ then
$$M^tvB\lambda _l^tt^iv$$
In fact the same upper bound holds if $`v(V_1\mathrm{}V_{l1})V_{l,i}`$.
Given $`i=0,\mathrm{},n_l1`$, if $`vV_{l,i}V_{l,i1}`$ then there exists $`C_v>0`$ such that if $`t1`$ then
$$M^tvC_v\lambda _l^tt^i$$
Given $`l=1,\mathrm{},L`$ with $`\lambda _l0`$, similar statements are true with negative values of $`t`$.
###### Proof.
We start with the case when $`M^t=e^{Jt}`$ is a 1-parameter Jordan subgroup, and the proposition follows by examining each Jordan block (3.1).
The second case we consider is when $`M^t`$ has all positive real eigenvalues. By Proposition 3.1 we have $`M^t=A^1e^{Jt}A`$, and Proposition 3.2 follows immediately from the first case applied to $`e^{Jt}`$, together with the fact that $`A`$ takes the Jordan decomposition of $`M^t`$ to the Jordan decomposition of $`e^{Jt}`$.
In the general case, applying Proposition 3.1 we have $`M^t=(A^1e^{Jt}A)P^t=\overline{M}^tP^t`$. We can the apply the second case to $`\overline{M}^t=A^1e^{Jt}A`$. Since $`P^t`$ commutes with $`\overline{M}^t`$ it follows that $`P^t`$ preserves the Jordan decomposition of $`\overline{M}^t`$. Proposition 3.2 then follows from the boundedness of $`P^t`$. ∎
## 4 The Solvable Lie Group $`G_M`$
Recall that $`\mathrm{GL}_\times (n,𝐑)`$ denotes those matrices in $`\mathrm{GL}(n,𝐑)`$ which lie on a 1-parameter subgroup of $`\mathrm{GL}(n,𝐑)`$. Also, each matrix in $`\mathrm{GL}_\times (n,𝐑)`$ has positive determinant.
Given a matrix $`M\mathrm{GL}_\times (n,𝐑)`$ lying on a 1-parameter subgroup $`M^t`$ of $`\mathrm{GL}(n,𝐑)`$, we associate a solvable Lie group denoted $`G_M`$. This is the semidirect product $`G_M=𝐑^n_M𝐑`$ with multiplication defined by
$$(x,t)(y,s)=(x+M^ty,t+s)$$
for all $`(x,t),(y,s)𝐑^n\times 𝐑`$. We will often identify $`G_M=𝐑^n_M𝐑`$ with the underlying set $`𝐑^n\times 𝐑`$.
###### Remark.
Although the Lie group $`G_M`$ depends on more than just the matrix $`M=M^1`$ itself—it depends on the entire 1-parameter subgroup $`M^t`$—we suppress this dependence in our notation $`G_M=𝐑^n_M𝐑`$. This is justified by the fact that the quasi-isometry type of $`G_M`$ depends only on $`M`$, not on the 1-parameter subgroup containing $`M`$ (see the remark after Proposition 4.1). Henceforth, when we say something like “Given $`M\mathrm{GL}_\times (n,𝐑)\mathrm{}`$”, we will either implicitly or explicitly choose a 1-parameter subgroup $`M^t<\mathrm{GL}(n,𝐑)`$ with $`M^1=M`$, which in turn determines $`G_M`$.
If $`M`$ has integer entries then there is a homomorphism $`\mathrm{\Gamma }_MG_M`$ taking the commuting generators $`a_1,\mathrm{},a_n`$ to the standard basis of the integer lattice $`𝐙^n\times 0𝐑^n\times 0𝐑^n\times 𝐑`$, and taking the stable letter $`t`$ to the generator $`(0,1)𝐑^n\times 𝐑`$. The relator $`ta_it^1=\varphi _M(a_i)`$ is checked by noting that
$$(0,1)(x,0)(0,1)=(Mx,0),\text{for all}x𝐑^n$$
Cocompactness of the image of this homomorphism is obvious. To see that $`\mathrm{\Gamma }_M`$ embeds in $`G_M`$ one checks that in the abelian-by-cyclic extension $`1A\mathrm{\Gamma }_M𝐙1`$, the group $`A`$ is identified with the nested union $`𝐙^nM^1(𝐙^n)M^2(𝐙^n)\mathrm{}`$, in $`𝐑^n`$. This also shows that discreteness of $`\mathrm{\Gamma }_M`$ in $`G_M`$ is equivalent to $`detM=1`$, which is equivalent to $`𝐙^n=M(𝐙^n)`$.
For the next several sections we will investigate the geometry of the solvable Lie group $`G_M`$. In this section we begin by showing that $`G_M`$ and $`G_N`$ are quasi-isometric if $`M,N`$ have powers with the same absolute Jordan form. Later in §7 we will see that when $`M`$ has integer entries, much of the geometry of $`\mathrm{\Gamma }_M`$ is reflected in the geometry of $`G_M`$.
We endow $`G_M`$ with the left invariant metric determined by taking the standard Euclidean metric at the identity of $`G_M𝐑^n\times 𝐑=𝐑^{n+1}`$. At a point $`(x,t)𝐑^n\times 𝐑G_M`$, the tangent space is identified with $`𝐑^n\times 𝐑`$, and the Riemannian metric is given by the symmetric matrix
$$\left(\begin{array}{cc}Q_M(t)\hfill & 0\hfill \\ 0\hfill & 1\hfill \end{array}\right)$$
where $`Q_M(t)=(M^t)^TM^t`$. For each $`t𝐑`$, the identification $`𝐑^n𝐑^n\times tG_M`$ induces in $`𝐑^n`$ the metric determined by the quadratic form $`Q_M(t)`$. This metric has distance formula
$$d_{M,t}(x,y)=M^t(xy)$$
##### Remarks
* When $`M`$ is a $`1\times 1`$ matrix with entry $`a>1`$, the group $`G_M`$ is isomorphic to $`\mathrm{Aff}(𝐑)`$, the group of affine transformations of $`𝐑`$, and as a Riemannian manifold $`G_M`$ is isometric to a scaled copy of the hyperbolic plane with constant sectional curvature depending on $`a`$.
* The eigenvalues of $`M`$ are greater than $`1`$ in absolute value if and only if all sectional curvatures of $`G_M`$ are negative (see \[Hei74\]).
###### Proposition 4.1 (How the metric on $`G_M`$ depends on choices).
Given 1-parameter subgroups $`M^t,N^t`$ in $`\mathrm{GL}(n,𝐑)`$, suppose there exist real numbers $`r,s>0`$ such that $`M^r`$ and $`N^s`$ have the same absolute Jordan form. Then the metric spaces $`G_M`$ and $`G_N`$ are quasi-isometric. To be explicit there exists $`A\mathrm{GL}(n,𝐑)`$ and $`K1`$ such that for each $`t𝐑`$, the map $`vA(v)`$ is a $`K`$-bilipschitz homeomorphism from the metric $`d_{M,t}`$ to the metric $`d_{N,\frac{s}{r}t}`$; it follows that the map from $`G_M=𝐑^n_M𝐑`$ to $`G_N=𝐑^n_N𝐑`$ given by
$$(x,t)(Ax,\frac{s}{r}t)$$
is a bilipschitz homeomorphism from $`G_M`$ to $`G_N`$, with bilipschitz constant $`sup\{K,\frac{s}{r},\frac{r}{s}\}`$.
###### Remark.
The absolute Jordan form of $`M^r`$ is uniquely determined by $`M`$ and $`r`$: it is the $`r^{\text{th}}`$ power of the absolute Jordan form of $`M`$. It follows in particular that the quasi-isometry type of $`G_M`$ depends only on the matrix $`M=M^1`$, not on the choice of 1-parameter subgroup $`M^t`$.
###### Proof of Proposition 4.1.
We proceed in cases.
Case 1: Assume that $`N^t=e^{Jt}`$ is the unique 1-parameter Jordan subgroup such that $`N=e^J`$ is conjugate to the absolute Jordan form of $`M`$. Applying Proposition 3.1 we have
$$M^t=(A^1N^tA)P^t$$
where $`A\mathrm{GL}(n,R)`$ and the 1-parameter subgroup $`P^t`$ is bounded.
Choose $`t𝐑`$ and $`v𝐑^n`$. We must show that the two numbers $`M^tv=P^t(A^1N^tA)v`$ and $`N^tAv`$ have ratio bounded away from $`0`$ and $`\mathrm{}`$, with bound independent of $`t,v`$. Setting $`u=N^tAv`$, it suffices to show that $`P^tA^1u`$ and $`u`$ have bounded ratio. But this is clearly true, with a bound of
$$\left(\underset{t}{sup}P^t\right)\mathrm{Max}\{A,\frac{1}{A}\}$$
since the 1-parameter subgoup $`P^t`$ is bounded.
Case 2: Assume that there exists $`a>0`$ such that $`M^t=N^{at}`$ for all $`t`$. Then the metrics $`d_{M,t}`$ and $`d_{N,at}`$ are identical.
General case: Applying Case 2 we may assume that $`detM=detN`$. Applying Case 1 twice we may go from $`G_M`$ to $`G_{e^J}`$ to $`G_N`$, where $`e^J`$ is conjugate to the absolute Jordan form of $`M`$ and of $`N`$. ∎
## 5 Dynamics of $`G_M`$, Part I: <br>Horizontal Respecting Quasi-isometries
In this section we begin studying the asymptotic geometry of the solvable Lie groups $`G_M`$ associated to 1-parameter subgroups $`M^t`$ of $`\mathrm{GL}(n,𝐑)`$. As we saw in Section 4, the quasi-isometry type of $`G_M`$ depends only on $`M`$, not on the choice of 1–parameter subgroup $`M^t`$ passing through $`M`$; see the remark after Proposition 4.1. We therefore continue to suppress the choice of 1-parameter subgroup in our notation. Further, we do not restrict the determinant to be $`>1`$: the results of this section hold even when $`detM=1`$.
### 5.1 Theorem 5.2 on horizontal respecting quasi-isometries
Let $`X,Y`$ be metric spaces. Let $``$ be a decomposition of $`X`$, that is, a collection of disjoint subsets of $`X`$ whose union is $`X`$. Let $`𝒢`$ be a decomposition of $`Y`$. Motivated by a foliation of a manifold, the elements of these decompositions are called *leaves* and the decomposition itself is called the *leaf space*. A quasi-isometry $`\varphi :XY`$ is said to *coarsely respect* the decompositions $`,𝒢`$ if there exists a number $`A0`$ and a map of leaf spaces $`h:𝒢`$ such that for each leaf $`L`$ we have
$$d_{}(\varphi (L),h(L))A$$
For example, consider the space $`G_M`$. The coordinate function $`G_M𝐑^n\times 𝐑𝐑`$ given by $`(x,t)t`$ is called the *time function* of $`G_M`$. The level sets $`P_t𝐑^n\times t`$ form the *horizontal foliation* of $`G_M`$, whose leaves are called *horizontal leaves* of $`G_M`$, and whose leaf space is $`𝐑`$. Notice that $`d_{}(P_s,P_t)=\left|st\right|`$, and so the time function induces an isometry between the horizontal leaf space equipped with the Hausdorff metric and $`𝐑`$.
Consider another matrix $`N\mathrm{GL}_\times (n,𝐑)`$, and denote the horizontal leaves of $`G_N`$ by $`P_t^{}`$.
###### Definition (Horizontal respecting).
A quasi-isometry $`\varphi :G_MG_N`$ is said to be *horizontal respecting* if it coarsely respects the horizontal foliations of $`G_M,G_N`$. That is, there exists a function $`h:𝐑𝐑`$ and $`A0`$ such that $`d_{}(\varphi (P_t^{}),P_{h(t)}^{})A`$ for all $`t𝐑`$.
The function $`h:𝐑𝐑`$ is called an *induced time change* for $`\varphi `$, with *Hausdorff constant* $`A`$.
If $`h,h^{}`$ are two induced time changes for $`\varphi `$ then $`sup_t\left|h(t)h^{}(t)\right|A+A^{}<\mathrm{}`$, where $`A,A^{}`$ are Hausdorff constants for $`h,h^{}`$ respectively. Also, if $`h:𝐑𝐑`$ is an induced time change for $`\varphi `$ with Hausdorff constant $`A`$, if $`A^{}0`$, and if $`h^{}:𝐑𝐑`$ is any function satisfying $`sup_{t𝐑}\left|h(t)h^{}(t)\right|A^{}`$, then $`h^{}`$ is also an induced time change for $`\varphi `$, with Hausdorff constant $`A+A^{}`$.
###### Lemma 5.1.
For each $`K,C,A`$ there exists $`C^{}`$ such that if $`\varphi :G_MG_N`$ is a horizontal respecting $`(K,C)`$ quasi-isometry, and $`h:𝐑𝐑`$ is an induced time change for $`\varphi `$ with Hausdorff constant $`A`$, then $`h`$ is a $`(K,C^{})`$ quasi-isometry of $`𝐑`$.
###### Proof.
We have $`\left|h(t)h(s)\right|d_{}(P_{h(t)},P_{h(s)})+2AK\left|ts\right|+C+2A`$. The reverse inequality is similar, and so $`h`$ is a quasi-isometric embedding. Since $`\varphi `$ is coarsely onto, an easy argument shows $`h`$ is coarsely onto. ∎
A $`(K,C^{})`$ quasi-isometry $`h:𝐑𝐑`$ induces a bijection of the two-point set $`\mathrm{Ends}(𝐑)=\{\mathrm{},+\mathrm{}\}`$: given $`\eta _1,\eta _2\mathrm{Ends}(𝐑)`$, we have $`h(\eta _1)=\eta _2`$ if and only if $`h`$ takes every sequence that diverges to $`\eta _1`$ to a sequence that diverges to $`\eta _2`$. The following two properties of $`h`$ are equivalent:
1. $`h`$ induces the identity on $`\mathrm{Ends}(𝐑)`$.
2. $`h`$ is *coarsely increasing*, that is there exists $`L>0`$ such that if $`t>s+L`$ then $`h(t)>h(s)`$.
That (2) implies (1) is obvious. The other direction is true with any $`L>2C^{}K`$, for if there existed $`ts+L`$ with $`h(t)h(s)`$, then since $`h`$ induces the identity on $`\mathrm{Ends}(𝐑)`$ there would exist $`t^{}>t`$ such that $`\left|h(s)h(t^{})\right|C^{}`$, but also $`\left|h(s)h(t^{})\right|\left|st^{}\right|/KC^{}L/KC^{}>C^{}`$, a contradiction.
If $`h:𝐑𝐑`$ is an induced time change of a horizontal respecting quasi-isometry $`\varphi :G_MG_N`$, and if $`h`$ satisfies the equivalent properties (1) and (2), then we say that $`\varphi `$ *coarsely respects the transverse orientation* of the horizontal foliations.
Terminology (time vs. height): In some contexts the vertical parameter which we have been calling “time” will also be called height, as sometimes this terminology is more suggestive, for example in discussing horizontal foliations.
Here is the main result, whose proof will occupy the remainder of this section and the next section.
###### Theorem 5.2 (Horizontal respecting quasi-isometries).
Let $`\varphi :G_MG_N`$ be a quasi-isometry which coarsely respects the transversely oriented horizontal foliations of $`G_M`$ and $`G_N`$. Then there exist real numbers $`r,s>0`$ so that $`M^r`$ and $`N^s`$ have the same absolute Jordan form.
Our proof of Theorem 5.2 proceeds in steps, following the outline given in the introduction.
### 5.2 Step 1a: Hyperbolic dynamics and the Shadowing Lemma
The Lie group $`G_M`$ has a natural flow which fits into the theory of partially hyperbolic dynamical systems. From the dynamics we find that the flow has several invariant foliations, the “weak stable, weak unstable, and center” foliations. In §§5.2,5.3, by using the Shadowing Lemma (\[HPS77\], Lemma 7.A.2, page 133), we prove that a horizontal respecting quasi-isometry $`G_MG_N`$ also respects the dynamically defined foliations of $`G_M`$, $`G_N`$.
From this result we obtain the first piece of our rigidity theorem, by showing that expanding, contracting, and unipotent parts of the absolute Jordan forms of $`M`$ and $`N`$ have the same ranks respectively, and that the unipotent parts are identical.
#### 5.2.1 Dynamically defined foliations
Consider a 1-parameter subgroup $`M^t\mathrm{GL}(n,𝐑)`$, with real Jordan form $`M^t=\overline{M}^tP^t`$. Consider the Jordan decomposition of $`\overline{M}`$, and group the root spaces according to whether the corresponding eigenvalue is $`<1`$, $`=1`$, or $`>1`$ (alternatively, a logarithm which is $`<0`$, $`=0`$, or $`>0`$), to obtain a decomposition $`𝐑^n=V^{}V^0V^+`$.
###### Remark.
It might happen that one or two of the factors $`V^{}`$, $`V^0`$, $`V^+`$ is trivial, that is, 0-dimensional, for instance when all eigenvalues of $`M`$ lie outside the unit circle.
Now consider the Lie group $`G_M=𝐑^n_M𝐑`$ determined by a 1-parameter subgroup $`M^t`$. Define the *vertical flow* $`\mathrm{\Phi }`$ on $`G_M`$ to be
$$\mathrm{\Phi }_t(x,s)=(x,s+t)$$
The tangent bundle $`TG_M`$ has a $`\mathrm{\Phi }`$-invariant splitting
$$TG_M=E^sE^cE^u$$
defined as follows. The tangent space at each point $`xG_M`$ is identified with $`𝐑^n𝐑`$, and we take
$$E_x^s=V^{}0,E_x^c=V^0𝐑,E_x^u=V^+0$$
It is evident from the construction that each of the distributions $`E^sE^c,E^uE^c,E^c`$ is integrable, tangent to foliations denoted $`W^s,W^u,W^c`$. We call these foliations the (weak) *stable, unstable, and center foliations*, respectively. The stable and unstable foliations are transverse, and the intersection of any stable leaf with any unstable leaf is a center leaf.
Applying the Exponential Lower Bound from Proposition 3.2, there exist constants $`A>0`$, $`\lambda >1`$ such that:
* If $`vE^u`$, then for $`t0`$ we have $`D\mathrm{\Phi }_tvA\lambda ^tv`$, and for $`t0`$ we have $`D\mathrm{\Phi }_tv\frac{1}{A}\lambda ^tv`$.
* If $`vE^s`$, then for $`t0`$ we have $`D\mathrm{\Phi }_tvA\lambda ^tv`$, and for $`t0`$ we have $`D\mathrm{\Phi }_tv\frac{1}{A}\lambda ^tv`$.
Also, applying the Exponential$``$Polynomial Upper Bound from Proposition 3.2, there exists $`B>0`$ and an integer $`n1`$ such that:
* If $`vE^c`$, then for $`\left|t\right|1`$ we have $`D\mathrm{\Phi }_tvB\left|t\right|^nv`$.
When we want to emphasize the dependence of the $`V`$’s and $`E`$’s on the 1-parameter subgroup $`M^t`$, we will append a subscript, e.g. $`V_M^+`$, $`E_M^s`$, etc.
#### 5.2.2 Shadowing Lemma
Consider a flow $`\mathrm{\Phi }`$ on a metric space $`X`$. We write $`xt`$ as an abbreviation for $`\mathrm{\Phi }_t(x)`$. Given $`ϵ,T>0`$, an *$`(ϵ,T)`$-pseudo-orbit* of $`\mathrm{\Phi }`$ consists of a sequence of flow segments $`(x_i[0,t_i])`$, where the index $`i`$ runs over an interval in $`𝐙`$, such that $`d_X(x_it_i,x_{i+1})<ϵ`$ and $`t_i>T`$ for all $`i`$.
###### Lemma 5.3 (Shadowing Lemma).
Consider a 1-parameter subgroup $`M^t`$ of $`\mathrm{GL}(n,𝐑)`$, and let $`\mathrm{\Phi }`$ be the vertical flow on $`G_M`$. For every $`ϵ,T>0`$ there exists $`\delta ,ϵ^{},T^{}>0`$ such that every $`(ϵ,T)`$-pseudo-orbit of $`\mathrm{\Phi }`$ is $`\delta `$-shadowed by an $`(ϵ^{},T^{})`$-pseudo-orbit of $`\mathrm{\Phi }`$ which is contained in some center leaf. That is, if $`(x_i[0,t_i])`$ is an $`(ϵ,T)`$-pseudo-orbit, then there is an $`(ϵ^{},T^{})`$-pseudo-orbit $`(y_i[0,t_i])`$ contained in some center leaf so that $`d(x_it,y_it)<\delta `$ for all $`i`$ and all $`t[0,t_i]`$.
###### Proof.
By construction, the foliations $`W^s`$ and $`W^u`$ are coordinate foliations in $`𝐑^{n+1}`$; this shows that the flow $`\mathrm{\Phi }`$ has a “global product structure” in the language of hyperbolic dynamical systems. The lemma now follows the proof of the Shadowing Lemma in \[HPS77\], Lemma 7.A.2, page 133. A direct proof is also easy to work out, and is left to the reader. ∎
### 5.3 Step 1b: Foliations rigidity
The Shadowing Lemma implies further rigidity properties of horizontal respecting quasi-isometries:
###### Proposition 5.4 (Foliations rigidity).
Suppose $`\varphi :G_MG_N`$ is a quasi-isometry which coarsely respects the horizontal foliations and their transverse orientations. Then $`\varphi `$ also coarsely respects the weak unstable foliations $`W_M^u`$, $`W_N^u`$, the weak stable foliations $`W_M^s,W_N^s`$, and the center foliations $`W_M^c`$, $`W_N^c`$. In particular
* $`dim(V_M^+)=dim(V_N^+)`$
* $`dim(V_M^{})=dim(V_N^{})`$
* $`dim(V_M^0)=dim(V_N^0)`$
##### Remarks
* In the case where neither $`M`$ nor $`N`$ has any eigenvalue on the unit circle, the center foliations of both $`G_M`$ and $`G_N`$ are simply the foliations by vertical flow lines, and Proposition 5.4 says that $`\varphi `$ respects these foliations. But in the general case, it is not true that $`\varphi `$ always respects the foliations by vertical flow lines. For a simple counterexample, consider the $`1\times 1`$ matrix $`M=N=(1)`$, which gives $`\mathrm{\Gamma }_M=\mathrm{\Gamma }_N=𝐙^2`$. There exist horizontal respecting quasi-isometries of $`𝐑^2=𝐑\times 𝐑`$ which do not respect the vertical foliation.
* If all eigenvalues of $`M`$ and $`N`$ are outside the unit circle, then both $`G_M`$ and $`G_N`$ are negatively curved, and the proposition follows from a standard fact: a quasigeodesic in a negatively curved space $`X`$ is Hausdorff close to a geodesic (this was the approach taken in \[FM98\] in the case of a $`1\times 1`$ matrix $`M`$, where $`G_M`$ is isometric to a scaled copy of the hyperbolic plane). This “fact” is unavailable when $`X=G_M`$ is not negatively curved, forcing us to study horizontal respecting quasi-isometries via the Shadowing Lemma.
Before proving Proposition 5.4, we use it to obtain some pieces of our classification theorem. Since $`\mathrm{rk}(J_M^{})=dim(V_M^{})`$ etc., we immediately have:
###### Corollary 5.5.
If there is a quasi-isometry from $`G_M`$ to $`G_N`$ which coarsely respects the transversely oriented horizontal foliations, then $`\mathrm{rk}(J_M^{})=\mathrm{rk}(J_N^{})`$, $`\mathrm{rk}(J_M^0)=\mathrm{rk}(J_N^0)`$, and $`\mathrm{rk}(J_M^+)=\mathrm{rk}(J_N^+)`$.∎
We also have:
###### Corollary 5.6.
The unipotent blocks of the absolute Jordan forms of $`M`$ and $`N`$ are identical.
###### Proof.
Let $`L`$ be some center leaf of $`G_M`$, of dimension $`k`$. From Proposition 5.4 it follows that $`\varphi (L)`$ is Hausdorff close to some center leaf $`L^{}`$ of $`G_N`$, also of dimension $`k`$. By composition with nearest point projection (which moves points a uniformly bounded amount) we get an induced map $`LL^{}`$. By Lemma 2.1 this map is a quasi-isometry. By Proposition 4.1, $`L`$ and $`L^{}`$ are quasi-isometric to the nilpotent Lie groups $`𝐑^{k1}_{J_M^0}𝐑`$ and $`𝐑^{k1}_{J_N^0}𝐑`$, respectively. As Bridson and Gersten have shown \[BG96\], Pansu’s invariant \[Pan89b\] may be used to prove that $`J_M^0=J_N^0`$. ∎
###### Proof of Proposition 5.4.
We begin with:
###### Claim 5.7.
For each vertical flow line $`\gamma =\mathrm{\Phi }_𝐑(x)`$ in $`G_M`$, there exists a center leaf $`\tau _\gamma `$ in $`G_N`$ such that $`\varphi (\gamma )`$ is contained in the $`\alpha `$-neighborhood of $`\tau _\gamma `$, where the constant $`\alpha >0`$ does not depend on $`\gamma `$.
Before proving the claim, we apply it to prove the proposition as follows.
Consider any two vertical flow lines $`\gamma _1,\gamma _2`$ in $`G_M`$. By the claim we have that $`\varphi (\gamma _1)`$ and $`\varphi (\gamma _2)`$ lie, respectively, in bounded neighborhoods of center leaves $`\sigma _1`$ and $`\sigma _2`$ of $`G_N`$. Since $`h(t)\pm \mathrm{}`$ as $`t\pm \mathrm{}`$, for each choice of sign $`+`$ or $``$ the following two statements are equivalent, and the second statement implies the third:
1. The distance between the points $`\gamma _1P_t`$ and $`\gamma _2P_t`$ in $`P_t`$ stays bounded as $`t\pm \mathrm{}`$.
2. The distance between the points $`\varphi (\gamma _1)P_{h(t)}`$ and $`\varphi (\gamma _2)P_{h(t)}`$ in $`P_{h(t)}`$ stays bounded as $`t\pm \mathrm{}`$.
3. The Hausdorff distance between the sets $`\sigma _1P_{h(t)}`$ and $`\sigma _2P_{h(t)}`$ in $`P_{h(t)}`$ stays bounded as $`t\pm \mathrm{}`$.
Using $``$ signs, the first statement is equivalent to saying that $`\gamma _1,\gamma _2`$ are contained in the same unstable leaf of $`G_M`$, and the third statement is equivalent to saying that $`\sigma _1,\sigma _2`$ are contained in the same unstable leaf of $`G_N`$. It follows that $`\varphi `$ takes every unstable leaf of $`G_M`$ into a bounded neighborhood of an unstable leaf of $`G_N`$. Applying the same argument to a coarse inverse $`\overline{\varphi }`$ of $`\varphi `$ gives the opposite inclusion. Since $`d(\overline{\varphi }\varphi ,\text{Id})<\mathrm{}`$ it follows that the image under $`\varphi `$ of any unstable leaf of $`G_M`$ lies a bounded Hausdorff distance from an unstable leaf of $`G_N`$, that is, $`\varphi `$ coarsely preserves the unstable foliations. A similar argument using $`+`$ signs shows that $`\varphi `$ coarsely preserves stable foliations. By taking intersections of stable and unstable leaves it follows that $`\varphi `$ coarsely preserves center foliations.
The final statements about dimensions follow from the fact that dimension is a quasi-isometry invariant, for leaves of the foliations in question; see \[Ger93\] or \[BW97\].
It remains to prove the claim. Applying Lemma 5.1, we have an induced time change $`h:𝐑𝐑`$ which is a $`(K,C^{})`$-quasi-isometry with Hausdorff constant $`A`$, where $`C^{}`$ depends only on $`K,C,A`$. Furthermore by Lemma 5.1 and the comments following it, the map $`h`$ is coarsely increasing: there exists $`L=L(K,C,A)>0`$ such that if $`ts+L`$ then $`h(t)>h(s)`$.
We can furthermore increase $`L`$, depending only on $`K,C^{},A`$, so that:
$$\begin{array}{c}\text{if }t^{}t+L,xP_t^{},yP_t,\varphi (x)P_s^{}\text{ and }\varphi (y)P_s,\hfill \\ \hfill \text{then }s^{}s+1.\end{array}$$
(5.1)
In fact taking $`L>(C^{}+2A+1)K`$ will do, for then we have
$`h(t^{})`$ $`h(t)+(t^{}t)/KC^{}h(t^{})+L/KC^{}`$
$`h(t)+2A+1`$
and, since $`P_s^{}`$ is $`A`$-Hausdorff close to $`P_{h(t^{})}`$ and $`P_s`$ is $`A`$-Hausdorff close to $`P_{h(t)}`$, it follows that $`s^{}s+1`$.
To prove the claim, we first show that $`\varphi (\gamma )`$ is Hausdorff close to some pseudo-orbit in $`G_N`$, and then we apply the Shadowing Lemma to show that the pseudo-orbit lies in a bounded neighborhood of some center leaf.
To be more precise, fix a point $`x_0\gamma `$ and consider the sequence $`x_i=\mathrm{\Phi }_{iL}(x_0)`$ for $`i𝐙`$. Let $`y_i=\varphi (x_i)`$, and let $`s_i`$ be such that $`y_iP_{s_i}`$. From (5.1) it follows that $`s_{i+1}s_i+1`$. Let $`t_i=s_{i+1}s_i1`$.
We claim that there exists $`ϵ>0`$, depending ultimately only on $`K,C`$, so that $`\left(y_i[0,t_i]\right)`$ is an $`(ϵ,1)`$-pseudo-orbit; in other words, $`d(y_it_i,y_{i+1})`$ is bounded. To see why, first note that
$`d(y_it_i,y_i)`$ $`=t_i=s_{i+1}s_i`$
$`2A+h(L(i+1))h(Li)`$
$`2A+KL+C^{}`$
and then
$$d(y_i,y_{i+1})Kd(x_i,x_{i+1})+CKL+C$$
so we may take $`ϵ=2A+2KL+C+C^{}`$.
Applying the Shadowing Lemma, there exists $`\beta ,ϵ^{},T^{}`$ such that $`\left(y_i[0,t_i]\right)`$ is $`\beta `$-Hausdorff close to an $`(ϵ^{},T^{})`$-pseudo-orbit $`\left(y_i^{}[0,t_i]\right)`$ contained in some center leaf of $`G_N`$. On the other hand, since every point of $`\gamma `$ is within distance $`L`$ of some $`x_i`$, it follows that $`\varphi (\gamma )`$ is uniformly Hausdorff close to $`\left(y_i[0,t_i]\right)`$, and so it is also uniformly close to the pseudo-orbit $`\left(y_i^{}[0,t_i]\right)`$. ∎
### 5.4 Step 2: Time rigidity
The main result of this subsection says that a horizontal respecting quasi-isometry has an induced time change function which is be affine.
###### Proposition 5.8 (Time rigidity).
Consider the Lie groups $`G_M,G_N`$ where $`M,N\mathrm{GL}_\times (n,𝐑)`$ each have an eigenvalue of absolute value greater than 1. Then there exists $`m𝐑_+`$ with the following properties. For all $`K1`$, $`C,A0`$ there exists $`A^{}0`$ such that if $`\varphi :G_MG_N`$ is a $`(K,C)`$ quasi-isometry which coarsely respects horizontal foliations and their transverse orientations, with an induced time change of Hausdorff constant $`A`$, then there exists $`b𝐑`$ such that $`h(t)=mt+b`$ is an induced time change with Hausdorff constant $`A^{}`$. In fact, $`m`$ can be computed as follows: Let $`\alpha `$ (resp. $`\beta `$) be the least eigenvalue greater than $`1`$ of the absolute Jordan form of $`M`$ (resp. $`N`$); the numbers $`\alpha ,\beta `$ exist by the assumption on eigenvalues. Then $`m=\mathrm{log}\alpha /\mathrm{log}\beta `$.
##### Remarks
* In the case of self-quasi-isometries of $`\mathrm{Aff}(R)=G_{(e^1)}=𝐇^2`$ which coarsely respect the horizontal foliation, this result is part of Proposition 5.3 of \[FM98\], where the conclusion is that the induced time change is a translation of $`𝐑`$.
* One of the delicate points in Gromov’s development of the inf$`\delta `$im invariant is the rescaling problem discussed at the beginning of Section 7.$`\mathrm{C}_1`$ of \[Gro93\]: rates of exponential growth change when the parameter is rescaled. Time Rigidity allows us to avoid the rescaling problem altogether, by showing that the time parameter is “natural” with respect to quasi-isometries.
###### Proof.
This proof will define a sequence of constants which will depend on $`K,C,A`$ and on the matrices $`M`$ and $`N`$. We will indicate the dependence on $`K,C,A`$ by writing, for example, $`C_1=C_1(K,C,A)`$, but we will suppress the dependence on $`M,N`$. Although each constant in the sequence will depend on previous constants in the sequence, by induction it will ultimately depend only on $`K,C,A,M,N`$.
###### Claim 5.9.
For each fixed time $`t_0`$, and for each $`tt_0`$, we have
$$h(t)m(tt_0)+h(t_0)C_1$$
for some $`C_1=C_1(K,C,A)0`$.
Accepting this claim for the moment, we prove the proposition. The idea is simply that the conclusion of the claim, applied to both $`h`$ and its coarse inverse $`\overline{h}`$, with $`t_0+\mathrm{}`$, implies the proposition.
Let $`s`$ be a time parameter for $`G_N`$. Let $`\overline{\varphi }:G_NG_M`$ be a coarse inverse for $`\varphi `$, also a quasi-isometry which coarsely respects the horizontal foliations and their transverse orientations, and with an induced time change $`\overline{h}(s)`$. The constants for $`\overline{\varphi }`$ and $`\overline{h}`$ depend only on $`K,C,A`$. The claim therefore applies as well to $`\overline{h}`$ and we obtain, for each fixed time $`s_0`$ and each $`ss_0`$,
$$\overline{h}(s)\frac{1}{m}(ss_0)+\overline{h}(s_0)C_2$$
for some $`C_2=C_2(K,C,A)0`$.
It is clear $`\overline{h}`$ is a coarse inverse for $`h`$, that is:
$$\left|\overline{h}(h(t))t\right|C_3,\left|h(\overline{h}(s))s\right|C_3$$
for some $`C_3=C_3(K,C,A)0`$.
Also, by Lemma 5.1 and the comments after it, the map $`h`$ is coarsely increasing: there exists $`L=L(K,C,A)0`$ such that if $`t^{}>t+L`$ then $`h(t^{})>h(t)`$.
We reverse the inequality in the claim as follows. Fix $`t_0`$. Let $`s_0=t_0`$. Consider for the moment some $`tt_0L`$. Letting $`s=h(t)`$ it follows that $`ss_0`$ and so we have
$$\overline{h}(h(t))\frac{1}{m}(h(t)h(t_0))+\overline{h}(h(t_0))C_2$$
But $`t+C_3\overline{h}(h(t))`$ and $`\overline{h}(h(t_0))t_0C_3`$ and so we obtain
$`t`$ $`{\displaystyle \frac{1}{m}}(h(t)h(t_0))+t_0(2C_3+C_2)`$
$`h(t)`$ $`m(tt_0)+h(t_0)+m(2C_3+C_2)`$
This has been derived only for $`tt_0L`$, but for $`t_0Ltt_0`$ we obtain a similar inequality with another constant in place of $`m(2C_3+C_2)`$. Therefore, for all $`tt_0`$ we obtain
$$m(tt_0)+h(t_0)C_4h(t)m(tt_0)+h(t_0)+C_4$$
for some $`C_4=C_4(K,C,A)`$. Note that this is true *for all* $`t_0`$, with $`C_4`$ independent of $`t_0`$.
In particular, taking $`t_0=0`$, for all $`t0`$ we obtain
$$mt+h(0)C_4h(t)mt+h(0)+C_4$$
Now take any $`t_10`$, and since $`0t_1`$ we obtain
$$m(0t_1)+h(t_1)C_4h(0)m(0t_1)+h(t_1)+C_4$$
and so
$$mt_1+h(0)C_4h(t_1)mt_1+h(0)+C_4$$
Taking $`b=h(0)`$, this proves that $`mt+b`$ is an induced time change for $`\varphi `$, with Hausdorff constant $`A^{}=C_4+A`$.
Now we turn to the proof of Claim 5.9.
Let $`M^t=\overline{M}^tQ^t`$, $`N^t=\overline{N}^tQ^{}^t`$ be the real Jordan forms. Let $`U`$ (resp. $`U^{}`$) be the root space with eigenvalue $`1`$ for $`\overline{M}`$ (resp. $`\overline{N}`$). Let $`W`$ (resp. $`W^{}`$) be the direct sum of root spaces with eigenvalue $`1`$ for $`\overline{M}`$ (resp. $`\overline{N}`$). Recall that $`\alpha `$ is the smallest eigenvalue $`>1`$ for $`\overline{M}`$, and $`\beta `$ is the smallest eigenvalue $`>1`$ for $`\overline{N}`$. Let $`V`$ be the direct sum of $`U`$ and the eigenspace with eigenvalue $`\alpha `$ for $`\overline{M}`$. We have $`UVW`$; let $`(U),(V),(W)`$ be the corresponding foliations of $`G_M𝐑^n\times 𝐑`$, whose leaves are parallel to $`U\times 𝐑,V\times 𝐑,W\times R`$ respectively. We also have $`U^{}W^{}`$; let $`(U^{}),(W^{})`$ be the corresponding foliations of $`G_N`$.
Here is the idea for proving Claim 5.9. Each leaf of $`(V)`$ is foliated by leaves of $`(U)`$. Because $`V`$ is the direct sum of $`U`$ with the $`\alpha `$ eigenspace of $`\overline{M}`$, it follows that as $`t\mathrm{}`$ distinct leaves of $`(U)`$ in $`(V)`$ diverge from each other *exactly* as $`\alpha ^t`$, measured in the time $`t`$ horizontal plane of $`G_M`$. This is a consequence of the Exponential Lower Bound and the Exponential$``$Polynomial Upper Bound in Proposition 3.2; notice that it is critical here that $`V`$ not be the direct sum of $`U`$ with the $`\alpha `$ *root space*, for then Exponential$``$Polynomial Upper Bound would be at best $`\alpha ^t`$ times some polynomial, which would mess up the following calculations. Mapping over via the quasi-isometry $`\varphi :G_MG_N`$, distinct leaves of $`(U)`$ in a single leaf of $`(V)`$ must (coarsely) map to distinct leaves of $`(U^{})`$ in a single leaf of $`(W^{})`$, which as $`s\mathrm{}`$ diverge from each other *at least as fast as* $`\beta ^s`$, by the Exponential Lower Bound. The time change map $`th(t)=s`$ therefore cannot grow slower than $`s=\frac{\mathrm{log}\alpha }{\mathrm{log}\beta }t`$, as $`t\mathrm{}`$.
To make this precise, pick a leaf $`L_V`$ of $`(V)`$, contained in some leaf $`L_W`$ of $`(W)`$. We use the symbol $`\gamma `$ to denote a general leaf of $`(U)`$, which we will typically take to be a subset of $`L_V`$. By Proposition 5.4, there exists a leaf $`L_W^{}`$ of $`(W^{})`$ such that
$$d_{}(f(L_W),L_W^{})C_5=C_5(K,C,A)$$
and for each leaf $`\gamma `$ of $`(U)`$ there exists a leaf $`\gamma ^{}`$ of $`(U^{})`$ such that
$$d_{}(f(\gamma ),\gamma ^{})C_5$$
Moreover, if $`\gamma L_V`$ then $`\gamma ^{}L_W^{}`$, because $`L_VL_W`$ and so $`\gamma ^{}`$ stays in a bounded neighborhood of $`L_W^{}`$, but any leaf of $`(U^{})`$ which is not a subset of $`L_W^{}`$ has points which are arbitrarily far from $`L_W^{}`$.
Let $`P_t`$ be the horizontal subset of $`G_M`$ at height $`t𝐑`$, and let $`d_t`$ denote Hausdorff distance in $`P_t`$ between closed subsets of $`P_t`$. Let $`P_s^{}`$ be the horizontal subset of $`G_N`$ at height $`s𝐑`$, and let $`d_s^{}`$ denote Hausdorff distance in $`P_s^{}`$.
Since the Hausdorff distance in $`G_N`$ between $`\varphi (P_t)`$ and $`P_{h(t)}^{}`$ is at most $`A`$, the vertical projection from $`\varphi (P_t)`$ to $`P_{h(t)}^{}`$ induces a quasi-isometry between $`P_t`$ and $`P_{h(t)}^{}`$; the multiplicative constant of this quasi-isometry is $`K`$, and its additive constant depends only on $`K,C,A`$. It follows that there exists a “coarseness constant” $`C_6=C_6(K,C,A)`$ so that for any $`t`$, and for any $`x,yP_t`$ with $`d_t(x,y)C_6`$, if $`x^{},y^{}P_{h(t)}^{}`$ are the vertical projections of $`\varphi (x),\varphi (y)`$ then
$$\frac{1}{2K}d_t(x,y)d_{h(t)}^{}(x^{},y^{})2Kd_t(x,y)$$
(5.2)
To prove Claim 5.9, fix a time $`t_0`$ and let $`s_0=h(t_0)`$. Let $`\gamma _1,\gamma _2`$ be two leaves of $`(U)`$ contained in $`L_V`$, and let $`\gamma _i^{}`$ be the unique leaf of $`(U^{})`$ within bounded Hausdorff distance of $`\varphi (\gamma _i)`$; this bound depends only on $`K,C,A`$, as shown in Proposition 5.4.
In $`G_M`$, apply the Exponential Lower Bound and the Exponential$``$Polynomial Upper Bound of Proposition 3.2, and so for all $`tt_0`$ we have
$`A\alpha ^{t+t_0}d_{t_0}(\gamma _1P_{t_0},\gamma _2P_{t_0})`$ $`d_t(\gamma _1P_t,\gamma _2P_t)`$
$`B\alpha ^{t+t_0}d_{t_0}(\gamma _1P_{t_0},\gamma _2P_{t_0})`$
where $`A,B`$ depend only on $`G_M`$ (note that $`t=t_0`$ gives $`A1B`$).
We want the distance between $`\gamma _1`$ and $`\gamma _2`$ in $`P_t`$ to be greater than the coarseness constant $`C_6`$, for each $`tt_0`$, in order that property (5.2) may be applied. We therefore impose a condition on $`\gamma _1`$ and $`\gamma _2`$, namely that
$$d_{t_0}(\gamma _1P_{t_0},\gamma _2P_{t_0})\frac{C_6}{A}$$
which implies, for all $`tt_0`$, that
$$d_t(\gamma _1P_t,\gamma _2P_t)C_6$$
and so
$`{\displaystyle \frac{1}{2K}}d_t(\gamma _1P_t,\gamma _2P_t)`$ $`d_{h(t)}^{}(\gamma _1^{}P_{h(t)}^{},\gamma _2^{}P_{h(t)}^{})`$
$`2Kd_t(\gamma _1P_t,\gamma _2P_t)`$
which implies
$`{\displaystyle \frac{A}{2K}}\alpha ^{t+t_0}d_{t_0}(\gamma _1P_{t_0},\gamma _2P_{t_0})`$ $`d_{h(t)}^{}(\gamma _1^{}P_{h(t)}^{},\gamma _2^{}P_{h(t)}^{})`$
$`2BK\alpha ^{t+t_0}d_{t_0}(\gamma _1P_{t_0},\gamma _2P_{t_0})`$
Next, applying the Exponential Lower Bound of Proposition 3.2 in $`G_N`$, for each $`ss_0`$ we have
$$d_s^{}(\gamma _1^{}P_s^{},\gamma _2^{}P_s^{})A\beta ^{s+s_0}d_{s_0}^{}(\gamma _1^{}P_{s_0}^{},\gamma _2^{}P_{s_0}^{})$$
Taking $`s=h(t)`$, and using the fact that $`s_0=h(t_0)`$, this implies
$$\beta ^{h(t)+h(t_0)}d_{h(t_0)}^{}(\gamma _1^{}P_{h(t_0)}^{},\gamma _2^{}P_{h(t_0)}^{})\frac{2BK}{A}\alpha ^{t+t_0}d_{t_0}(\gamma _1P_{t_0},\gamma _2P_{t_0})$$
Therefore,
$$\beta ^{h(t)+h(t_0)}d_{t_0}(\gamma _1P_{t_0},\gamma _2P_{t_0})\frac{4BK^2}{A}\alpha ^{t+t_0}d_{t_0}(\gamma _1P_{t_0},\gamma _2P_{t_0})$$
Now divide both sides by $`d_{t_0}(\gamma _1P_{t_0},\gamma _2P_{t_0})`$, and take logarithms, obtaining
$$(h(t)+h(t_0))\mathrm{log}(\beta )\mathrm{log}\left(\frac{4BK^2}{A}\right)+(t+t_0)\mathrm{log}(\alpha )$$
and so
$$h(t)\frac{\mathrm{log}(\alpha )}{\mathrm{log}(\beta )}(tt_0)+h(t_0)\frac{\mathrm{log}\left(\frac{4BK^2}{A}\right)}{\mathrm{log}(\beta )}$$
proving Claim 5.9 and therefore completing the proof of Proposition 5.8. ∎
### 5.5 Interlude: The induced boundary map
The *upper boundary* $`^uG_M`$ is defined to be the leaf space of the weak stable foliation; this leaf space is identified with $`V^+`$. The *lower boundary* $`_{\mathrm{}}G_M`$ is the leaf space of the weak unstable foliation, identified with $`V^{}`$. The internal boundary $`_{\mathrm{int}}G_M`$ is defined as
$$_{\mathrm{int}}G_M=_{\mathrm{}}G_M\times ^uG_M=V^{}\times V^+𝐑^n/V^0$$
which is identified with the leaf space of the center foliation.
As a consequence of Proposition 5.4, a quasi-isometry $`\varphi :G_MG_L`$ which respects the transversely oriented horizontal foliations induces a bijection
$$_{\mathrm{int}}\varphi :_{\mathrm{int}}G_M_{\mathrm{int}}G_L$$
which preserves the factors, that is,
$$_{\mathrm{int}}\varphi =_l\varphi \times ^u\varphi :_{\mathrm{}}G_M\times ^uG_M_{\mathrm{}}G_L\times ^uG_L$$
Recall the 1-parameter family of metrics $`d_{M,t}`$ on $`𝐑^n`$ given by the quadratic form $`Q_{M,t}=(M^t)^TM^t`$. The internal boundary $`_{\mathrm{int}}G_M`$ is identified with $`𝐑^n/V^0`$ and with $`V^{}\times V^+`$, and we consider two 1-parameter families of metrics.
First, regarding points of $`𝐑^n/V^0`$ as affine subspaces parallel to $`V^0`$, there is a 1-parameter family of Hausdorff metrics induced from $`d_{M,t}`$ which we denote $`dh_{M,t}`$. Second, restrict the action of $`M^t`$ to the subspace $`V^{}\times V^+`$ to get a 1-parameter subgroup of $`\mathrm{GL}(V^{}\times V^+)`$, and by choosing a basis for $`V^{}\times V^+`$ we obtain a 1-parameter subgroup $`\widehat{M}^t`$ of $`\mathrm{GL}(k,𝐑)`$, where $`k`$ is the dimension of $`V^{}\times V^+`$. We obtain a 1-parameter family of metrics $`d_{\widehat{M},t}`$. There is a canonical identification $`V^{}\times V^+𝐑^n/V^0`$, and with respect to this identification the metrics $`d_{\widehat{M},t}`$ and $`dh_{M,t}`$ are bilipschitz equivalent, with a uniform bilipschitz constant independent of $`t`$.
Note that the absolute Jordan form of $`\widehat{M}`$ is identical with the nonunipotent part of the absolute Jordan form of $`M`$, and similarly for $`N`$.
###### Lemma 5.10.
Given two 1-parameter subgroups $`M^t,N^t`$ of $`\mathrm{GL}(n,𝐑)`$, for all $`K1`$, $`C,A0`$, there exist $`K^{}1,C^{}0`$ with the following properties. If $`\varphi :G_MG_N`$ is a $`K,C`$ quasi-isometry which coarsely respects the transversely oriented horizontal foliations, with Hausdorff constant $`A`$, then for every $`t𝐑`$ the induced map $`_{\mathrm{int}}\varphi :_{\mathrm{int}}G_M_{\mathrm{int}}G_N`$ is a $`K^{},C^{}`$ quasi-isometry from the metric $`d_{\widehat{M},t}`$ to the metric $`d_{\widehat{N},h(t)}`$.
###### Proof.
With what we know, the proof is mostly a matter of chasing through definitions.
The quasi-isometry $`\varphi `$ is a bounded distance from a quasi-isometry $`\psi :G_MG_N`$ which takes the horizontal leaf $`P_t`$ to the horizontal leaf $`P_{h(t)}^{}`$, and which simultaneously takes center leaves of $`G_M`$ to center leaves of $`G_N`$. Now restrict the center foliations of $`G_M,G_N`$ to $`P_t`$, $`P_{h(t)}^{}`$, and denote the respective leaf spaces as $`Q_t,Q_{h(t)}^{}`$.
In order to apply Lemma 2.1, consider each horizontal leaf $`P_t`$ of $`G_M`$ as a geodesic metric space, with respect to the Riemannian metric induced by restriction from $`G_M`$. The inclusion map $`P_tG_M`$ is evidently $`(1,0)`$ coarsely lipschitz, and it is uniformly proper, with a uniformity function $`s(r)=a^r`$ where $`a>1`$ is larger than the maximum of the absolute values of all eigenvalues of $`M`$ and their multiplicative inverses. Note in particular that the coarse Lipschitz constants and the uniformity functions of the maps $`P_tG_M`$ depend only on $`K,C,A`$ and on the matrix $`M`$, but not on $`t`$. Similar remarks apply to the inclusion map $`P_{h(t)}^{}G_N`$. Applying Lemma 2.1, restricting $`\psi `$ to $`P_t`$ results in a map $`\psi _t:P_tP_{h(t)}^{}`$ which is a quasi-isometry. There is in turn an induced map $`\theta _t:Q_tQ_{h(t)}^{}`$ which is a quasi-isometry with respect to the associated Hausdorff metric. The quasi-isometry constants of the maps $`\psi _t`$ and $`\theta _t`$ depend only on $`K,C,A`$.
Now consider the coordinate identifications $`G_M𝐑^n\times 𝐑`$, $`G_N𝐑^n\times 𝐑`$. By construction of the left invariant metrics, for each $`t`$ the space $`P_t`$ is identified with $`𝐑^n\times t𝐑^n`$ with the metric $`d_{M,t}`$, and the space $`P_{h(t)}^{}`$ is identified with $`𝐑^n`$ with metric $`d_{N,h(t)}`$, and so the maps $`\psi _t:𝐑^n𝐑^n`$ are uniform quasi-isometries from $`d_{M,t}`$ to $`d_{N,h(t)}`$ for all $`t`$. Also, $`Q_t`$ is identified with $`𝐑^n/V_M^0`$ with the associated Hausdorff metric $`dh_{M,t}`$, and $`Q_{h(t)}^{}`$ is identified with $`𝐑^n/V_N^0`$ with the associated Hausdorff metric $`dh_{N,h(t)}`$, and so the maps $`\theta _t:𝐑^n/V_M^0𝐑^n/V_N^0`$ are uniform quasi-isometries from $`dh_{M,t}`$ to $`dh_{N,h(t)}`$ for all $`t`$. This implies that $`\theta _t:V_M^{}\times V_M^+V_N^{}\times V_N^+`$ is a quasi-isometry from $`d_{\widehat{M},t}`$ to $`d_{\widehat{N},t}`$ for all $`t`$. But for all $`t`$ the map $`\theta _t`$ is identical to $`_{\mathrm{int}}\varphi :_{\mathrm{int}}G_M_{\mathrm{int}}G_N`$, proving the lemma. ∎
### 5.6 Step 3: Reduction to Theorem 5.11 on 1-parameter subgroup rigidity
Assume the hypotheses of Theorem 5.2, namely that we have 1-parameter subgroups $`M^t,N^t`$, and a quasi-isometry $`\varphi :G_MG_N`$ which coarsely respects the transversely oriented horizontal foliations. Applying Proposition 5.8, there is an induced time change of the form $`h(t)=mt+b`$ with $`m>0`$. Applying Proposition 4.1, there is a horizontal respecting quasi-isometry $`G_NG_{N^m}`$ with an induced time change of the form $`ss/m`$. By composition we obtain a horizontal respecting quasi-isometry $`G_MG_{N^m}`$ with an induced time change of the form $`tt+b^{}`$. Changing the coordinates in $`G_M`$ by a translation of the time coordinate $`t`$, we have a horizontal respecting quasi-isometry $`G_MG_{N^m}`$ for which the identity map $`tt`$ is an induced time change. Applying Lemma 5.10, we obtain a map $`_{\mathrm{int}}\varphi :𝐑^n𝐑^n`$ which, for each $`t`$, is a $`(K^{},C^{})`$-quasi-isometry from $`d_{\widehat{M},t}`$ to $`d_{\widehat{N}^m,t}`$.
Now apply the following theorem (with $`N`$ in place of $`N^m`$), which will be proved in the next section:
###### Theorem 5.11 (1-parameter subgroup rigidity).
Let $`M^t,N^t`$ be 1-parameter subgroups of $`\mathrm{GL}(n,𝐑)`$, such that $`M=M^1`$ and $`N=N^1`$ have no eigenvalues on the unit circle. If there exists a bijection $`f:𝐑^n𝐑^n`$ and constants $`K1,C0`$ such that for each $`t𝐑`$ and $`p,q𝐑^n`$ we have
$$C+\frac{1}{K}d_{M,t}(p,q)d_{N,t}(f(p),f(q))Kd_{M,t}(p,q)+C$$
then $`M`$ and $`N`$ have the same absolute Jordan form.
Returning to the previous discussion, this theorem allows us to conclude that $`\widehat{M}`$ and $`\widehat{N}^m`$ have the same absolute Jordan form, and so the nonunipotent parts of the absolute Jordan forms of $`M,N^m`$ are identical. We have already proved in Corollary 5.6 that the unipotent parts are identical, and so $`M`$ and $`N^m`$ have the same absolute Jordan forms, finishing the proof of Theorem 5.2. ∎
## 6 Dynamics of $`G_M`$, Part II: <br>1-parameter subgroup rigidity
In this section we give a proof of Theorem 5.11.
Let $`M^t`$, $`N^t`$ be 1-parameter subgroups of $`\mathrm{GL}(n,𝐑)`$ with no eigenvalues on the unit circle. Let $`M^t=\overline{M}^tP^t`$, $`N^t=\overline{N}^tQ^t`$ be the real Jordan forms, so $`\overline{M}`$ and $`\overline{N}`$ have all positive eigenvalues, none equal to $`1`$. Let $`f:𝐑^n𝐑^n`$ be a bijection which satisfies
$$C+\frac{1}{K}d_{M,t}(p,q)d_{N,t}(f(p),f(q))Kd_{M,t}(p,q)+C$$
(6.1)
for all $`t𝐑,p,q𝐑^n`$.
The bijection $`f:𝐑^n𝐑^n`$ must in fact be a homeomorphism. To see why, for each $`p𝐑^n`$, $`R>0`$, $`T>0`$ let
$$F_{p,R}(T)=\{q𝐑^n|d_{M,t}(p,q)<R\text{for all}t(T,T)\}$$
In other words, $`F_{p,R}(T)`$ is the intersection of open balls of radius $`R`$ about $`p`$ in each of the metrics $`d_{M,t}`$, for $`t(T,T)`$. Since the eigenvalues of $`\overline{M}`$ are all positive real numbers, none equal to $`1`$, it follows from Proposition 3.2 that for each $`p𝐑^n`$ and each $`R>0`$ the collection of sets $`F_{p,R}(T)`$ as $`T`$ ranges in $`(0,\mathrm{})`$ is a neighborhood basis for $`p`$, in the standard topology on $`𝐑^n`$. We define a similar neighborhood basis using matrix $`N`$, denoted $`G_{p,R}(T)`$. Since $`f(F_{p,R}(T))G_{f(p),KR+C}(T)`$ for each $`p𝐑^n,R>0,T>0`$, it follows that $`f`$ is continuous. The same argument applies to $`f^1`$, and so $`f`$ is a homeomorphism.
The idea of the proof of Theorem 5.11 is to show that $`f`$ respects certain “flags of foliations” which are closely related to the Jordan decompositions of $`𝐑^n`$ with respect to $`M^t`$ and $`N^t`$. We begin by setting up the notation needed to define and study these foliations.
###### Definition (Flags of foliations).
If $`V`$ is a vector subspace of $`𝐑^n`$, define a foliation $`(V)`$ of $`𝐑^n`$ whose leaves are the affine subspaces of $`𝐑^n`$ parallel to $`V`$. Given a flag of subspaces $`V_1\mathrm{}V_r`$, it follows that if $`1i<jr`$ then each leaf of $`(V_i)`$ is contained in some leaf of $`(V_j)`$; we denote this relation by saying that $`(V_1)\mathrm{}(V_r)`$ is a *flag of foliations* of $`𝐑^n`$.
Recall the root space decompositions of $`𝐑^n`$ with respect to $`\overline{M}`$ and $`\overline{N}`$. We denote the eigenvalues of $`\overline{M}`$ and $`\overline{N}`$ by
$$0<\mu _m^{}<\mathrm{}<\mu _1^{}<1<\mu _1^+<\mathrm{}<\mu _r^+$$
and
$$0<\nu _n^{}<\mathrm{}<\nu _1^{}<1<\nu _1^+<\mathrm{}<\nu _s^+$$
respectively. The corresponding root space decompositions are denoted
$$V_m^{}\mathrm{}V_1^{}V_1^+\mathrm{}V_r^+$$
and
$$W_n^{}\mathrm{}W_1^{}W_1^+\mathrm{}W_s^+.$$
As in section 4 we set
$`V^{}=V_m^{}\mathrm{}V_1^{},`$ $`V^+=V_1^+\mathrm{}V_r^+`$
$`W^{}=W_n^{}\mathrm{}W_1^{},`$ $`W^+=W_1^+\mathrm{}W_s^+`$
Define the *root space flags*
$`U_i^{}`$ $`=V_i^{}\mathrm{}V_1^{},i=1,\mathrm{},m`$
$`U_j^+`$ $`=V_1^+\mathrm{}V_j^+,j=1,\mathrm{},r`$
$`Y_i^{}`$ $`=W_i^{}\mathrm{}W_1^{},i=1,\mathrm{},n`$
$`Y_j^+`$ $`=W_1^+\mathrm{}W_j^+,j=1,\mathrm{},s`$
and by convention we take $`U_0^{}`$, $`U_0^+`$, $`Y_0^{}`$, $`Y_0^+`$ each to be the trivial subspace. Associated to the root space flags we have *root space foliation flags*
$`(U_1^{})`$ $`\mathrm{}(U_m^{})=(V^{})`$
$`(U_1^+)`$ $`\mathrm{}(U_r^+)=(V^+)`$
$`(Y_1^{})`$ $`\mathrm{}(Y_n^{})=(W^{})`$
$`(Y_1^+)`$ $`\mathrm{}(Y_s^+)=(W^+)`$
#### Step 1: $`f`$ respects contracting and expanding foliations.
First we show that $`f((V^{}))=(W^{})`$ and $`f((V^+))=(W^+)`$.
Given $`p,q𝐑^n`$ we have the following chain of equivalences:
1. $`p,q`$ are in the same leaf of $`(V^+)`$.
2. $`d_{M,t}(p,q)=M^t(pq)0`$ as $`t+\mathrm{}`$.
3. $`d_{M,t}(p,q)`$ is bounded for $`t[0,+\mathrm{})`$.
4. $`d_{N,t}(f(p),f(q))`$ is bounded for $`t[0,+\mathrm{})`$.
5. $`d_{N,t}(f(p),f(q))=N^t(f(p)f(q))0`$ as $`t+\mathrm{}`$.
6. $`f(p),f(q)`$ are in the same leaf of $`(W^+)`$.
The equivalence of (1–3) follows from Proposition 3.2, and similarly for (4–6). The equivalence of (3) and (4) follows from (6.1). This shows $`f((V^+))=(W^+)`$. A similar argument with $`t(\mathrm{},0]`$ shows $`f((V^{}))=(W^{})`$.
#### Step 2: $`f`$ respects root space foliation flags.
Next we show:
###### Claim 6.1.
$`f:𝐑^n𝐑^n`$ respects the root space foliation flags, and corresponding root spaces have the same eigenvalues. More precisely we have:
1. $`r=s`$.
2. $`\mu _j^+=\nu _j^+`$ for $`j=1,\mathrm{},r`$.
3. $`f((U_j^+))=(Y_j^+)`$ for $`j=1,\mathrm{},r`$.
4. $`m=n`$.
5. $`\mu _i^{}=\nu _i^{}`$ for $`i=1,\mathrm{},m`$.
6. $`f((U_i^{}))=(Y_i^{})`$ for $`i=1,\mathrm{},m`$.
It follows that $`M,N`$ have the same eigenvalues with the same multiplicities.
We give the proof of 1,2,3; the proof of 4,5,6 is similar.
We know by Step 1 that $`f((V^+))=(W^+)`$. Consider points $`p,q`$ in the same leaf of $`(V^+)`$, so $`f(p),f(q)`$ are in the same leaf of $`(W^+)`$. From Proposition 3.2 it follows that as $`t\mathrm{}`$ both of the quantities $`d_{M,t}(p,q)`$ and $`d_{N,t}(f(p),f(q))`$ approach $`+\mathrm{}`$. It follows that for sufficiently large $`t`$, in the inequality (6.1) we can absorb the additive constant $`C`$, yielding
$$\frac{1}{K+1}d_{M,t}(p,q)d_{N,t}(f(p),f(q))(K+1)d_{M,t}(p,q)$$
(6.2)
Define displacement vectors $`v=pq`$, $`w=f(p)f(q)`$. Taking natural logarithms, dividing by $`t`$, and taking limsup, we have
$`\underset{t\mathrm{}}{lim\; sup}{\displaystyle \frac{\mathrm{log}\left(d_{M,t}(p,q)\right)}{t}}`$ $`=\underset{t\mathrm{}}{lim\; sup}{\displaystyle \frac{\mathrm{log}\left(d_{N,t}(f(p),f(q))\right)}{t}}`$
$`\underset{t+\mathrm{}}{lim\; sup}{\displaystyle \frac{\mathrm{log}M^tv}{t}}`$ $`=\underset{t+\mathrm{}}{lim\; sup}{\displaystyle \frac{\mathrm{log}N^tw}{t}}`$ (6.3)
To evaluate these limits, let $`I(p,q)=I(v)`$ be the unique integer such that
$$vU_{I(v)}^+U_{I(v)1}^+$$
or, equivalently, the unique integer such that $`p,q`$ are in the same leaf of $`(U_{I(p,q)}^+)`$ but not in the same leaf of $`(U_{I(p,q)1}^+)`$ (recall the convention that $`U_0^+=0`$, and so $`I(p,q)=0`$ if and only if $`p=q`$). Define $`J(f(p),f(q))=J(w)`$ similarly by
$$wY_{J(w)}^+Y_{J(w)1}^+$$
Applying Proposition 3.2 we have
$`\underset{t+\mathrm{}}{lim\; sup}{\displaystyle \frac{\mathrm{log}M^tv}{t}}`$ $`=\mu _{I(v)}^+`$
$`\underset{t+\mathrm{}}{lim\; sup}{\displaystyle \frac{\mathrm{log}N^tw}{t}}`$ $`=\nu _{J(w)}^+`$
and so by (6.3) we have
$$\mu _{I(p,q)}^+=\mu _{I(v)}^+=\nu _{J(w)}^+=\nu _{J(f(p),f(q))}^+$$
Since $`f`$ is a bijection from each leaf of $`(V^+)`$ to some leaf of $`(W^+)`$, items (1) and (2) of Claim 6.1 now follow, and it also follows that
$$I(p,q)=J(f(p),f(q))$$
for all $`p,q`$ contained in the same leaf of $`(V^+)`$.
We now prove item (3) of Claim 6.1 by induction on $`j`$. If $`p,q`$ are in the same leaf of $`(U_1^+)`$ then $`I(p,q)=1`$ and so $`J(p,q)=1`$ which implies that $`f(p),f(q)`$ are in the same leaf of $`(Y_1^+)`$. A similar argument with $`f^1`$ proves that $`f((U_1^+))=(Y_1^+)`$, proving the base step of the induction. Now assume that $`f((U_j^+))=(Y_j^+)`$, and suppose $`p,q`$ are in the same leaf of $`(U_{j+1}^+)`$. There are two cases to consider. If $`p,q`$ lie in the same leaf of $`(U_j^+)`$ then by the induction hypothesis $`f(p),f(q)`$ lie in the same leaf of $`(Y_j^+)`$, in particular they lie in the same leaf of $`(Y_{j+1}^+)`$. If $`p,q`$ do not lie in the same leaf of $`(U_j^+)`$ then $`I(p,q)=j+1`$ and so $`J(f(p),f(q))=j+1`$ and so $`f(p),f(q)`$ lie on the same leaf of $`(Y_{j+1}^+)`$. A similar argument with $`f^1`$ shows that $`f((U_{j+1}^+))=Y_{j+1}^+`$, completing the induction.
As mentioned earlier, (4–6) are proved similarly, completing the proof of Claim 6.1.
#### Step 3: $`f`$ respects Jordan foliation flags.
From Step 2, for each fixed $`j=1,\mathrm{},r`$ the matrices $`M,N`$ have $`\mu _j^+`$ root spaces $`V_j^+,W_j^+`$ respectively. As part of their root space flags we have
$`U_j^+`$ $`=U_{j1}^+V_j^+`$
$`Y_j^+`$ $`=Y_{j1}^+W_j^+`$
Let $`c_j`$ be the index of nilpotency of $`\mu _jIM`$, and let $`d_j`$ be the index of nilpotency of $`\mu _jIN`$. Then we have Jordan filtrations
$`V_{j,0}^+\mathrm{}V_{j,c_j}^+`$ $`=V_j^+`$
$`W_{j,0}^+\mathrm{}W_{j,d_j}^+`$ $`=W_j^+`$
and we set $`U_{j,k}^+=U_{j1}^+V_{j,k}^+`$ and $`Y_{j,k}^+=Y_{j1}^+W_{j,k}^+`$, yielding subspace flags
$`U_{j1}^+U_{j,0}^+\mathrm{}U_{j,c_j1}^+`$ $`=U_j^+`$
$`Y_{j1}^+Y_{j,0}^+\mathrm{}Y_{j,d_j1}^+`$ $`=Y_j^+`$
Corresponding to these subspace flags are foliation flags,
$`(U_{j1}^+)(U_{j,0}^+)\mathrm{}(U_{j,c_j1}^+)`$ $`=(U_j^+)`$
$`(Y_{j1}^+)(Y_{j,0}^+)\mathrm{}(Y_{j,d_j1}^+)`$ $`=(Y_j^+)`$
called the *expanding Jordan foliation flags* associated to the corresponding root space foliations $`(U_j^+)`$, $`(Y_j^+)`$ respectively. The *contracting Jordan foliation flags* associated to each root space foliation $`(U_i^{})`$, $`(Y_i^{})`$ are similarly defined.
###### Claim 6.2.
$`f:𝐑^n𝐑^n`$ respects the Jordan foliation flags associated to corresponding root space foliations. More precisely, for each $`j=1,\mathrm{},r`$ we have:
1. $`c_j=d_j`$.
2. $`f((U_{j,k}^+))=(Y_{j,k}^+)`$ for $`k=0,\mathrm{},c_j1`$.
and similarly for the contracting Jordan foliation flags.
From this claim, for each $`j=1,\mathrm{},r`$ it immediately follows that $`\overline{M},\overline{N}`$ have the same Jordan blocks with eigenvalue $`\mu _j^+`$, and so the expanding parts of the Jordan forms for $`\overline{M},\overline{N}`$ are identical; similarly for the contracting parts. Since $`M,N`$ have no eigenvalues on the unit circle, it now follows that $`M,N`$ have the same absolute Jordan forms, completing the proof of Theorem 5.11.
###### Proof of Claim 6.2.
Consider $`p,q𝐑^n`$ in the same leaf of $`(U_j^+)`$ but not in the same leaf of $`(U_{j1}^+)`$, and so $`f(p),f(q)`$ are in the same leaf of $`(Y_j^+)`$ but not in the same leaf of $`(Y_{j1}^+)`$. Define displacement vectors $`v=pq`$, $`w=f(p)f(q)`$, and so $`vU_j^+U_{j1}^+`$ and $`wY_j^+Y_{j1}^+`$. We know that
$$\underset{t+\mathrm{}}{lim\; sup}\frac{\mathrm{log}M^tv}{t}=\frac{\mathrm{log}N^tw}{t}=\mu _j^+$$
We also know that (6.2) is true for $`t`$ sufficiently close to $`\mathrm{}`$, and so for $`t`$ sufficiently close to $`+\mathrm{}`$ we have
$$\frac{1}{K+1}M^tvN^tw(K+1)M^tv$$
(6.4)
By induction on $`k=0,1,\mathrm{}`$, we shall prove that $`vU_{j,k}^+`$ if and only if $`wY_{j,k}^+`$, or equivalently that $`f((U_{j,k}^+))=(Y_{j,k}^+)`$.
For the basis step $`k=0`$, divide the inequality (6.4) by $`\mu ^t`$ to obtain, for all $`t`$ sufficiently close to $`+\mathrm{}`$:
$$\frac{1}{K+1}\frac{M^tv}{\mu ^t}\frac{N^tw}{\mu ^t}(K+1)\frac{M^tv}{\mu ^t}$$
(6.5)
By the Exponential Lower Bound and the Exponential$``$Polynomial Upper Bound of Proposition 3.2, the quantity $`{\displaystyle \frac{M^tv}{\mu ^t}}`$ is bounded for $`t0`$ if and only if $`vU_{j,0}^+`$; and the quantity $`{\displaystyle \frac{N^tw}{\mu ^t}}`$ is bounded on $`t0`$ if and only if $`wY_{j,0}^+`$. However by (6.5) the boundedness of these two quantities on $`t0`$ are equivalent.
For the induction step, assume that $`f((U_{j,k1}^+))=(Y_{j,k1}^+)`$, that is, $`vU_{j,k1}^+`$ if and only if $`wY_{j,k1}^+`$. We must prove that $`vU_{j,k}^+U_{j,k1}^+`$ if and only if $`wY_{j,k}^+Y_{j,k1}^+`$. From (6.4), for $`t`$ sufficiently close to $`+\mathrm{}`$ we have
$$\frac{1}{K+1}\frac{M^tv}{\mu ^tt^k}\frac{N^tw}{\mu ^tt^k}(K+1)\frac{M^tv}{\mu ^tt^k}$$
(6.6)
and
$$\frac{1}{K+1}\frac{M^tv}{\mu ^tt^{k1}}\frac{N^tw}{\mu ^tt^{k1}}(K+1)\frac{M^tv}{\mu ^tt^{k1}}$$
(6.7)
By the Exponential$``$Polynomial Upper and Lower Bounds of Proposition 3.2, the following two statements are equivalent:
* $`vU_{j,k}^+U_{j,k1}^+`$
* For $`t0`$, the quantity $`{\displaystyle \frac{M^tv}{\mu ^tt^k}}`$ is bounded, but the quantity $`{\displaystyle \frac{M^tv}{\mu ^tt^{k1}}}`$ is not bounded.
Similarly, the following two statements are equivalent:
* $`wY_{j,k}^+Y_{j,k1}^+`$.
* For $`t0`$, the quantity $`{\displaystyle \frac{N^tw}{\mu ^tt^k}}`$ is bounded, but the quantity $`{\displaystyle \frac{N^tw}{\mu ^tt^{k1}}}`$ is not bounded.
But by inequalities (6.6) and (6.7), statements (2) and (4) are equivalent, and so statements (1) and (3) are equivalent, completing the inductive proof of item 2 of Claim 6.2 for all $`k0`$.
The foliation flag $`(U_{j,0}^+)\mathrm{}(U_{j,k}^+)\mathrm{}`$ must terminate at $`(U_j^+)`$ for the same value of $`k`$ for which the flag $`(Y_{j,0}^+)\mathrm{}(Y_{j,k}^+)\mathrm{}`$ terminates at $`(Y_j^+)`$, proving that $`c_j=d_j`$, and completing the proof of Claim 6.2. ∎
Our proof of Theorem 5.11 actually provides for some regularity of $`f`$. We record the statement here, although it is not used at all in this paper.
###### Proposition 6.3 (Regularity).
With the assumptions as in Theorem 5.11, $`f`$ is a homeomorphism which respects the contracting and expanding root space foliation flags of $`\overline{M},\overline{N}`$, and for each corresponding pair of root space foliations $`f`$ also respects the associated Jordan foliation flags.∎
###### Remark.
Even stronger regularity properties should hold. For instance, $`f`$ should satisfy lipschitz conditions in directions parallel to a root space, by arguments similar to the results of \[FM98\]. Understanding what happens transverse to root spaces will require new ideas.
## 7 Quasi-isometries of $`\mathrm{\Gamma }_M`$ via Coarse Topology
Recall the notation for abelian-by-cyclic Lie groups: given $`M\mathrm{GL}_\times (m,𝐑)`$, a 1-parameter subgroup $`M^t\mathrm{GL}(m,𝐑)`$ with $`M^1=M`$ determines a Lie group denoted $`G_M=𝐑^m_M𝐑`$.
This entire section will be devoted to a proof of the following.
###### Proposition 7.1 (Induced quasi-isometries of $`G_M`$).
Consider integral matrices $`M\mathrm{GL}_\times (m,𝐑)`$, $`N\mathrm{GL}_\times (n,𝐑)`$ and suppose that $`detM,detN>1`$. If there exists a quasi-isometry $`f:\mathrm{\Gamma }_M\mathrm{\Gamma }_N`$ then $`m=n`$ and there exists a quasi-isometry $`\varphi :G_MG_N`$ which coarsely respects horizontal foliations and their transverse orientations. Furthermore, all associated constants for $`\varphi `$ depend only on those for $`f`$.
### 7.1 A geometric model for $`\mathrm{\Gamma }_M`$
Let $`M\mathrm{GL}_\times (m,𝐑)`$ be an integral matrix lying on a 1-parameter subgroup $`M^t`$ of $`\mathrm{GL}(m,𝐑)`$ with $`M^1=M`$ and with associated Lie group $`G_M`$. We assume that $`detM>1`$ and we denote $`d=detM`$.
We start by constructing a contractible, $`(m+1)`$-dimensional metric complex $`X_M`$ on which $`\mathrm{\Gamma }_M`$ acts properly discontinuously and cocompactly by isometries, and so the group $`\mathrm{\Gamma }_M`$ will be quasi-isometric to the geodesic metric space $`X_M`$.
The description of $`\mathrm{\Gamma }_M`$ as an ascending HNN extension shows that $`\mathrm{\Gamma }_M`$ is the fundamental group of the mapping torus of an injective endomorphism of the $`m`$-dimensional torus. Let $`X_M`$ be the universal cover of this mapping torus. Topologically, there is a fibration
where $`T_M`$ is the homogeneous directed tree with one edge coming into each vertex and $`d=detM`$ edges going out of each vertex. Hence $`X_M`$ is a topological product $`X_M𝐑^{n1}\times T_M`$.
The action of $`\mathrm{\Gamma }_M`$ on $`X_M`$ by deck tranformations induces an action of $`\mathrm{\Gamma }_M`$ on $`T_M`$. This action is equivalent to the usual action of the HNN extension $`\mathrm{\Gamma }_M`$ on its Bass-Serre tree $`T_M`$.
Before constructing a metric on $`X_M`$, let us describe the essential properties of such a metric. These are best described by giving the isometry types of natural subcomplexes of $`X_M`$.
###### Definition (Doubled horoballs).
We define a doubled $`G_M`$ horoball, denoted $`H_M`$, to be the metric space obtained by identifying two copies of $`\{(x,t)G_M|t0\}`$ along $`\{(x,0)G_M\}`$, endowed with the path metric.
###### Definition (Hyperplanes in $`X_M`$).
Let $`P_{\mathrm{}}=\pi _M^1(\mathrm{})`$, where $`\mathrm{}`$ is a bi-infinite line in the directed tree $`T_M`$. We call $`P_{\mathrm{}}`$ a hyperplane in $`X_M`$. There are two cases to consider:
* $`\mathrm{}`$ is coherently oriented in $`T_M`$. In this case $`P_{\mathrm{}}`$ is isometric to $`G_M`$, and we call $`P_{\mathrm{}}`$ a coherent hyperplane in $`X_M`$.
* $`\mathrm{}`$ is not coherently oriented in $`T_M`$, and thus switches orientation precisely once. In this case $`P_{\mathrm{}}`$ is isometric to $`H_M`$, and we call $`P_{\mathrm{}}`$ an incoherent hyperplane in $`X_M`$.
This definition nearly determines a metric on $`X_M`$. To specify a metric on $`X_M`$, one proceeds as follows. Fix a path metric on $`T_M`$ so that each edge has length $`1`$. Fix a base vertex on $`T_M`$. These choices determine a unique height function $`T_M𝐑`$ taking the base vertex to the origin, and taking each edge to a segment of length $`1`$ via an orientation preserving isometry. We have also defined a height function $`G_M𝐑`$. Note that the height function on $`G_M`$ was previously called the “time function”; we will use both terms.
The complex $`X_M`$ is the fiber product of the two height functions $`T_M𝐑`$, $`G_M𝐑`$, as shown in the following diagram:
There are induced projections $`g_M:X_MG_M`$ and $`\pi _M:X_MT_M`$, and an induced height function $`X_M𝐑`$. There is a unique path metric on $`X_M`$ so that each continuous cross section $`G_MX_M`$ of $`g_M`$ is a path-isometric embedding; and hence each coherent hyperplane in $`X_M`$ is an isometrically embedded copy of $`G_M`$.
###### Definition (horizontal leaf).
A horizontal leaf $`L`$ in $`X_M`$ is a subset of the form $`L=\pi _M^1(v)`$ where $`vT_M`$.
Note that the collection of horizontal leaves on $`X_M`$, equipped with the Hausdorff metric, forms a metric space which is isometric to $`T_M`$ via the projection map $`\pi _M:X_MT_M`$.
Note that each hyperplane in $`X_M`$ comes equipped with a foliation by horizontal leaves. For coherent hyperplanes $`P`$ in $`X_M`$, which are isometric to $`G_M`$, the notion of horizontal leaf in $`P`$ coincides with that of a horizontal leaf in $`G_M`$, given in §5.1.
### 7.2 Proof of Proposition 7.1 on induced quasi-isometries of $`G_M`$
Let $`M,N`$ be as in the statement of the proposition.
We begin by showing that $`M`$ and $`N`$ have the same size. Suppose that $`M\mathrm{GL}(m,𝐑)`$ and $`N\mathrm{GL}(n,𝐑)`$. In §7.1 we constructed finite classifying spaces for $`\mathrm{\Gamma }_M`$ and $`\mathrm{\Gamma }_N`$ of dimensions $`m+1,n+1`$ respectively, and by Lemma 5.2 of \[FM99b\] these numbers are the virtual cohomological dimensions of $`\mathrm{\Gamma }_M,\mathrm{\Gamma }_N`$. By a result of Block-Weinberger \[BW97\] and Gersten \[Ger93\], virtual cohomological dimension is a quasi-isometry invariant for groups with finite classifying spaces. It follows that $`m=n`$.
Now $`\mathrm{\Gamma }_M`$ acts properly discontinuously, freely, and cocompactly on $`X_M`$. This action is by isometries, because $`\mathrm{\Gamma }_M`$ acts on $`G_M`$, on $`T_M`$, and on $`𝐑`$ by isometries, and the fiber product diagram is equivariant with respect to these actions. It follows that $`\mathrm{\Gamma }_M`$ in any word-metric is quasi-isometric to $`X_M`$. Henceforth we will freely interchange $`\mathrm{\Gamma }_M`$ and $`X_M`$ when discussing quasi-isometry type. The same discussion applies to $`\mathrm{\Gamma }_N`$ and $`X_N`$, and so the quasi-isometry $`f:\mathrm{\Gamma }_M\mathrm{\Gamma }_N`$ gives a quasi-isometry (perhaps with bigger constants) $`f:X_MX_N`$.
Proposition 7.1 generalizes the case when $`M`$ and $`N`$ are $`1\times 1`$ matrices, done in §4 and §5 of \[FM98\]. The proof here is more difficult, and the steps must be proved in different order. In Steps 1 and 2 we prove (in a more general context; see Theorem 7.7) that a quasi-isometry $`X_MX_N`$ coarsely respects hyperplanes and horizontal sets. However, we must still distinguish between coherent and incoherent hyperplanes. This is easy in the $`1\times 1`$ case handled in \[FM98\], where $`G_M`$ and $`G_N`$ are (scaled versions of) $`𝐇^2`$, and a doubled $`𝐇^2`$ horoball is evidently not quasi-isometric to $`𝐇^2`$. In general we are unable to distinguish the quasi-isometry types of coherent and incoherent hyperplanes. To get around this, in Step 3, Proposition 7.11, we prove that there is no horizontal respecting quasi-isometry between a coherent and an incoherent hyperplane.
##### Step 1. Quasi-isometrically embedded hyperplanes are close to hyperplanes:
Given integral matrices $`M,N\mathrm{GL}_\times (n,𝐑)`$, if $`P=G_M`$ or $`H_M`$, then for all $`K1,C0`$ there exists $`A0`$ such that if $`\varphi :PX_N`$ is a $`K,C`$-quasi-isometric embedding then there is a unique hyperplane $`QX_N`$ with $`d_{}(\varphi (P),Q)A`$.
This was proved for $`1\times 1`$ matrices in \[FM98\]. Our proof of Step 1, while following the same outline as the $`1\times 1`$ case, will actually apply in a much broader setting. The generalized versions of Steps 1 and 2, given in Theorem 7.3 and Theorem 7.7, are used for example in \[FM99a\] to study surface-by-free groups, and also in \[MSW\] to prove quasi-isometric rigidity theorems for various “homogeneous” graphs of groups (see the remark after Theorem 7.7).
The generalization of Step 1 given in Theorem 7.3 will require moving from the category of quasi-isometric embeddings into the category of uniformly proper embeddings. After a fair amount of work to establish the new setting, we then quote some theorems of coarse algebraic topology and follow the proof of \[FM98\].
Consider a finite graph $`\mathrm{\Gamma }`$ of finitely generated groups; each edge $`e`$ is oriented, with initial and final vertices $`i(e)`$, $`f(e)`$. We say that $`\mathrm{\Gamma }`$ is *geometrically homogeneous* if each edge-to-vertex injection is a quasi-isometry with respect to the word metric space, or equivalently, has finite index image. Ideally we would like to have a version of Step 1 for any geometrically homogeneous graph of groups in which each vertex and edge group is the fundamental group of a closed, aspherical $`n`$-manifold, or even more generally, an $`n`$-dimensional Poincaré duality group. This should come from a more careful reading of results in coarse algebraic topology such as \[KK\], but meanwhile we will use Theorems 7.5 and 7.6, which require us to impose additional assumptions on $`\mathrm{\Gamma }`$.
Suppose that we have a category $`𝒞`$ of aspherical, closed, smooth manifolds such that $`𝒞`$ is closed under finite coverings and satisfies *smooth rigidity*, meaning that any homotopy equivalence between manifolds in $`𝒞`$ is homotopic to a diffeomorphism. Such categories include: the $`n`$-torus, $`n1`$; hyperbolic surfaces; all other irreducible, nonpositively curved, locally symmetric spaces, by Mostow’s Rigidity Theorem \[Mos73\]; solvmanifolds, by earlier work of Mostow \[Mos54\]; nilmanifolds, by still earlier work of Malcev \[Mal49\]; and various generalizations due to Farrell and Jones \[FJ89\], \[FJ97\].
We shall assume that $`\mathrm{\Gamma }`$ is a geometrically homogeneous graph of groups where each vertex group $`\mathrm{\Gamma }_v`$ is the fundamental group of a manifold $`M_v`$ in the category $`𝒞`$. Construct a *graph of aspherical manifolds* $`M_\mathrm{\Gamma }`$, with fundamental group $`\pi _1\mathrm{\Gamma }`$, as follows. For each edge $`e`$, the two injections $`\mathrm{\Gamma }_e\mathrm{\Gamma }_{i(e)}`$, $`\mathrm{\Gamma }_e\mathrm{\Gamma }_{t(e)}`$ determine two finite covering spaces of $`M_v`$ each of whose fundamental group is identified with $`\mathrm{\Gamma }_e`$, and so we obtain a diffeomorphism between the two covering spaces; identify these covering spaces and let $`M_e`$ be the resulting smooth manifold. We have smooth, finite covering maps $`M_eM_{i(e)}`$, $`M_eM_{t(e)}`$ inducing the corresponding edge-to-vertex group injections. Form $`M_\mathrm{\Gamma }`$ from the disjoint union
$$\left(\underset{v}{}M_v\right)\left(\underset{e}{}M_e\times e\right)$$
by gluing $`M_e\times i(e)`$ to $`M_{i(e)}`$ and $`M_e\times f(e)`$ to $`M_{f(e)}`$ via the finite covering maps $`M_eM_{i(e)}`$ and $`M_eM_{f(e)}`$. From the construction of $`M_\mathrm{\Gamma }`$ we obtain a map $`M_\mathrm{\Gamma }\mathrm{\Gamma }`$, such that each fiber $`M_x`$, $`x\mathrm{\Gamma }`$, is a manifold in the category $`𝒞`$.
Let $`X_\mathrm{\Gamma }`$ be the universal cover of $`M_\mathrm{\Gamma }`$. There is a $`\mathrm{\Gamma }`$-equivariant fiber bundle $`X_\mathrm{\Gamma }T_\mathrm{\Gamma }`$ over the Bass-Serre tree $`T_\mathrm{\Gamma }`$ of $`\mathrm{\Gamma }`$, whose fiber is a contractible $`n`$-manifold. Any geodesic metric on $`M_\mathrm{\Gamma }`$ lifts to a $`\pi _1\mathrm{\Gamma }`$-equivariant geodesic metric on $`X_\mathrm{\Gamma }`$. Smoothness allows us to impose additional geometric structure on $`X_\mathrm{\Gamma }`$ which we now describe.
A geodesic metric space is *proper* if closed balls are compact. A *bounded-geometry metric simplicial complex* is a simplicial complex $`\mathrm{\Sigma }`$ equipped with a proper, geodesic metric such that for some constants $`0<C_1<C_2`$ each positive dimensional simplex has diameter between $`C_1`$ and $`C_2`$, and for some constant $`C>0`$ the link of each simplex has $`C`$ simplices. A subset $`S`$ of $`\mathrm{\Sigma }`$ is *rectifiable* if for any $`p,qS`$ there exists a path in $`S`$ between $`p`$ and $`q`$ which is rectifiable in $`\mathrm{\Sigma }`$, and which has the shortest $`\mathrm{\Sigma }`$-length among all paths in $`S`$ between $`p`$ and $`q`$. The length of such a path defines a geodesic metric on $`S`$. A *$`D`$-homotopy* in $`\mathrm{\Sigma }`$ is a homotopy whose tracks all have diameter $`D`$. The space $`\mathrm{\Sigma }`$ is *uniformly contractible* if there exists a function $`\delta :[0,\mathrm{})[0,\mathrm{})`$, such that for every bounded subset $`S\mathrm{\Sigma }`$, the inclusion map $`S\mathrm{\Sigma }`$ is $`\delta \left(\mathrm{diam}(S)\right)`$-homotopic to a constant map. More precisely we say that $`\mathrm{\Sigma }`$ is *$`\delta `$-uniformly contractible*.
Let $`T`$ be a bounded geometry, metric simplicial tree, let $`X`$ be a proper, geodesic metric, and let $`\pi :XT`$ be a surjective map. Denote $`X_A=\pi ^1(A)`$ for each $`AT`$. The map $`\pi `$ is called a *metric fibration* if:
1. $`X`$ is a uniformly contractible, bounded geometry, metric simplicial complex.
2. For each subtree $`T^{}T`$, the subset $`X_T^{}`$ is a subcomplex of $`X`$ and is rectifiable in $`X`$.
3. For each $`tT`$ the subspace $`X_t`$ is uniformly contractible and is a bounded geometry, metric simplicial complex, with bounded geometry constants and uniform contractiblity data independent of $`t`$.
4. The map $`\pi :XT`$ is distance nonincreasing.
5. There is a homeomorphism $`\mathrm{\Theta }:XF\times T`$ such that:
1. For all $`tT`$, $`\mathrm{\Theta }(X_t)=F\times t`$.
2. For all $`xF`$, the map $`Tx\times T\stackrel{\mathrm{\Theta }^1}{}X`$ is a locally isometric embedding.
3. There exists $`K1`$ such that for all edges $`e`$ of $`T`$ and $`te`$, the retraction $`r:et`$ induces a projection $`X_e\stackrel{\Theta }{}F\times e\stackrel{\text{Id}\times r}{}F\times t\stackrel{\mathrm{\Theta }^1}{}X_t`$ which is $`K`$-lipschitz.
Each fiber $`X_t`$, $`tT`$, is called a *horizontal leaf* in $`X`$. If $`L`$ is a bi-infinite line in $`T`$ then $`X_L`$ is called a *hyperplane* in $`X`$. Items 4 and 5b combine to show that the map of item 5b is an isometric embedding; the image $`\mathrm{\Theta }^1(x\times T)`$ is called a *vertical leaf* in $`X`$. For each subtree $`T^{}T`$, the closest point retraction $`r:TT^{}`$ induces a map
$$X\stackrel{\Theta }{}F\times T\stackrel{\text{Id}\times r}{}F\times T^{}\stackrel{\mathrm{\Theta }^1}{}X_T^{}$$
called *vertical projection* of $`X`$ to $`X_T^{}`$.
###### Remark.
Suppose $`\mathrm{\Gamma }`$ is a graph of groups taken from a category $`𝒞`$ as above. Let $`M_\mathrm{\Gamma }`$ and $`X_\mathrm{\Gamma }T_\mathrm{\Gamma }`$ be as constructed above starting from $`\mathrm{\Gamma }`$. Then elementary constructions produce a metric and a simplicial structure on $`M_\mathrm{\Gamma }`$ which lifts to a $`\mathrm{\Gamma }`$-equivariant metric and simplicial structure on $`X`$ such that $`XT_\mathrm{\Gamma }`$ is a metric fibration. Item (1) follows by compactness of $`M_\mathrm{\Gamma }`$.
###### Remark.
The definition has some redundancy: item (1) is a formal consequence of item (3), as can be seen by elementary but mildly tedious arguments. But by the previous remark we may dispense with these arguments for the examples at hand.
The following lemma, applied to a bi-infinite line in $`T`$, gives good geometric properties for hyperplanes:
###### Lemma 7.2.
If $`\pi :XT`$ is a metric fibration then there exist functions $`\delta ^{}:[0,\mathrm{})[0,\mathrm{})`$ and $`\rho :[0,\mathrm{})[0,\mathrm{})`$ with $`\underset{t\mathrm{}}{lim}\rho (t)=\mathrm{}`$, such that for any subtree $`T^{}T`$ we have:
* The embedding $`X_T^{}X`$ is $`\rho `$-uniformly proper.
* The geodesic metric space $`X_T^{}`$ is $`\delta ^{}`$-uniformly contractible.
###### Proof.
To prove (1), consider $`x,yX_T^{}`$, let $`D=d_X(x,y)`$, and let $`\gamma :[0,D]X`$ be a geodesic connecting $`x`$ and $`y`$. Let $`N_D(T^{})`$ be the $`D`$-neighborhood of $`T^{}`$ in $`T`$, so $`\gamma X_{N_D(T^{})}`$. Applying item 5 iteratively, projecting inward starting from the edges of $`N_D(T^{})`$ furthest from $`T^{}`$, it follows that vertical projection $`X_{N_D(T^{})}X_T^{}`$ distorts any distance $`r`$ by at worst $`K^Dr`$, and so $`d_{X_T^{}}(x,y)K^DD`$.
To prove (2), suppose that $`AX_T^{}`$ and $`\mathrm{diam}_{X_T^{}}(A)R`$ and so $`A`$ is $`R^{}`$-homotopic to a constant in $`X`$ where $`R^{}`$ depends on $`R`$ but not on $`A`$. This homotopy may then be mapped back to $`X_T^{}`$ by vertical projection, distorting diameters of homotopy tracks by an amount bounded in terms of $`R^{}`$ as we saw above. The result is an $`R^{\prime \prime }`$-homotopy of $`A`$ to a constant in $`T^{}`$, with $`R^{\prime \prime }`$ depending only on $`R`$ and not on $`A`$. ∎
Here is our generalization of Step 1. It applies to any metric fibration of the form $`X_\mathrm{\Gamma }T_\mathrm{\Gamma }`$, where $`\mathrm{\Gamma }`$ is a finite, geometrically homogeneous graph of fundamental groups of manifolds in any of the categories $`𝒞`$ described earlier.
###### Theorem 7.3.
Let $`\pi :XT`$ be a metric fibration whose fibers are contractible $`n`$-manifolds for some $`n`$. Let $`P`$ be a contractible $`(n+1)`$-manifold which is a uniformly contractible, bounded geometry, metric simplicial complex. Then for any uniformly proper embedding $`\varphi :PX`$, there exists a unique hyperplane $`QX`$ such that $`\varphi (P)`$ and $`Q`$ have finite Hausdorff distance in $`X`$. The bound on Hausdorff distance depends only on the metric fibration data for $`\pi `$, the uniform contractibility data and bounded geometry data for $`P`$, and the uniform properness data for $`\varphi `$.
###### Proof.
Uniqueness of $`Q`$ follows obviously from the fact that distinct hyperplanes in $`X`$ have infinite Hausdorff distance.
For existence of $`Q`$ we follow closely the proof of Proposition 4.1 of \[FM98\], concentrating on details needed to explicate the difference between the “quasi-isometric” setting of \[FM98\] and the present “uniformly proper” setting.
Using the bounded geometry of $`P`$, uniform contractibility of $`X`$, and uniform properness of $`\varphi `$, we may replace $`\varphi `$ by a continuous, uniformly proper map, moving values of $`\varphi `$ a bounded distance. Henceforth we shall assume $`\varphi `$ is continuous.
Pick a topologically proper embedding of $`T`$ in an open disc $`D`$. For each component $`U`$ of $`DT`$, the frontier of $`U`$ in $`D`$ is a bi-infinite line $`L(U)`$ in $`T`$. There is a homeomorphism of pairs $`(\overline{U},L(U))(L(U)\times [0,\mathrm{}),L(U)\times 0)`$.
Consider the topologically proper embedding $`X\stackrel{\Theta }{}F\times TF\times D`$. Note that $`F\times D`$ is a contractible $`(n+2)`$-manifold. For each component $`U`$ of $`DT`$ we have a homeomorphism
$`F\times \overline{U}`$ $`\stackrel{}{}F\times (L(U)\times [0,\mathrm{}))\stackrel{}{}(F\times L(U))\times [0,\mathrm{})`$
$`\underset{\mathrm{\Theta }\times \text{Id}}{\overset{}{}}X_{L(U)}\times [0,\mathrm{})`$
The frontier of this set in $`F\times D`$ is $`F\times L(U)X_{L(U)}`$. Put a product metric and a product simplicial structure on $`X_{L(U)}\times [0,\mathrm{})`$ and glue to $`F\times L(U)`$. Doing this for each $`U`$, we impose a proper geodesic metric on $`F\times D`$ for which the inclusion $`XF\times D`$ is an isometric embedding.
The simplicial structure on $`F\times D`$ evidently has bounded geometry. Also, the metric space $`F\times D`$ is uniformly contractible. To see this, let $`AF\times D`$ have diameter $`r`$. If $`AX\mathrm{}`$ then homotoping along product lines of $`X_{L(U)}\times [0,\mathrm{})`$ for each $`U`$ we obtain an $`r`$-homotopy of $`A`$ into $`F\times TX`$, and then we use uniform contractibility of $`X`$. Whereas if $`AX=\mathrm{}`$, then $`AF\times UX_{L(U)}\times (0,\mathrm{})`$ for some component $`U`$ of $`DT`$; there is an $`r`$-homotopy of $`A`$ into some $`X_{L(U)}\times x`$, and the latter is uniformly contractible by Lemma 7.2.
We now plug this setup into the coarse separation and packing methods of Farb–Schwartz \[FS96\] and Schwartz \[Sch96\]. We’ll use a generalization of the Coarse Separation Theorem with more easily applied hypotheses, due to Kapovich–Kleiner \[KK\]. We denote the $`r`$-ball about a subset $`A`$ of a metric space $`M`$ by $`B_r(A;M)`$. In a metric space $`Z`$, a subset $`UZ`$ is *deep in $`Z`$* if for each $`r>0`$ there exists $`xU`$ such that $`B_r(x;Z)U`$. A subset $`AZ`$ *coarsely separates* $`Z`$ if for some $`D>0`$ there are at least two components of $`ZN_D(A;Z)`$ which are deep in $`Z`$; the constant $`D`$ is called a *coarse separation constant* for $`A`$. Note that if subsets $`A`$ and $`B`$ of $`Z`$ have bounded Hausdorff distance from each other, then $`A`$ coarsely separates $`Z`$ if and only if $`B`$ does.
Here is an elementary consequence of the definitions:
###### Lemma 7.4.
Let $`f:XY`$ be a quasi-isometry between geodesic metric spaces. If $`AX`$ coarsely separates $`X`$ then $`f(A)`$ coarsely separates $`Y`$, with separation constant depending only on the quasi-isometry constants of $`f`$ and the separation constant for $`A`$. ∎
Here is the version of the Coarse Separation Theorem that we will use.
###### Theorem 7.5 (\[KK\]).
Let $`P`$ be a contractible $`(n+1)`$-manifold, $`Z`$ a contractible $`(n+2)`$-manifold, and suppose that $`P,Z`$ are uniformly contractible, bounded geometry, metric simplicial complexes. Let $`\mathrm{\Phi }:PZ`$ be a uniformly proper map. Then $`\mathrm{\Phi }(P)`$ coarsely separates $`Z`$, with coarse separation constant $`D`$ depending only on the uniform contractibility and bounded geometry data for $`P`$ and $`Z`$ and the uniform properness data for $`\mathrm{\Phi }`$. Moreover if $`\mathrm{\Phi }`$ is continuous then we may take $`D=0`$, that is, $`Z\mathrm{\Phi }(P)`$ has at least two components which are deep in $`Z`$. ∎
###### Remark.
In fact there are exactly two components of $`ZN_D(\mathrm{\Phi }(P);Z)`$ which are deep in $`Z`$ (see \[KK\]).
Following \[FS96\] we have a corollary:
###### Theorem 7.6 (Packing Theorem).
Let $`Q,P`$ be contractible $`(n+1)`$-manifolds, which are uniformly contractible, bounded geometry, metric simplicial complexes. Let $`\psi :QP`$ be a uniformly proper map. Then there exists $`R>0`$ such that $`N_R(\psi (Q);P)=P`$. The constant $`R`$ depends only on the uniform contractibility data and bounded geometry data for $`Q,P`$ and the uniform properness data for $`\psi `$.
###### Proof.
If no such $`R`$ exists then the image of the map $`Q\stackrel{𝜓}{}PP\times 𝐑`$ does not coarsely separate $`P\times 𝐑`$, violating Theorem 7.5. ∎
Continuing with the proof of Theorem 7.3, compose the continuous, uniformly proper map $`\varphi :PX`$ with the isometric embedding $`XF\times D`$ to obtain a continuous, uniformly proper map $`\mathrm{\Phi }:PF\times D`$. By the Coarse Separation Theorem it follows that $`(F\times D)\mathrm{\Phi }(P)`$ has at least two components which are deep in $`F\times D`$.
Now take the argument of \[FM98\], Step 1, pages 426–427, and apply it verbatim, to produce a hyperplane $`QX`$ such that $`Q\mathrm{\Phi }(P)`$. Next take the argument of Step 2, pages 427–428, and apply it verbatim, replacing “quasi-isometric embeddings” with “uniformly proper maps” and using the Packing Theorem above, to show the existence of $`R^{}`$ such that $`\varphi (P)N_R^{}(Q;X)`$, where $`R^{}`$ depends only on the metric fibration data for $`\pi `$, the uniform contractibility and bounded geometry data for $`P`$, and the uniform properness data for $`\varphi `$.
This finishes the proof of Theorem 7.3 and of Step 1. ∎
##### Step 2. A quasi-isometry takes hyperplanes and horizontal leaves in $`X_M`$ to hyperplanes and horizontal leaves in $`X_N`$:
Consider integral matrices $`M,N\mathrm{GL}_\times (n,𝐑)`$ with $`detM,detN>1`$, and let $`f:X_MX_N`$ be a quasi-isometric embedding. Then there is a constant $`A0`$, depending only on $`X_M,X_N`$ and the quasi-isometry constants of $`f`$, such that:
1. For each hyperplane $`PX_M`$ there exists a unique hyperplane $`QX_N`$ such that $`d_{}(f(P),Q)A`$.
2. For each horizontal leaf $`L`$ of $`X_M`$ there exists a horizontal leaf $`L^{}`$ of $`X_N`$ such that $`d_{}(f(L),L^{})A`$.
The proof of this step is the first place in our arguments where the assumption that $`detM,detN>1`$ is crucial. Again we will investigate this step in the general setting of metric fibrations over trees.
Consider a metric fibration $`\pi :XT`$. The tree $`T`$ is *bushy* if there exists a constant $`\beta `$ such that each point of $`T`$ is within distance $`\beta `$ of some vertex $`v`$ such that $`Tv`$ has at least 3 unbounded components. Note that if $`M`$ is an integer matrix in $`\mathrm{GL}_\times (n,𝐑)`$, and if $`X_MT_M`$ is the associated metric fibration over the Bass-Serre tree $`T_M`$ of the group $`\mathrm{\Gamma }_M`$, then $`T_M`$ is bushy if and only if $`detM>1`$. In fact, for any graph of finitely generated groups, the Bass-Serre tree is either bounded, quasi-isometric to a line, or bushy, and the question of which alternative holds is easily decided by inspection of the graph of groups.
Here is our generalization of Step 2:
###### Theorem 7.7.
Let $`\pi :XT`$, $`\pi ^{}:X^{}T^{}`$ be metric fibrations over $`\beta `$-bushy trees $`T,T^{}`$, such that the fibers of $`\pi `$ and $`\pi ^{}`$ are contractible $`n`$-manifolds for some $`n`$. Let $`f:XX^{}`$ be a quasi-isometry. Then there exists a constant $`A`$, depending only on the metric fibration data of $`\pi ,\pi ^{}`$, the quasi-isometry data for $`f`$, and the constant $`\beta `$, such that:
1. For each hyperplane $`PX`$ there exists a unique hyperplane $`QX^{}`$ such that $`d_{}(f(P),Q)A`$.
2. For each horizontal leaf $`LX`$ there is a horizontal leaf $`L^{}X^{}`$ such that $`d_{}(f(L),L^{})A`$.
###### Remark.
This result is used in \[MSW\] to prove quasi-isometric ridigity for fundamental groups of geometrically homogeneous graphs of groups whose vertex groups are fundamental groups of manifolds in a category $`𝒞`$ as above, as long as that class of groups is itself quasi-isometrically rigid. For example, quasi-isometric rigidity is proved for graphs of $`𝐙`$’s, $`𝐙^n`$’s, surface groups, lattices in semisimple Lie groups, nilpotent groups, etc.
###### Proof.
To prove (1), by Lemma 7.2 the inclusion map $`PX`$ is uniformly proper and $`P`$ is uniformly contractible, and clearly $`P`$ is a contractible $`(n+1)`$-manifold. Composing with $`f`$ we obtain a uniformly proper map $`PX^{}`$. Now apply Theorem 7.3.
The idea of the proof of (2) is that bushiness of the tree allows one to gain quasi-isometric control over horizontal leaves by considering them as “coarse intersections” of hyperplanes.
###### Definition (Coarse intersection).
A subset $`W`$ of a metric space $`X`$ is a coarse intersection of subsets $`U,VX`$, denoted $`W=U_CV`$, if there exists $`K_0`$ such that for every $`KK_0`$ there exists $`K^{}0`$ so that
$$d_{}(\mathrm{Nbhd}_K(U)\mathrm{Nbhd}_K(V),W)K^{}$$
Note that although such a set $`W`$ may not exist, when it does exist then any two such sets are a bounded Hausdorff distance from each other.
The following fact is an elementary consequence of the definitions.
###### Lemma 7.8.
For any quasi-isometry $`f:XY`$ of metric spaces, and $`U,VX`$, if $`U_CV`$ exists then $`f(U_CV)`$ is a coarse intersection of $`f(U)`$, $`f(V)`$, with constants depending only on the quasi-isometry constants for $`f`$ and the coarse intersection constants for $`U`$ and $`V`$. ∎
Consider now a metric fibration $`\pi :XT`$. A subset of $`X`$ of the form $`X_\sigma =\pi ^1(\sigma )`$, where $`\sigma `$ is an infinite ray in $`T`$, will be called a half-plane in $`X`$. The next lemma is an easy observation—see Figure 1.
###### Lemma 7.9.
Let $`\pi :XT`$ be a metric fibration over a tree $`T`$. Let $`P_1`$ and $`P_2`$ be distinct hyperplanes in $`X`$. Then $`P_1_CP_2`$ exists and is a bounded Hausdorff distance from either a half-plane or a horizontal leaf in $`X`$. Moreover, $`P_1_CP_2`$ is a bounded Hausdorff distance from a half-plane if and only if $`P_1P_2`$ is a half-plane. ∎
We remark that $`P_1_CP_2`$ can be an arbitrarily large finite Hausdorff distance from a horizontal leaf; see Figure 1(b,d).
###### Lemma 7.10.
Let $`\pi :XT`$, $`\pi ^{}:X^{}T^{}`$ be metric fibrations. Let $`f:XX^{}`$ be a quasi-isometry. Suppose $`P_1`$ and $`P_2`$ are distinct hyperplanes in $`X`$ which intersect in a half-plane. Then $`f(P_1)`$ and $`f(P_2)`$ are a uniformly bounded Hausdorff distance from distinct hyperplanes $`Q_1,Q_2`$ in $`X^{}`$ which intersect in a half-plane in $`X^{}`$.
###### Proof.
By Theorem 7.3, there exists a constant $`A`$ so that $`f(P_i)`$ is within Hausdorff distance $`A`$ of a unique hyperplane $`Q_i`$ in $`X^{}`$. Since $`P_1,P_2`$ are distinct they have infinite Hausdorff distance, so $`Q_1`$ and $`Q_2`$ have infinite Hausdorff distance and hence $`Q_1Q_2`$.
By Lemma 7.9, it is enough to prove that $`Q_1_CQ_2`$ is not a bounded Hausdorff distance from a horizontal leaf in $`X^{}`$. If $`Q_1_CQ_2`$ is a bounded Hausdorff distance from a horizontal leaf, then since any horizontal leaf in $`Q_1`$ coarsely separates $`Q_1`$ it must be that $`Q_1_CQ_2`$ coarsely separates $`Q_1`$. But $`P_1_CP_2`$ does not coarsely separate $`P_1`$. This contradicts Lemma 7.4. ∎
We now prove Theorem 7.7. Consider the quasi-isometry $`f:XX^{}`$. Since $`T`$ is bushy, any horizontal leaf $`L`$ in $`X`$ can be realized as a coarse intersection of three hyperplanes $`P_1,P_2,P_3`$, such that the pairwise intersections $`P_1P_2`$, $`P_2P_3`$, $`P_3P_1`$ form three half-planes, any two of which have infinite Hausdorff distance. Moreover, $`d_{}(L,P_1P_2P_3)\beta `$ where $`\beta `$ is a bushiness constant for $`T`$ (see Figure 2).
Consider the unique hyperplane $`Q_i`$ which lies a Hausdorff distance of at most $`A`$ from $`f(P_i)`$, $`i=1,2,3`$. By Lemma 7.10, the pairwise intersections $`Q_1Q_2`$, $`Q_2Q_3`$, $`Q_3Q_1`$ are all half-planes, any two of which have infinite Hausdorff distance. The following elementary fact about trees, applied to $`T^{}`$, now shows that $`Q_1Q_2Q_3`$ is a horizontal leaf $`L^{}`$ in $`X^{}`$:
Fact about trees: Let $`\mathrm{}_1,\mathrm{}_2,\mathrm{}_3`$ be bi-infinite lines in a simplicial tree $`T^{}`$, such that the pairwise intersections $`\mathrm{}_1\mathrm{}_2`$, $`\mathrm{}_2\mathrm{}_3`$, $`\mathrm{}_3\mathrm{}_1`$ are all infinite rays in $`T^{}`$, any two of which have infinite Hausdorff distance. Then $`\mathrm{}_1\mathrm{}_2\mathrm{}_3`$ is a vertex of $`T^{}`$.
Since $`LN_\beta (P_i)`$ it follows that
$$f(L)N_{K\beta +C}(f(P_i))N_{K\beta +C+A}(Q_i),i=1,2,3$$
But clearly we have $`_{i=1}^3N_{K\beta +C+A}(Q_i)=N_{K\beta +C+A}(L^{})`$.
To summarize, given a horizontal leaf $`L`$ of $`X`$, we have found a horizontal leaf $`L^{}`$ of $`X^{}`$ such that $`LN_A^{}(L^{})`$ where $`A^{}=K\beta +C+A`$. A similar argument using a coarse inverse for $`f`$ provides the desired bound for $`d_{}(f(L),L^{})`$. This completes the proofs of Theorem 7.7 and of Step 2. ∎
##### Step 3. A quasi-isometry takes coherent hyperplanes in $`X_M`$ to coherent hyperplanes in $`X_N`$.
Let $`M,N`$ be as in the statement of Proposition 7.1, and fix a quasi-isometry $`f:X_MX_N`$.
Let $`P`$ be any coherent hyperplane in $`X_M`$. By Step 2 it follows that $`f(P)`$ is within a Hausdorff distance $`A`$ from a unique hyperplane $`Q`$ in $`X_N`$. By composing $`f|P`$ with vertical projection $`X_NQ`$ we obtain a map $`\varphi :PQ`$. The inclusion maps $`PX_M`$ and $`QX_N`$ are coarsely lipschitz and uniformly proper; indeed they are isometric embeddings with respect to the induced path metrics on $`P,Q`$. By Lemma 2.1, $`\varphi `$ is a quasi-isometry, with quasi-isometry constants depending only on those for $`f`$. By Step 2, $`f`$ coarsely respects the horizontal foliations of $`X_M`$ and $`X_N`$; vertical projection $`X_NQ`$ takes horizontal leaves to horizontal leaves, and so $`\varphi `$ coarsely respects the horizontal foliations of $`P`$ and $`Q`$, with a coarseness constant depending only on the quasi-isometry constants of $`f`$.
Since $`P`$ is a coherent hyperplane it is isometric to $`G_M`$. Since $`Q`$ is a hyperplane it is isometric to either $`G_N`$ or $`H_N`$, and we now show that the second possibility cannot occur.
###### Proposition 7.11.
Given matrices $`M,N\mathrm{GL}_\times (n,𝐑)`$ with $`detM,detN>1`$, there is no quasi-isometry $`\varphi :G_MH_N`$ which coarsely respects horizontal foliations.
###### Proof.
The idea of the proof is to compare the growth types of the filling area functions for “quasivertical bigons” in $`G_M`$ and in $`H_N`$. In $`G_M`$ this growth type will be quadratic, while in $`H_N`$ it will be exponential.
Let $`H=G_M`$, $`H_M`$, $`G_N`$, or $`H_N`$. There is a quotient map $`H𝐑`$ whose point pre-images give the horizontal foliation of $`H`$, and such that the Hausdorff distance between two horizontal leaves equals the distance between the corresponding points in $`𝐑`$. A path $`\gamma `$ in $`H`$ is said to be $`(K,C)`$-*quasivertical* if its projection to $`𝐑`$ is a $`(K,C)`$-quasigeodesic. Define a $`(K,C)`$-*quasivertical bigon* in $`H`$ to be a pair of $`(K,C)`$-quasivertical paths $`\gamma ,\gamma ^{}`$ which begin and end at the same point.
If $`K,C`$ are fixed, we define a filling area function $`A(L)`$ for $`(K,C)`$-quasivertical bigons in $`H`$. Given a $`(K,C)`$-quasivertical bigon $`\gamma ,\gamma ^{}`$, its *filling area* is the infimal area of a Lipschitz map $`D^2H`$ whose boundary is a reparameterization of the closed curve $`\gamma ^1\gamma ^{}`$; such a map $`D^2H`$ is called a *filling disc* for $`\gamma ^1\gamma ^{}`$. For each $`L0`$ define $`𝒜(L)`$ to be the supremal filling area over all $`(K,C)`$-quasivertical bigons $`\gamma ,\gamma ^{}`$ in $`H`$ such that $`\mathrm{Length}(\gamma )+\mathrm{Length}(\gamma ^{})L`$.
Suppose there is a quasi-isometry $`\varphi :G_MH_N`$ which coarsely respects horizontal foliations. Let $`\overline{\varphi }:H_NG_M`$ be a coarse inverse for $`\varphi `$, also coarsely respecting horizontal foliations. Clearly $`\overline{\varphi }`$ takes any $`K,C`$-quasivertical bigon in $`H_N`$ to a $`(K^{},C^{})`$-quasivertical bigon in $`G_M`$, distorting lengths by at worst an affine function; this affine function, and the constants $`K^{},C^{}`$, depend only on $`K,C`$, the quasi-isometry constants for $`\varphi `$, and the Hausdorff constant for the induced height function. Fill the resulting bigon in $`G_M`$ as efficiently as possible, and map back to $`H_N`$ via $`\varphi `$, distorting area by at worst an affine function which again has the same dependencies. We thereby obtain a filling of the original bigon in $`H_N`$. If $`𝒜_1(L)`$ denotes the filling area function for $`(K^{},C^{})`$-quasivertical bigons in $`G_M`$, and if $`𝒜_2(L)`$ denotes the filling area function for $`(K,C)`$-quasivertical bigons in $`H_N`$, it follows that the growth type of $`𝒜_2(L)`$ is dominated by the growth type of $`𝒜_1(L)`$, that is,
$$𝒜_2(L)\alpha 𝒜_1(\beta L+\delta )+\zeta $$
for some positive constants $`\alpha ,\beta ,\delta ,\zeta `$ independent of $`L`$.
However, we shall now show that $`𝒜_1(L)`$ has a quadratic upper bound while $`𝒜_2(L)`$ has an exponential lower bound, contradicting the above inequality.
Consider a $`K^{},C^{}`$-quasivertical bigon $`\gamma ,\gamma ^{}`$ in $`G_M`$. Applying the argument of Claim 5.7, there are center leaves $`\tau ,\tau ^{}`$ in $`G_M`$ and quasivertical paths $`\rho \tau ,\rho ^{}\tau ^{}`$ which stay uniformly close to $`\gamma ,\gamma ^{}`$, respectively. The initial points of $`\rho ,\rho ^{}`$ are at a uniformly bounded distance, as are the terminal points, and it follows that $`\rho ^{}`$ stays uniformly close to a quasivertical path $`\rho ^{\prime \prime }\tau `$. Connecting initial and terminal endpoints with short paths $`\eta ,\eta ^{}`$ we thus obtain a closed curve $`\rho ^1\eta \rho ^{\prime \prime }\eta ^{}`$, contained in a center leaf of $`G_M`$, which stays uniformly close to $`\gamma ^1\gamma ^{}`$. Since center leaves of $`G_M`$ are isometric to Euclidean space, in which the filling function is quadratic, it follows that $`𝒜_1(L)`$ has a quadratic upper bound.
To show that $`𝒜_2(L)`$ has an exponential lower bound, we now construct quasivertical bigons in $`H_N`$ which can be filled only by discs of exponential area. In the case where $`N`$ is a $`1\times 1`$ matrix such loops are given explicitly in \[ECH<sup>+</sup>92\], Chapter 7.4; examples for general $`N`$ are simple modifications of this example. To be explicit, choose an eigenvalue of $`N`$ of absolute value $`\alpha >1`$; such an eigenvalue exists because $`detN>1`$. Choose an affine subspace $`A𝐑^n`$ parallel to the $`\alpha `$-eigenspace of $`N`$. Consider the subspace $`A\times 𝐑𝐑^n\times 𝐑G_N`$.
For each fixed $`L0`$, choose two vertical segments $`g,g^{}`$ in $`A\times [0,\mathrm{})`$ whose upper endpoints are in $`A\times L`$ and whose lower endpoints are in $`A\times 0`$, and so that the distance in $`A\times L`$ between the upper endpoints, measured using the Riemannian metric on $`G_N`$, is equal to $`1`$; it follows that the distance in $`A\times 0`$ between the lower endpoints, measured using the Riemannian metric on $`G_N`$, is within a constant multiple of $`\alpha ^L`$.
Now double this picture, in the doubled $`G_N`$ horoball $`H_N`$, to get a closed loop in $`H_N`$, that is: in one horoball go up $`g`$, across $`1`$ unit, and down $`g^{}`$, and then in the other horoball go up $`g^{}`$, across $`1`$ unit, and down $`g`$; let $`\rho `$ be the resulting closed curve in $`H_N`$. We have $`\mathrm{Length}(\rho )=4L+2`$. To see that the filling area of $`\rho `$ is exponential in $`L`$, note that any filling disc for $`\rho `$ must contain a path in $`A\times 0`$ connecting the lower endpoints of $`g,g^{}`$, because $`A\times 0`$ separates the two halves of $`\rho `$ in $`H_N`$. This path has length exponential in $`L`$; and a neighborhood of this path in the filling disc has area exponential in $`L`$. ∎
#### Step 4. A horizontal respecting quasi-isometry preserves transverse orientation
Let $`M`$, $`N`$, and $`f:\mathrm{\Gamma }_M\mathrm{\Gamma }_N`$ be as in the statement of Proposition 7.1. By Step 3 there is a quasi-isometry $`\varphi :G_MG_N`$, and by Step 2, $`\varphi `$ coarsely respects the horizontal foliations of $`G_M`$ and $`G_N`$. Suppose that $`\varphi `$ reverses the transverse orientation. There is a quasi-isometry $`G_NG_{N^1}`$ which coarsely respects horizontal foliations, *reversing* transverse orientations. Precomposing with $`\varphi :G_MG_N`$ and applying Steps 1–3, we obtain a quasi-isometry $`G_MG_{N^1}`$ which coarsely respects the transversely oriented horizontal foliations. Applying Theorem 5.2, it follows that $`M`$ and $`N^1`$ have positive real powers with the same absolute Jordan form, and so these powers also have the same determinant. But each positive power of $`M`$ has determinant $`>1`$, whereas every positive power of $`N^1`$ has determinant $`<1`$, a contradiction showing that $`\varphi `$ must preserve the transverse orientation.
This completes the proof of Proposition 7.1.∎
###### Remark.
Note in the proof of Proposition 7.1 that different choices of coherent hyperplanes in $`X_M`$ yield different quasi-isometries $`\varphi `$. In some cases $`\varphi `$ is well-defined up to some constant $`A`$, that is, for any two choices of coherent hyperplane in $`X_M`$, the induced maps $`\varphi _1,\varphi _2:G_MG_N`$ satisfy $`sup_xd(\varphi _1(x),\varphi _2(x))A`$. This is true, for example, in the “centerless” case where $`M,N`$ have no eigenvalues on the unit circle. In the general case, the best that can be said is that the map induced by $`\varphi `$ from the center leaf space of $`G_M`$ to the center leaf space of $`G_N`$ is well defined up to a constant, with respect to the Hausdorff metrics on the center leaf spaces.
## 8 Finding the Integers
In this section we prove Theorem 1.1. Let $`M,N`$ be integral $`(n\times n)`$ matrices with $`\left|detM\right|,\left|detN\right|>1`$. We must prove that $`\mathrm{\Gamma }_M`$ is quasi-isometric to $`\mathrm{\Gamma }_N`$ if and only if there exist positive integers $`a,b`$ such that $`M^a`$ and $`N^b`$ have the same absolute Jordan form.
First we show that the groups $`\mathrm{\Gamma }_{M^a}`$ and $`\mathrm{\Gamma }_M`$ are quasi-isometric, for any positive integer $`a`$, by showing that $`\mathrm{\Gamma }_{M^a}`$ is a subgroup of finite index in $`\mathrm{\Gamma }_M`$, specifically of index $`a`$. To see why, consider the presentations
$`\mathrm{\Gamma }_M`$ $`=𝐙^n,t|t^1xt=M(x),x𝐙^n`$
$`\mathrm{\Gamma }_{M^a}`$ $`=𝐙^n,s|s^1xs=M^a(x),x𝐙^n`$
Define a homomorphism $`\mathrm{\Gamma }_M𝐙/a𝐙`$ by $`𝐙^n0,t1`$. This homomorphism is onto, and its kernel is generated by $`𝐙^n,t^a`$. This kernel is isomorphic to $`\mathrm{\Gamma }_{M^a}`$ under the injection $`\mathrm{\Gamma }_{M^a}\mathrm{\Gamma }_M`$ given by $`xx,st^a`$.
Similarly, $`\mathrm{\Gamma }_{N^b}`$ is quasi-isometric to $`\mathrm{\Gamma }_N`$, for any positive integer $`b`$.
By squaring $`M,N`$ if necessary, we may therefore assume that $`detM,detN>1`$, and that $`M`$ and $`N`$ lie on 1-parameter subgroups; we continue with this assumption up through the end of the proof in §8.2. Choose 1-parameter subgroups $`M^t,N^t`$ of $`\mathrm{GL}(n,𝐑)`$ with $`M=M^1,N=N^1`$, let $`G_M,G_N`$ be the associated Lie groups constructed in §4, and let $`X_M,X_N`$ be the associated geodesic metric spaces constructed in §7. The group $`\mathrm{\Gamma }_M`$ is quasi-isometric to $`X_M`$, and $`\mathrm{\Gamma }_N`$ is quasi-isometric to $`X_N`$.
### 8.1 The first half of the classification
Assuming that $`M^a`$ and $`N^b`$ have the same absolute Jordan form, where $`a,b`$ are positive integers, we must prove that $`\mathrm{\Gamma }_M`$ and $`\mathrm{\Gamma }_N`$ are quasi-isometric. We have shown above that $`\mathrm{\Gamma }_{M^a}`$ and $`\mathrm{\Gamma }_M`$ are quasi-isometric, and that $`\mathrm{\Gamma }_{N^b}`$ and $`\mathrm{\Gamma }_N`$ are quasi-isometric. Replacing $`M`$ by $`M^a`$ and $`N`$ by $`N^b`$, we may therefore assume that $`M,N`$ have the same absolute Jordan form. We shall prove that $`\mathrm{\Gamma }_M,\mathrm{\Gamma }_N`$ are quasi-isometric by constructing a bilipschitz homeomorphism between $`X_M`$ and $`X_N`$.
Since the absolute Jordan forms of $`M,N`$ are equal it follows that $`detM=detN`$; let $`d`$ be the common value. Applying Proposition 4.1, there is a bilipschitz homeomorphism from $`G_M=𝐑^n_M𝐑`$ to $`G_N=𝐑^n_M𝐑`$ of the form $`(x,t)(Ax,t)`$ for some $`A\mathrm{GL}(n,𝐑)`$. In the fiber product description of $`X_M`$, $`X_N`$, the trees $`T_M`$ and $`T_N`$ may both be identified with the homogeneous, oriented tree $`T_d`$ with one incoming and $`d`$ outgoing edges at each vertex. The bilipschitz homeomorphism $`G_MG_N`$ and the identity homeomorphism $`T_dT_d`$ both respect the height functions, and so these two homeomorphisms combine to give the desired bilipschitz homeomorphism $`X_MX_N`$.
### 8.2 Quasi-isometric implies integral powers have same absolute Jordan forms
Assuming $`\mathrm{\Gamma }_M,\mathrm{\Gamma }_N`$ are quasi-isometric, there is a quasi-isometry $`f:X_MX_N`$. Combining Proposition 7.1 and Theorem 5.2 gives $`r𝐑_+`$ such that $`M^r`$ and $`N`$ have the same absolute Jordan form. We must show that there exist $`a,b𝐙_+`$ so that $`M^a`$ and $`N^b`$ have the same absolute Jordan form.
Since $`M^r`$ and $`N`$ have the same absolute Jordan form, listing the absolute values of the eigenvalues of $`M`$ and $`N`$ in increasing order we obtain
$`\mu _a\mathrm{}\mu _01<\mu _1=:\alpha _M\mathrm{}\mu _b`$
$`\nu _a\mathrm{}\nu _01<\nu _1=:\alpha _N\mathrm{}\nu _b`$
with $`\mu _i^r=\nu _i`$, $`aib`$. From this it follows that
$$\frac{\mathrm{log}\alpha _N}{\mathrm{log}\alpha _M}=r=\frac{\mathrm{log}detN}{\mathrm{log}detM}$$
Let $`𝐐_M`$ denote the set of coherent hyperplanes in $`X_M`$, and let $`h_M`$ denote the height function on $`M`$. We define a metric on $`𝐐_M`$ as follows: given coherent hyperplanes $`P_1,P_2`$, let $`L`$ denote the horizontal leaf $`L=(P_1P_2)`$. Then we set
$$d_{𝐐_M}(P_1,P_2)=\left(detM\right)^{h_M(L)}$$
It is easy to check that this defines a metric on $`𝐐_M`$, and since the tree $`T_M`$ branches $`m=detM`$ times as $`h_M`$ increases by $`1`$, the metric space $`(𝐐_M,d_{𝐐_M})`$ is isometric to the $`m`$-adic rational numbers in their usual metric of Hausdorff dimension 1. Similarly, attached to $`X_N`$ is a metric space $`(𝐐_N,d_{𝐐_N})`$ isometric to the $`n`$-adic rational numbers, with $`n=detN`$.
From Step 3 in the proof of Proposition 7.1 (see §7.2), the quasi-isometry $`f:X_MX_N`$ takes each coherent hyperplane in $`X_M`$ to within a uniform Hausdorff distance of a unique coherent hyperplane in $`X_N`$, hence induces a bijection $`\psi :𝐐_M𝐐_N`$. For each $`\mathrm{}𝐐_M`$, setting $`\mathrm{}^{}=\psi (\mathrm{})`$, there is an induced horizontal-respecting quasi-isometry $`P_{\mathrm{}}P_{\mathrm{}^{}}^{}`$, and by Time Rigidity (Proposition 5.8) this quasi-isometry has an induced time change of the form $`tmt+b`$ where
$$m=\frac{\mathrm{log}\alpha _M}{\mathrm{log}\alpha _N}=1/r$$
and where $`b`$ depends ostensibly on $`\mathrm{}`$. However, for another $`\mathrm{}_1`$, $`P_{\mathrm{}}`$ and $`P_\mathrm{}_1`$ coincide below some value of $`t`$, and so $`tmt+b`$ is an induced time change for both $`P_{\mathrm{}}P_{\mathrm{}^{}}^{}`$ and $`P_\mathrm{}_1P_\mathrm{}_1^{}^{}`$, possibly with a larger coarseness constant (this argument is taken from Claim 6.3 on p. 436 of \[FM98\]). Therefore, there is a uniform induced time change $`tmt+b`$ with $`b`$ independent of $`\mathrm{}`$, and with a uniform Hausdorff constant $`A`$.
We now claim that $`\psi `$ is a bilipschitz homeomorphism. To this end, let $`P_1,P_2𝐐_M`$ be given. Let $`L=(P_1P_2)`$ and let $`L^{}=(\psi (P_1)\psi (P_2))`$. Then
$$h_N(L^{})mh_M(L)+bA$$
Hence
$`{\displaystyle \frac{d_{𝐐_N}(\psi (P_1),\psi (P_2))}{d_{𝐐_M}(P_1,P_2)}}`$ $`={\displaystyle \frac{(detN)^{h_N(L^{})}}{(detM)^{h_M(L)}}}`$
$`{\displaystyle \frac{(detN)^{mh_M(L)b+A}}{(detM)^{h_M(L)}}}`$
$`={\displaystyle \frac{\left((detN)^{(\mathrm{log}detM/\mathrm{log}detN))}\right)^{h_M(L)}(detN)^{b+A}}{(detM)^{h_M(L)}}}`$
$`=(detN)^{b+A}`$
which is a constant not depending on $`P_1`$ or $`P_2`$. Hence $`\psi `$ is Lipschitz. The same argument applied to $`\psi ^1`$ shows that $`\psi `$ is bilipschitz.
Applying Cooper’s Theorem (appendix to \[FM98\], Corollary 10.11) on bilipschitz homeomorphisms of Cantor sets, we obtain that there exist integers $`a,b>0`$ such that $`(detM)^a=(detN)^b`$. Since $`M^r`$ and $`N`$ have the same absolute Jordan form, we have
$$\frac{b}{a}=\frac{\mathrm{log}detM}{\mathrm{log}detN}=r$$
and so $`(M^r)^a=M^b`$ and $`N^a`$ have the same absolute Jordan form.
## 9 Quasi-isometric rigidity
In this section we prove Theorem 1.2 in a series of steps. Recall the hypotheses: $`M`$ is an integer matrix in $`\mathrm{GL}(n,𝐑)`$ with $`\left|detM\right|>1`$, and $`G`$ is a finitely generated group quasi-isometric to $`\mathrm{\Gamma }_M`$. By squaring $`M`$ if necessary we may assume that $`M\mathrm{GL}_\times (n,𝐑)`$ and $`detM>1`$, and therefore $`\mathrm{\Gamma }_M`$ is quasi-isometric to $`X_M`$. It follows that $`G`$ is quasi-isometric to $`X_M`$.
##### Step 1.
The action of $`G`$ on itself by left multiplication can be conjugated by the quasi-isometry $`GX_M`$ to give a proper, cobounded quasi-action of $`G`$ on $`X_M`$ (see \[FM99b\], Proposition 2.1). Since $`detM>1`$ we may apply Theorem 7.7, concluding that the quasi-action of $`G`$ on $`X_M`$ coarsely respects the fibers of the uniform metric fibration $`X_MT_M`$.
##### Step 2.
Now we use the following result of \[MSW\]. Suppose $`\pi :XT`$ is a uniform metric fibration over a bushy tree $`T`$. If $`G`$ is a finitely presented group with a cobounded, proper quasi-action on $`X`$, and if the quasi-action coarsely respects the fibers, then $`G`$ is the fundamental group of a graph of groups whose vertex and edge groups are quasi-isometric to a fiber $`X_t=\pi ^1(t)`$.
By Step 1, this result applies to the quasi-action of $`G`$ on $`X_M`$, because $`G`$ is quasi-isometric to the finitely presented group $`\mathrm{\Gamma }_M`$ and so $`G`$ is finitely presented. The fibers of the map $`X_MT_M`$ are isometric to $`𝐑^n`$, and it follows that $`G`$ is the fundamental group of a graph of groups with each vertex and edge group quasi-isometric to $`𝐑^n`$.
##### Step 3.
Any finitely-generated group quasi-isometric to $`𝐑^n`$ is virtually $`𝐙^n`$ (see \[Ger95\]), and so $`G`$ is the fundamental group of a graph of groups whose vertex and edge groups are virtually $`𝐙^n`$.
##### Step 4.
Applying the argument in Section 5 of \[FM99b\] to $`G`$ gives that either $`G`$ contains a noncyclic free group or $`G`$ is an ascending HNN extension of the form
$$G=A_\varphi =A,t|tat^1=\varphi (a),aA$$
where $`A`$ is virtually $`𝐙^n`$ and $`\varphi :AA`$ is an injective endomorphism. Since $`\mathrm{\Gamma }_M`$ is amenable, and since $`G`$ is quasi-isometric to $`\mathrm{\Gamma }_M`$, then $`G`$ is amenable, and so $`G`$ cannot contain a noncyclic free group. The second possibility must therefore occur: $`G=A_\varphi `$ as above.
##### Step 5.
Now we turn to an analysis of injective endomorphisms of virtually abelian groups. Suppose $`A`$ is a finitely generated, virtually abelian group. Any injective endomorphism of $`A`$ has finite index image.
A subgroup $`BA`$ is *characteristic for endomorphisms* if, for any injective endomorphism $`\varphi :AA`$, we have $`\varphi (B)B`$.
Given a group $`A`$ and $`gA`$, the centralizer of $`g`$ in $`A`$ is denoted $`C_A(g)`$. The *virtual center* of $`A`$, denoted $`V(A)`$, is the set of all $`gA`$ such that $`[A:C_A(g)]<\mathrm{}`$. This is a subgroup, because if $`g,hV(A)`$ then the subgroup $`C_A(gh)`$, which contains $`C_A(g)C_A(h)`$, has finite index.
###### Lemma 9.1 (Some characteristic subgroups).
Let $`A`$ be a finitely generated, virtually abelian group. Then the virtual center $`V(A)`$, its center $`ZV(A)`$, and its torsion subgroup $`TZV(A)`$, are all characteristic for endomorphisms of $`A`$. Moreover, $`V(A)`$ and $`ZV(A)`$ both have finite index in $`A`$, whereas $`TZV(A)`$ is finite.
Lemma 9.1 is proved below.
##### Step 6.
Consider the HNN extension $`G=A_\varphi `$ above. Let $`V(A)`$, $`ZV(A)`$, $`TZV(A)`$ be as in Lemma 9.1, so all these subgroups are taken into themselves by $`\varphi `$. Since $`TZV(A)`$ is finite we in fact have $`\varphi (TZV(A))=TZV(A)`$, and so $`K=TZV(A)`$ is a finite, normal subgroup of $`G`$.
Replacing $`G`$ by $`G/K`$, we may assume that $`TZV(A)`$ is trivial, and it follows that $`ZV(A)`$ is torsion-free abelian, and so is isomorphic to $`𝐙^n`$. Since $`\varphi (ZV(A))ZV(A)`$, the action of $`\varphi `$ on $`ZV(A)`$ is given by some $`n\times n`$ matrix of integers $`N`$. Thus, $`G/K`$ has a finite-index subgroup isomorphic to $`\mathrm{\Gamma }_N`$, finishing the proof of Theorem 1.2.
##### Proof of Lemma 9.1.
To see $`[A:V(A)]<\mathrm{}`$, if $`B`$ is any finite-index abelian subgroup of $`A`$ then obviously $`BV(A)`$.
Consider an endomorphism $`\varphi :AA`$. We now show that $`\varphi (V(A))V(A)`$. Consider $`gV(A)`$, so $`[A:C_A(g)]<\mathrm{}`$. It follows that $`[\varphi (A):C_{\varphi (A)}(\varphi (g))]<\mathrm{}`$, and so $`[A:C_{\varphi (A)}(\varphi (g))]<\mathrm{}`$. But $`C_{\varphi (A)}(\varphi (g))C_A(\varphi (g))`$, and so $`\varphi (g)V(A)`$.
Next we claim that $`V(V(A))=V(A)`$. To see why, if $`gV(A)`$ then $`[A:C_G(g)]<\mathrm{}`$, and so $`[V(A):C_G(g)V(A)]<\mathrm{}`$. But $`C_G(g)V(A)C_{V(A)}(g)`$, and so $`[V(A):C_{V(A)}(g)]<\mathrm{}`$, i.e. $`gV(V(A))`$.
Next we claim that $`[V(A):ZV(A)]<\mathrm{}`$. In fact if $`V`$ is any finitely generated group which is its own virtual center, then $`[V:ZV]<\mathrm{}`$ (the converse is also true, trivially). To see why, let $`g_1,\mathrm{},g_k`$ be a generating set for $`V`$. Since $`V(V)=V`$, each of the groups $`C_V(g_1),\mathrm{},C_V(g_k)`$ has finite index in $`V`$. It follows that their intersection has finite index in $`V`$; but their intersection is precisely $`ZV`$.
Now we claim that $`ZV(A)`$ is characteristic for endomorphisms of $`V(A)`$ (and so is also characteristic for endomorphisms of $`A`$). In fact, if $`V`$ is any finitely generated group whose center $`ZV`$ has finite index, then $`ZV`$ is characteristic for any injective endomorphism $`\varphi :VV`$ whose image has finite index. To see why, we have $`Z(\varphi (V))=\varphi (ZV)`$, and so
$$[\varphi (V):Z(\varphi (V))]=[\varphi (V):\varphi (ZV)]=[V:ZV]<\mathrm{}$$
Clearly $`\varphi (V)ZVZ(\varphi (V))`$, and so
$$[\varphi (V):Z(\varphi (V))][\varphi (V):\varphi (V)ZV]$$
The quotient group $`V/ZV`$ is finite, and the quotient homomorphism $`VV/ZV`$, when restricted to the subgroup $`\varphi (V)`$, has kernel $`\varphi (V)ZV`$. It follows that
$$[\varphi (V):\varphi (V)ZV]|V/ZV|=[V:ZV]=[\varphi (V):Z(\varphi (V))]$$
All of the above inequalities are therefore equalities, and so
$$\varphi (ZV)=Z(\varphi (V))=\varphi (V)ZV$$
which implies $`\varphi (ZV)ZV`$.
Finally, it is clear that for any finitely generated abelian group, the torsion subgroup is characteristic for injective endomorphisms.
## 10 Questions
### 10.1 Remarks on the polycyclic case
Given an integer matrix $`M\mathrm{GL}(n,𝐑)`$, the group $`\mathrm{\Gamma }_M`$ is polycyclic if and only if $`\left|detM\right|=1`$, and if $`M\mathrm{GL}_\times (n,𝐑)`$ this occurs if and only if $`\mathrm{\Gamma }_M`$ is a cocompact discrete subgroup of $`G_M`$. In this case it follows that $`\mathrm{\Gamma }_M`$ is quasi-isometric to $`G_M`$, and the notion of horizontal-respecting quasi-isometry clearly transfers to $`\mathrm{\Gamma }_M`$. The techniques of this paper do not provide a quasi-isometric classification in this case, however they do yield the following partial result:
###### Theorem 10.1.
If $`M,N\mathrm{SL}(n,𝐙)`$ lie on 1-parameter subgroups of $`\mathrm{GL}(n,𝐑)`$, then there is a horizontal respecting quasi-isometry $`\mathrm{\Gamma }_M\mathrm{\Gamma }_N`$ if and only if there is a horizontal respecting quasi-isometry $`G_MG_N`$, and this occurs if and only if there are real numbers $`a,b0`$ such that $`M^a,N^b`$ have the same absolute Jordan form. ∎
This raises the question: Is every quasi-isometry $`\mathrm{\Gamma }_M\mathrm{\Gamma }_N`$ horizontal respecting? Equivalently, is every quasi-isometry $`G_MG_N`$ horizontal respecting? The answer is obviously no, for example when $`M,N`$ are identity matrices and $`G_M,G_N`$ are Euclidean spaces. However, we conjecture:
###### Conjecture 10.2.
If $`M,N\mathrm{SL}(n,𝐙)`$ lie on 1-parameter subgroups of $`\mathrm{GL}(n,𝐑)`$, and if $`M,N`$ have no eigenvalues on the unit circle, then any quasi-isometry $`G_MG_N`$ is horizontal respecting.
Moreover, Theorem 10.1 and Conjecture 10.2 together would imply the following (see \[FM99c\]):
###### Conjecture 10.3.
Suppose $`M\mathrm{SL}(n,𝐙)`$ has no eigenvalues on the unit circle. If $`G`$ is any finitely generated group quasi-isometric to $`\mathrm{\Gamma }_M`$, then there is a finite normal subgroup $`F`$ of $`G`$ so that $`G/F`$ is abstractly commensurable to $`\mathrm{\Gamma }_N`$, for some $`N\mathrm{SL}(n,𝐙)`$ with no eigenvalues on the unit circle.
### 10.2 The quasi-isometry group of $`\mathrm{\Gamma }_M`$
Given a finitely generated group $`G`$, the set of quasi-isometries from $`G`$ to itself, modulo the identification of quasi-isometries which differ by a bounded amount, forms a group called the *quasi-isometry group* of $`G`$, denoted $`\mathrm{QI}(G)`$. Given a $`1\times 1`$ matrix $`M=(m)`$ with $`m2`$, the quasi-isometry group of the solvable Baumslag-Solitar group $`\mathrm{\Gamma }_M\mathrm{BS}(1,m)`$ was computed in \[FM98\]:
$$\mathrm{QI}(\mathrm{BS}(1,m))\mathrm{Bilip}(𝐑)\times \mathrm{Bilip}(𝐐_m)$$
where $`𝐐_m`$ is the metric space of $`m`$-adic rational numbers, and $`\mathrm{Bilip}(X)`$ denotes the group of bilipschitz self maps of a metric space $`X`$.
###### Problem 10.4.
Compute the quasi-isometry group of $`\mathrm{\Gamma }_M`$ in general.
The strongest result we have on this problem so far is Proposition 6.3, but see the remarks after that proposition.
In \[FM99b\] the computation of $`\mathrm{QI}(\mathrm{BS}(1,m))`$ was applied to prove quasi-isometric rigidity of $`\mathrm{BS}(1,m)`$, using techniques of Hinkkanen \[Hin85\] and Tukia \[Tuk86\]. While quasi-isometric rigidity of $`\mathrm{BS}(1,m)`$ now has a completely different proof \[MSW\], which we have here generalized to $`\mathrm{\Gamma }_M`$, one might still pursue:
###### Problem 10.5.
Give a proof of quasi-isometric rigidity of $`\mathrm{\Gamma }_M`$, generalizing the results of \[FM99b\].
This should lead to a deeper understanding of the geometry of $`\mathrm{\Gamma }_M`$. For example, Tukia \[Tuk86\] characterizes subgroups of the quasiconformal group of a sphere which are conjugate into the Mobiüs group. We have analogous results for lattices in three-dimensional solv-geometry, and there should be generalizations to solvable Baumslag-Solitar groups and to $`\mathrm{\Gamma }_M`$.
Benson Farb:
Department of Mathematics
University of Chicago
5734 University Ave.
Chicago, Il 60637
farb@math.uchicago.edu
Lee Mosher:
Department of Mathematics and Computer Science
Rutgers University, Newark
Newark, NJ 07102
mosher@andromeda.rutgers.edu
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# New type astrophysical solution to the solar neutrino problems and its predictions to the SNO
## 1 Introduction
Two different approaches seem to be promising in exploring the roots of the solar neutrino problems. One is the traditional particle physics approach, suggesting that neutrino oscillations are responsible for the missing solar neutrinos. Unfortunately, the small mixing angle, large mixing angle and vacuum oscillation solutions all suffer from being poor fits to the observations. Eleven years ago Bahcall (1989) wrote in his book, that the MSW neutrino oscillation solution of the solar neutrino problem (SNP) is attractive since the mixing angles and mass differences can each vary by orders of magnitude, and so it does not need fine tuning as the vacuum solution does. But not so much later Paterno and Scalia (1994) noticed that the allowed region shrinkened to a point-like area, and so ”hypotheses on non-conventional neutrino properties are strongly disfavoured, except for the matter neutrino oscillations, the latter surviving within very narrow limits”. Now even this remained point-like area (and the surrounding small 95 $`\%`$ C.L. region, which extends only due to the theoretical and experimental errors) do not have a high allowance. The probability of these solutions decreased from 100 $`\%`$ to cca. 10$`\%`$ for the widely regarded best fit of small mixing angle (SMA) solution (and similarly for the LMA and VAC solutions), when the SuperKamiokande (SK) rates, spectral and day-night effect data are also taken into account.
One can evaluate the acceptability of the presented solutions by the probability $`P(\chi ^2)`$ belonging to the calculated $`\chi _{min}^2`$. The rule of thumb telling that ”the value of $`\chi ^2`$ for a ”moderately” good fit is $`\chi ^2d.o.f.`$” (Press et al., 1992) may be useful. Bahcall, Krastev and Smirnov (1998) obtained for the SMA global fit $`\chi _{min}^2=26.5/17d.o.f.`$ which is acceptable at the probability $`P(\chi ^2)`$ = 7 $`\%`$ C.L. Suzuki (1998) obtained $`\chi _{min}^2=50.2/31d.o.f.`$ for the case when oscillations are not allowed. The corresponding probability is $`P(\chi ^2)=1.6\%`$. The best fit he obtained is for the vacuum oscillation (VAC), with $`\chi _{min}^2=38.7/31`$ d.o.f., with $`P(\chi ^2)20\%`$. This latter case we recognize also as a solution worse than moderately good, since $`\chi _{min}^2>d.o.f.`$ At the same time, Bahcall, Krastev and Smirnov (1998) obtained a global fit (when the constraints from the rates, the spectrum shape and the Day-Night asymmetry are all included) to for the best fit vacuum oscillation solution $`\chi _{min}^2=28.4`$ for 18 d.o.f., which is acceptable only at 6$`\%`$ C.L.
It may be disturbing that two different and contrasting confidences are introduced to characterise the acceptance of a fit. One is $`P(\chi ^2)`$, the probability of a fit. The other one, $`P(\mathrm{\Delta }\chi ^2)`$, characterises how far the (allegedly normal) distribution of the experimental and theoretical errors extends. Since the larger errors leads to larger allowed regions, and the farther iso-probability contours belong to increasing $`n`$ characterising the distance as $`n\sigma `$, therefore the higher is this second probability, the worse is the fit. Therefore, the higher is the former probability $`P(\chi ^2)`$, the better is the fit; while, on the contrary, the higher is the latter probability, $`P(\mathrm{\Delta }\chi ^2)`$, the worse is the fit.
In Gonzalez-Garcia and Pena-Garay (2000) Table 3, the larger C.L. belongs to a higher $`\chi ^2/d.o.f`$ and so it shows a poorer fit. They found for the SMA a C.L. $`83\%`$, which tells us that it is somewhere outside the 1.5 $`\sigma `$ region around the best fit which has a low (cca. 15 $`\%`$) probability itself. Langacker (1999) noted that the LMA solution and the no oscillation hypothesis has the same rate of goodness, i.e. both of them is disfavored at the 95-99 $`\%`$ CL level, while the LMA solution is also a very poor fit, although it is allowed at 95$`\%`$ CL. In other terms, this means that SMA is refuted with more than 1.5 $`\sigma `$.
Recently Maris and Petcov (2000) found that ”the conservative SMA around the point $`\mathrm{\Delta }m^2=4\times 10^6`$, $`sin^22\mathrm{\Theta }=0.0085`$ is ruled out at 1.5 sigma” by the degree of the day-night effect observed by the SK. Now Haxton (2000) noted that ”One puzzling aspect of atmospheric, solar, and LSND neutrino results is that they require three independent $`\delta m^2`$s. That is, they do not respect the relation $`\delta m_{21}^2+\delta m_{31}^2+\delta m_{32}^2=0`$, thus either one or more of the neutrino experiments must be attributed to some phenomena other than neutrino oscillations, or a fourth neutrino is required.”
We learned that the neutrino oscillation solutions of the SNP actually are poor fits. Moreover, we recognise that even if they would be moderately good fits, it would be still an urgent need to find new physics outside from the neutrino oscillations to explain the results of the neutrino experiments. In this Letter we suggest to look for the new physics in a conservative region, in the field of astrophysics. The reason is the (not yet recognised) significance of the facts corresponding to the coupling of the solar core to the surface activity phenomena (one can find the presentation of the astrophysical core-related problems in Grandpierre, 1996, 1999, 2000). It is known, that the solar core rotates so slowly that its rotation rate is the surface rate $`\pm 30\%`$ (Elsworth et al., 1995). At the same time, the so-called ’best solar model’ of Pinsonneault et al. (1989) predicted a rotation rate of the core which is 4-15 times the surface rate. The discrepancy may be surmounted only when allowing a coupling between the solar core and the surface active regions. But if the core may participate in the surface activity cycle, the energy production has to be touched, and so the neutrino production may be influenced as well. In this way we found a way how the new physics arises in the field of astrophysics.
The objections raised against a possible astrophysical solution to the solar neutrino problems (see e.g. Bludman, Hata and Langacker 1994) are valid only for non-standard models without new physics (without changing the spectral shape of the individual neutrino fluxes). But it is known that the solution of the SNP needs new physics (Bludman, Hata and Langacker, 1994). The often used term ”model-independence” refers to general models in which the individual pp, $`Be^7`$, CNO and $`B^8`$ neutrino fluxes are allowed to vary as free parameters. But if new physics is present, another kind of nuclear reactions may also contribute to the energy and neutrino production of the Sun. In this paper we attempt to show how the presence of a high-temperature energy source (Grandpierre, 1996, 2000) could contribute to the neutrino detector data.
It is a general view that the principal neutrino sources are pp, pep, $`{}_{}{}^{7}Be`$, $`{}_{}{}^{8}B`$, $`{}_{}{}^{13}N`$, and $`{}_{}{}^{15}O`$. As Bahcall and Krastev (1996) remarked, this picture has become so accepted that it is sometimes referred to as ”model-independent”. Nevertheless, if other neutrino sources do exist, then these previously thought ”model-independent” models may all fail at a certain rate. Therefore, the related luminosity constraint also may prove insufficient and over-constrained when compared to the actual Sun, since the hot bubbles may produce a smaller or larger part of the solar luminosity. The luminosity constraint expresses the fact that the energy productions is related to the neutrino production since both is produced by the nuclear reactions:
$`L_{}/4\pi R^2={\displaystyle \underset{\alpha }{}}(Q/2<E>_\alpha )\varphi (\alpha ),`$ (1)
where R=1 A. U. (1.469$`\times 10^{13}`$ cm), $`<E>_\alpha `$ is the average neutrino energy and $`\varphi (\alpha )`$ is the $`\alpha `$th flux ($`\alpha `$=pp, pep, $`{}_{}{}^{7}Be,^8B,`$…), and $`Q`$ is the energy released in the fusion reaction $`4p+2e\alpha +2\nu `$. The most general luminosity constraint we found in the literature is presented in Minakata and Nunokawa, 1999 in the following form:
$`L_{}/4\pi R^2`$ $`=13.1\varphi (pp)+11.92\varphi (pep)+12.5\varphi (^7Be)+`$ (2)
$`6.66\varphi (^8B)+3.46\varphi (^{13}N)+21.57\varphi (^{15}O)`$
$`+2.36\varphi (^{17}F)+10.17\varphi (hep).`$
We found it more suitable to convert these neutrino fluxes to fractional fluxes $`\mathrm{\Phi }`$, normalised to the BP98 SSM fluxes:
$`1`$ $`=0.9119\mathrm{\Phi }(pp)+0.001966\mathrm{\Phi }(pep)+0.0758\mathrm{\Phi }(Be)`$ (3)
$`+5.193\times 10^5\mathrm{\Phi }(B)+3\times 10^3\mathrm{\Phi }(CNO).`$
Neglecting the small terms, and including the pep into the pp, the CNO into the Be term,
$`0.99=0.914\mathrm{\Phi }_1+0.076\mathrm{\Phi }_7,`$ (4)
where $`\mathrm{\Phi }_1`$ refers to the pp+pep fluxes, and $`\mathrm{\Phi }_7`$ to the Be+CNO fluxes (where we modified the left side with the condition that for $`\mathrm{\Phi }_1`$=$`\mathrm{\Phi }_7`$=1 the equation must be valid).
We used the fractional neutrino fluxes following Minakata and Nunokawa (1998). The chlorine equation is:
$`2.56=1.8\mathrm{\Phi }_7+5.9\mathrm{\Phi }_8.`$ (5)
The gallium-equation is:
$`72.4=69.6\mathrm{\Phi }_1+46.9\mathrm{\Phi }_7+12.4\mathrm{\Phi }_8.`$ (6)
It is easy to derive from these equations the lines of the Fig. 1 of Fiorentini and Ricci (1998), shown here as Fig. 1. It is confirmed that the non standard cool Sun models do not give an acceptable fit, as they are farther than $`3\sigma `$ from the intersection of the acceptance zones.
Now we are prepared to consider the case when hot bubbles are present. If the energy production of the hot bubbles is not negligible when compared to the total solar luminosity, we should involve it into the luminosity constraint. Regarding the energy budget of the hot bubbles, it depends on their temperature, which is indicated to be in the range of $`10^8`$ to $`10^{11}`$ K (Grandpierre, 1996, 2000). One can think that the bubbles may produce energy through the hot CNO cycle, triple alpha cycle (this cycle does not produce neutrinos) and other nova-type nuclear reactions (Audouze, Truran, Zimmerman, 1973). From these reactions the one which may produce the largest number of events in the neutrino detectors may be the hot CNO cycle. Writing the luminosity constraint for the hot bubble separately, for the case when the CNO cycle gives ten percent of the total solar luminosity as an upper limit, from (1) we obtain:
$`0.1L_{}=12.525\mathrm{\Phi }^b(CNO).`$ (7)
From this constraint, we can derive an upper limit for the bubble CNO neutrino flux,
$`\mathrm{\Phi }^b(CNO)<6\times 10^9cm^2s^1,`$ (8)
which is close to the value of the SSM beryllium-neutrino flux. This value is compatible with the present-day global constraints $`0.0<\mathrm{\Phi }(^7Be)<6.35`$ (Bahcall and Krastev, 1996).
Now we can calculate how Fig. 1 is modified when the hot bubbles are producing 22$`\%`$ of the total solar luminosity. In this case the luminosity constraint for the quiet solar core without the bubbles should be formulated as:
$`0.77=0.914\mathrm{\Phi }_1+0.076\mathrm{\Phi }_7.`$ (9)
Using this constraint in the Ga-equation (6), we observe that the effect of the hot bubble energy generation, or, more precisely, the constraint that the quiet solar core should produce a less than total solar luminosity, is to shift the Ga-zone upwards. This effect is helpful in obtaining better and physical (i.e. $`\mathrm{\Phi }(Be)>0`$) fluxes, and so to resolve one of the solar neutrino problems, the problem of beryllium flux.
Now one may consider what happens with the Kamiokande zone if hot bubbles are present. Since the bubbles may reach very high temperatures, they may produce mu and tau neutrinos, which can be observed by the Kamiokande and not by the other detectors. Therefore, since Fig. 2 describes the quiet solar core only, the Kam-zone may be shifted to the left. Actually, the rate of the necessary shift is determined by the intersection of the Cl and Ga zones. We observe from Fig. 2 that the shift occurs around $`\mathrm{\Phi }_8`$=0.26, which means that the bubbles has to produce a $`\nu _{\mu ,\tau }`$ neutral current with a contribution to the SK
$`\mathrm{\Phi }_{\mu ,\tau }^b=0.21.`$ (10)
We note here, that since the hot bubbles may contribute also to the high-energy excess observed at the SK, therefore the amount of $`\mathrm{\Phi }_{\mu ,\tau }^b`$ should not be so large. The Cl-zone of Fig. 1 will not be modified when going to Fig. 2.
## 2 Discussion and Conclusions
It is interesting that the dynamic solar model (DSM) obtained with the inclusion of the hot bubbles into the standard solar model (SSM) modifies the Ga, Cl, Kam and non-standard zones in a way the create shifting an overlap region in a physical range ($`\mathrm{\Phi }_7>0`$. Remarkably, the cool Sun model also overlaps with this overlap region. For a bubble-luminosity around 0.22 (of the total solar luminosity), they all fit with cca. 1.2 $`\sigma `$. Fig. 2 shows that the three kinds of neutrino detectors actually do not contradict to the standard neutrino picture. The situation is that the combined neutrino results actually indicate a lower than standard central temperature of the Sun with a remarkable confidence. Fig. 2 gives a better fit (cca. 1.2 $`\sigma `$) than that of the presently favoured SMA solution (higher than 1.5$`\sigma `$). The best fit may be reached when both the bubbles and the oscillations are taken into account.
One can observe from Fig. 2 that the new astrophysics may offer powerful perspectives to influence the solution of the SNP. It seems to be not true, that we know the basic physics of the Sun enough. Even if one can consider the presence of the hot bubbles in the solar core as yet not established, the frequent statements that astrophysical solutions are ruled out, prove to be unfounded in the here presented more general basis.
In Fig. 2 the non standard low-temperature solutions are also presented. Although in Fig. 1 they are found far from the 1 $`\sigma `$ zones, in the case of Fig. 2 they are around 1.2 $`\sigma `$ from the overlap region when $`T_c=0.942`$. This value of central temperature is quite consistent with the 0.78 solar luminosity produced by the quiet solar core.
We do not attempt to suggest that the hot bubbles permanently produce a significant part of the solar luminosity. Our attempt is more narrowly confined: it is to point out that the bubbles in principle may create a situation in which the astrophysical solution may be alone enough to solve the problem of the missing solar neutrinos. It is quite plausible, that the electron neutrinos do oscillate, and this phenomenon is responsible for a large part of the SNP and other apparent ”neutrino anomalies”. Nevertheless, the point is that if we ignore the role played by the dynamics of the solar core in the production of solar neutrinos, we may found ourselves in the uncomfortable situation of poor fits of the MS parameters and inconsistent consequences of atmospheric, LSND and solar neutrino problems.
To resolve the astrophysical problems of the solar core it is not necessary that the hot bubbles produce a significant part of the solar luminosity. Actually, they may produce only a negligible part of the solar luminosity and they could prove still be able to trigger and influence the surface solar activity. Gorbatzky (1964) already calculated that hot bubbles arising from point explosions with an initial surplus energy around $`10^{35}`$ ergs may be able to reach the stellar surfaces from 0.1 stellar radius. He ignored completely the energy production of these hot bubbles with temperatures $`T>10^8`$ K. But even a slight influence of the solar core dynamics might be able to improve significantly the fit and consistency of the neutrino problems.
The result red from Fig. 2 has predictive value for the future neutrino detectors. It is easy to estimate the consequences of this dynamic solar model (for standard neutrinos) regarding the SNO observations. Our picture modifies the conclusion that the \[NC\]/\[CC\] is larger than unity (\[i\]=observed rate/standard solar model rate, see Bahcall, Krastev, Smirnov, 2000) will definitely indicate the presence of neutrino observations. Since at high temperatures like $`10^{10}10^{11}K`$ the hot bubbles may produce mu and tau neutrinos independently of the presence of neutrino oscillations, they may contribute to \[NC\] increasing it. For example, if we take $`T=0.942`$, than we will have $`R_{SK}^{qc}=0.23`$, and so \[NC\]/\[CC\]=2.04. If we allow a value of T closer to the SSM value 1, the ratio $`([NC]/[CC])_{DSM}`$ decreases towards unity. This conjecture of the DSM should be taken into account at the interpretations of the SNO and SK observations.
Regarding the helioseismic context, we note that sound speed is not sensitive to the nuclear reactions. Bahcall and Ulrich (1988) and Basu, Pinsonneault and Bahcall (2000) remarked that even when switching out the $`He^3`$ \+ $`He^4`$ reaction (producing 14 $`\%`$ of the total solar luminosity), the sound speed differ only by 0.1 $`\%`$ from the sound speed obtained from the standard solar model. On the other hand, the calculations of the dynamic solar model show that when bubbles and neutrino oscillations are both present the SMA and LMA regions (of the purely MSW solutions) shift significantly even for such a small change as $`T=0.995`$. Moreover, the many-body effects of the particles in the nuclear reactions led to solar models which may compensate the cool solar models towards an SSM central temperature (Lavagno, Quarati 1999).
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# A norm on homology of surfaces and counting simple geodesics
## 1 Simple geodesics do not enter cusps
The result of this section goes back to Poincaré, while the sharp version for the torus stated below in Theorem 1 is proved in Greg McShane’s 1991 Warwick thesis . First, a definition:
Definition. A cusp region is a neighborhood of the cusp of $`S,`$ bounded by a horocycle.
Theorem 0.2. Let $`ϵ>0.`$ Any punctured torus has a cusp region with bounding curve of length $`4ϵ,`$ and this bound is optimal. No simple closed geodesic intersects a cusp region with boundary curve of length $`4ϵ.`$
For current purposes it will be sufficient to prove the following very easy theorem, weaker than Theorem 1 in the case of a punctured torus.
Theorem 0.3. Any cusped hyperbolic surface $`S`$ has a cusp region with bounding curve of length $`2.`$ No simple closed geodesic intersects a cusp region with boundary curve of length $`2.`$
Proof. In the upper half-plane model of $`H^2,`$ consider a fundamental domain of $`S,`$ arranged in such a way that the parabolic element preserving the cusp in question is $`\lambda :zz+1.`$ A simple closed geodesic $`g`$ of $`S`$ lifts to a geodesic in $`H^2,`$ which is represented by a semicircle $`\stackrel{~}{g}`$ in the half-space model. Since $`g`$ is simple, it follows that $`\stackrel{~}{g}\lambda \left(\stackrel{~}{g}\right)=\mathrm{}.`$ Thus, the radius of the semicircle representing $`\stackrel{~}{g}`$ is smaller than $`\frac{1}{2}.`$ The same goes for any non-vertical boundary component of the fundamental domain of $`S,`$ and thus the horocyclic arc joining (for example) $`\frac{1}{2}i`$ and $`1+\frac{1}{2}i`$ is entirely contained in the fundamental domain of $`S,`$ and meets no lift of any closed geodesic of $`S.`$ This horocyclic arc (which projects to a closed horocycle in $`S`$ has length 2.
## 2 Primitive elements in the fundamental group of the punctured torus
The reader should recall that the fundamental group $`\pi _1\left(T\right)`$ of the punctured torus is the free group on two generators $`F_2=s,t,`$ where $`s`$ and $`t`$ are the standard generators.
Definition. An element $`\gamma F_2`$ is called primitive, if there exists an automorphism $`\varphi `$ of $`F_2,`$ such that $`\psi \left(\gamma \right)=s.`$ If $`\varphi \left(\delta \right)=t,`$ then $`\gamma `$ and $`\delta `$ are called associated primitives.
Fact 0.4. The outer automorphism group of $`F_2`$ is isomorphic to the mapping class group of the punctured torus.
This was apparently first observed by Max Dehn, though Fact 2 has now passed into the folklore. One argument proceeds roughly as follows:
The outer automorphism group is generated by the so-called Nielsen transformation, which either permute the basis $`x,y,`$ or transform it into $`xy,y,`$ or $`xy^1,y.`$ In the case of a punctured torus these transformations can be topologically realized by Dehn twists.
Fact 0.5. Let $`\psi :F_2\mathrm{𝖹𝖹}^2`$ be the canonical abelianizing homomorphism. Then, if $`\gamma _1`$ and $`\gamma _2`$ are two primitive elements in $`F_2`$ such that $`\psi \gamma _1=\psi \gamma _2,`$ then $`\gamma _1`$ and $`\gamma _2`$ are conjugate.
Fact 2 follows easily from the work of Nielsen. For the proof see .
The following result, and its proof seem to go back to Poincaré.
Fact 0.6. Let $`S`$ be a hyperbolic surface of finite volume, and let $`\gamma `$ be a non-trivial simple closed curve, whose corresponding covering transformation $`\mathrm{\Gamma }`$ is hyperbolic. Then there is a unique geodesic freely homotopic to $`\gamma ;`$ this geodesic is simple.
Proof. The existence of a unique geodesic $`\stackrel{~}{\gamma }`$ freely homotopic to $`\gamma `$ follows by a completely standard straightening argument: since $`\mathrm{\Gamma }`$ is hyperbolic, it has two fixed points on the circle at infinity of $`H^2,`$ as do all of its translates. The sought after geodesic $`\stackrel{~}{\gamma }`$ is the unique geodesic in $`H^2`$ joining those two fixed points. To show that $`\stackrel{~}{\gamma }`$ is simple is also not hard. Indeed, the simplicity of $`\gamma `$ implies that for any covering transformation $`\beta ,`$ the fixed points of $`\beta \mathrm{\Gamma }\beta ^1`$ do not separate those of $`\mathrm{\Gamma }`$ in the cyclic order at infinity (otherwise the corresponding lifts of $`\gamma `$ would have intersected.) However, this immediately implies that the straightened curves do not intersect, by looking at it in the projective model of $`H^2.`$
Lemma 0.7. No simple geodesic on the punctured torus separates.
Proof. Suppose there was a separating simple geodesic $`\gamma .`$ Let it cut the torus into two components $`T_1`$ and $`T_2,`$ and assume that $`T_1`$ is the component which does not contain the cusp. Double $`T_1`$ along its boundary to obtain a compact oriented hyperbolic surface $`S=2T_1.`$ The area of $`S`$ is at least $`4\pi ,`$ and thus the area of $`T_1`$ is at least $`2\pi .`$ But the area of the whole punctured torus is exactly $`2\pi ,`$ and thus $`\gamma `$ could not have been separating.
Note. The above argument can be easily adapted to show that there are no simple closed geodesics on the thrice-punctured sphere.
The following lemma is well known:
Lemma 0.8. For any pair of non-separating simple closed curves $`\gamma _1`$ and $`\gamma _2`$ on the punctured torus $`T,`$ there exists a homeomorphism of $`T`$ taking $`\gamma _1`$ to $`\gamma _2.`$
Proof. This is immediate from the classification of surfaces.
Fact 2 and Lemmas 2 and 2 combine to show that is that there is a one-to-one correspondence between simple geodesics and primitive homology classes on the punctured torus (a primitive homology class is a class $`(m,n),`$ such that $`m`$ and $`n`$ are relatively prime), and it is geometrically obvious that this geodesic is really none other than the familiar $`(m,n)`$ torus knot, or the $`(m,n)`$ geodesic on the flat torus (without the puncture). The precise result is given by the following construction of Osborne and Zieschang :
Let $`m`$ and $`n`$ be a pair of relatively prime non-negative integers. Define a function $`f_{m,n}:\mathrm{𝖹𝖹}\{1,2\}`$ as follows: $`f_{m,n}\left(k\right)=f_{m,n}\left(k^{}\right),`$ if $`k=k^{}`$ modulo $`m+n;`$ for $`k`$ in $`\{1,\mathrm{},m\}`$ let $`f_{m,n}\left(k\right)=1,`$ and for $`k`$ in $`\{m+1,\mathrm{},m+n\}`$ let $`f_{m,n}\left(k\right)=2.`$
Now let
$$W_{m,n}(x_1,x_2)=\underset{i=0}{\overset{m+n1}{}}x_{f_{m,n}\left(1+im\right)}.$$
If $`m<0,`$ let $`W_{m,n}(x_1,x_2)=W_{m,n}(x_1^1,x_2),`$ and if $`n<0,`$ let $`W_{m,n}(x_1,x_2)=W_{m,n}(x_1,x_2^1).`$
The main theorem of is the following:
Theorem 0.9. If $`x_1`$ and $`x_2`$ are associated primitives in $`F_2=s,t,`$ and if $`m`$ and $`n`$ are relatively prime, then $`W_{m,n}(x_1,x_2)`$ is a primitive. furthermore, if $`mqnp=1`$ then $`W_{m,n}`$ and $`W_{p,q}`$ are associated primitives. In particular, up to conjugation the primitives of $`F(s,t)`$ are $`\{W_{m,n}(s,t)|m,n\mathrm{𝖹𝖹},(m,n)=1\}.`$
It can be seen that the words $`W_{m,n}(s,t)`$ are all cyclically reduced, hence of minimal word length in their conjugacy class. Since the word length of $`W_{m,n}`$ is equal to $`m+n`$ the following combinatorial version of Theorem A norm on homology of surfaces and counting simple geodesics holds:
Theorem 0.10. The number of conjugacy classes of primitive elements in $`F(s,t)`$ with reduced length not exceeding $`L`$ is asymptotic to $`L^2.`$
Proof. The number conjugacy classes in question equals
$$f\left(L\right)=|\left\{(m,n)\right|m+nL,(m,n)=1\}.$$
It follows from elementary number theory that $`f\left(L\right)`$ is asymptotically equal to $`L^2/\left(2\zeta \left(2\right)\right).`$
Remark 0.11. For free groups of higher rank the analogue of Theorem 2 fails spectacularly, in that the number of primitive elements of reduced length not exceeding $`L`$ grows exponentially, as demonstrated by the following construction due to Casson . Let $`F_n`$ be the free group generated by $`x_1,\mathrm{},x_n.`$ Now let $`y_i,`$ $`1in`$ be defined as follows: $`y_i=x_i`$ when $`i<n,`$ and $`y_n=x_nx_{i_k}^{j_k},`$ where $`i_k<n`$ for all $`k,`$ $`i_ki_{k+1}.`$ Clearly, $`\left\{y_i\right\}`$ is a generating set for $`F_n,`$ and also each $`y_i`$ is a cyclically reduced word. It is also clear that the number of possibilities for $`y_n`$ of length not exceeding $`L`$ grows exponentially in $`L,`$ for $`n>2.`$
## 3 Geometry of group actions
First, some definitions and a theorem, all directly from .
Definition. A geometry is a metric space in which each bounded set has compact closure. A group $`G`$ acts geometrically on a space $`X`$ if $`X`$ is a geometry, and there is a homomorphism $`\varphi `$ (usually suppressed) from $`G`$ into the isometry group of $`X`$ such that the $`G`$-action is properly discontinuous and cocompact (properly discontinuous means that for each compact set $`K`$ in $`X,`$ the set
$$\left\{gG\right|\mathrm{}KgK\}$$
is finite), and cocompact means that the orbit space $`X/G`$ is compact).
Definition. An intrinsic metric on a path space $`X`$ is one where the distance between two points is the infimum of path lengths between those points.
Definition. A relation $`R:XY`$ between spaces $`X`$ and $`Y`$ is said to be quasi-Lipschitz if $`R`$ is everywhere defined, and there exist positive numbers $`K`$ and $`L,`$ such that for each $`AX,`$
$$diamR\left(A\right)KdiamA+L.$$
Definition. Relations $`R:XY`$ and $`S:YX`$ are quasi-inverses if they are everywhere defined and there exists a constant $`M>0,`$ such that $`d(SR,id_X)<M`$ an $`d(RS,id_Y)<M.`$
Definition. A relation $`R:XY`$ is a quasi-Lipschitz equivalence if there is a quasi-inverse $`S:YX,`$ such that both $`R`$ and $`S`$ are quasi-Lipschitz.
Finally, a theorem:
Theorem 0.12. If a group acts geometrically on two geometries $`X`$ and $`Y`$ with intrinsic metrics, then $`X`$ and $`Y`$ are quasi-Lipschitz equivalent.
Having dispensed with the generalities, an observation:
Observation 0.13. If a finitely presented group $`G`$ is equipped with the word metric, then $`G`$ is a geometry, on which $`G`$ acts geometrically by multiplication from the right.
Now, restrict to the case of current interest, where $`G=F_2`$ acts on the hyperbolic plane $`H^2`$ with quotient the punctured torus. If this action was geometric, in the sense of the above definitions, then Theorem A norm on homology of surfaces and counting simple geodesics would follow from Theorem 2, Theorem 3, and the Observation 3. However, the action is not geometric, since the hyperbolic punctured torus is not compact. It is not hard to construct the right geometry.
Consider the tessellation of $`H^2`$ by fundamental domains for the given $`F_2`$ action, and truncate each fundamental domain by a closed horocycle of length $`2`$ (as is possible by Theorem 1). The geometry $`X`$ will be the subset of $`H^2`$ equal to the union of all these truncated fundamental domains, and equipped with path metric. Clearly $`X`$ is a geometry, on which $`F_2`$ acts geometrically, and so $`X`$ is quasi-isometric (quasi-Lipschitz equivalent) to $`F_2`$ equipped with its word metric.
While the general minimal paths on $`X`$ may be quite far from hyperbolic geodesics, by Theorem 1, the minimal paths corresponding to lifts of simple closed geodesics are unaffected by the truncation, and so their length are within a constant factor of the length of the corresponding reduced words in the word metric.
## 4 Geometry of multicurves
Let $`S`$ be a hyperbolic surface of finite volume with at most one cusp. We define a multicurve $`m`$ on $`S`$ to be a map from a (not necessarily connected) $`1`$-manifold $`M`$ to $`S.`$ We define the length of $`m`$ to be the sum of the lengths of the images of components of $`M.`$ We say that a multicurve $`m`$ is embedded if the image of $`m`$ is the union of simple closed curves $`\gamma _1,\mathrm{},\gamma _k`$ on $`S.`$ Note that the map $`m`$ may cover some of the components $`\gamma _i`$ multiple times, so this does not coincide with the usual meaning of embedding. A multicurve defines a singular chain, which, in turn, defines a homology class in $`H_1(S,\mathrm{𝖹𝖹}).`$
The first observation is the following:
Theorem 0.14. Let $`hH_1(S,\mathrm{𝖹𝖹})`$ be a non-trivial homology class. There exists a multicurve $`m`$ representing $`h`$ of minimal length, and $`m`$ is embedded, with all components geodesic.
Proof. First, we show the existence. If $`S`$ is a compact surface, this follows by a standard Arzela-Ascoli argument. If $`S`$ has cusps, let $`l_h`$ be the infimum of the lengths of multicurves representing $`h;`$ clearly $`l_h>0.`$ Let $`m_1,\mathrm{},m_l,\mathrm{}`$ be multicurves whose lengths approach $`l_h.`$ In order to apply the Arzela-Ascoli theorem, it is enough to show that the diameters of the images of the $`m_i`$ are uniformly bounded. Since the lengths of $`m_i`$ approach $`l_h,`$ it follows that the diameters of all the components of $`m_i,`$ for $`i`$ sufficiently great, are uniformly bounded by $`2l_h.`$ On the other hand, we can assume that no component of $`m_i,`$ for $`i`$ sufficiently large, can be contained entirely in a horodisk of area $`2`$ surrounding each cusp, since such a component is homologically trivial, and thus we can delete such a component with no change to homology, and replace $`m_i`$ by the resulting multicurve $`m_i^{}.`$ Hence, we can assume that all of the $`m_i`$ are contained in a compact subset of $`S,`$ and thus the existence follows.
Now let $`m`$ be a multicurve of minimal length representing $`h.`$ It is clear that each component of $`m`$ is geodesic. Suppose that $`m`$ is not embedded, thus two components $`c_1`$ and $`c_2`$ of $`m`$ intersect. We can assume that $`c_1`$ and $`c_2`$ are not multiply covered (if they are, we can split off one circle off each), and since they are geodesic, the intersection is transverse. Let $`O`$ be an intersection of $`c_1`$ and $`c_2,`$ and let $`A_1OB_1`$ and $`A_2OB_2`$ be small directed segments of $`c_1`$ and $`c_2,`$ respectively, surrounding $`O.`$ By cutting and pasting, we can replace these by $`s_1=A_1OB_2`$ and $`s_2=A_2OB_1,`$ without changing the homology, and then by smoothing $`s_1`$ and $`s_2`$ at $`O,`$ we obtain a shorter multicurve than $`m`$ representing $`h,`$ thus arriving at a contradiction.
Corollary 1. Let $`T`$ be a punctured torus equipped with a hyperbolic structure. Then, the shortest multicurve representing a non-trivial homology class $`h`$ is a simple closed geodesic if $`h`$ is a primitive homology class (that is, not a multiple of another class), and a multiply covered geodesic otherwise. In addition, the shortest multicurve representing $`h`$ is unique.
Proof. Since there is exactly one hyperbolic geodesic in any homotopy class, it is enough to observe that any two non-homotopic curves on a punctured torus intersect. Theorem 4 then implies that the shortest multicurve representing $`h`$ has one component, which is multiply covered if $`h`$ is not primitive (if $`h=(m,n)`$, with $`d=\mathrm{gcd}(m,n),`$ then the shortest multicurve representing $`h`$ is covered $`d`$ times). Since there is at most one simple closed geodesic in a homology class (see, eg section 2), the uniqueness follows.
## 5 A norm on homology of the punctured torus
First, let us define a valuation $`\mathrm{}`$ on $`H_1(T,\mathrm{𝖹𝖹}),`$ where $`\mathrm{}\left(h\right)`$ is defined to be the length of the shortest multicurve representing $`h.`$ The valuation of the trivial homology class is defined to be $`0.`$ Corollary 4 implies that
$$\mathrm{}\left(nh\right)=n\mathrm{}\left(h\right),$$
((0.4))
and that
$$\mathrm{}\left(h+g\right)\mathrm{}\left(h\right)+\mathrm{}\left(g\right),$$
((0.9))
where the inequality is strict if $`h`$ and $`g`$ are not both multiples of the same homology class. The latter inequality holds, because the union of the shortest multicurves corresponding to $`h`$ and $`g`$ is not embedded, and hence, the shortest multicurve corresponding to $`h+g`$ is shorter then the union. It follows that $`\mathrm{}`$ can be extended to the rational homology $`H_1(T,\text{ }\mathrm{Q}),`$ by linearity (equation (0.5)), and further, to $`H_1(T,R)`$ by continuity, which follows from equations (0.5) and (0.10). Since $`\mathrm{}\left(0\right)=0,`$ and by equations (0.5) and (0.10), $`\mathrm{}`$ is a pseudo-norm on the two-dimensional vector space $`H_1(T,R).`$ By the results of sections 1, 2 and 3, $`\mathrm{}`$ is actually a norm, since the results of those sections imply that $`0<c_1<\mathrm{}\left(h\right)/h_1<c_2<\mathrm{},`$ where $`h_1=\left|m\right|+\left|n\right|`$ for $`h=(m,n).`$ It follows that the unit ball of the norm $`\mathrm{}`$ is a compact, convex figure $`B_{\mathrm{}}`$ in the plane, and the number of simple geodesics of length not exceeding $`L`$ on the torus $`T`$ is equal to the number of primitive lattice points in $`LB_{\mathrm{}}.`$ Thus, Theorem A norm on homology of surfaces and counting simple geodesics follows.
## 6 Further investigations.
The geometry of the unit ball of the norm $`\mathrm{}`$ is very interesting; the authors will discuss it in a future paper. In particular, the study of the unit ball can be used to show that the error term in Theorem A norm on homology of surfaces and counting simple geodesics is essentially sharp. See .
The methods of this paper can be extended without much difficulty to define a norm on homology of surfaces of higher genus, by extending the length of shortest multicurve valuation, and this can be used to count minimal multicurves. It seems non-trivial to use this to count simple geodesics on surfaces of higher genus. Using the theory of train tracks and the methods of this note it can be shown (McShane and Rivin, in preparation) that the number of geodesics of length not exceeding $`L`$ on a closed surface of genus $`g`$ with $`c`$ cusps has order of growth $`L^{6g6+2c}.`$
Ecole Normale Superieure, Lyon Melbourne University, Victoria
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# FUSE OBSERVATIONS OF THE LOW-REDSHIFT Ly𝛽 FOREST
## 1. INTRODUCTION
Since the discovery of the high-redshift Ly$`\alpha `$ forest over 25 years ago, these abundant absorption features in the spectra of QSOs have been used as evolutionary probes of the intergalactic medium (IGM), galactic halos, large-scale structure, and chemical evolution. The Hubble Space Telescope (HST/FOS) Key Project on QSO absorption systems found that Ly$`\alpha `$ absorbers persist to low redshift in surprisingly large numbers (Bahcall et al. 1991; Morris et al. 1991; Jannuzi et al. 1998; Weymann et al. 1998). In this paper, we assume that the Ly$`\alpha `$ (and Ly$`\beta `$) lines are intergalactic. Richards et al. (1999) discuss the possibility that some C IV absorption systems could be intrinsic to the AGN or ejected at relativistic velocities.
The Colorado group has used HST to conduct a major survey of Ly$`\alpha `$ absorbers at low redshift ($`z0.07`$) along 15 AGN sightlines, using the moderate-resolution (19 km s<sup>-1</sup>) GHRS spectrograph (Stocke et al. 1995; Shull, Stocke, & Penton 1996; Penton et al. 2000a,b). Additional moderate-resolution HST/STIS data along 13 sightlines were taken during HST cycle 7. These observations measured Ly$`\alpha `$ absorbers down to equivalent widths of 10–20 mÅ and determined distributions of the low-$`z`$ Ly$`\alpha `$ absorbers in H I column density for $`12.3<\mathrm{log}N_{\mathrm{HI}}<14.0`$ and in line width for $`15<b_{\mathrm{width}}<100`$ km s<sup>-1</sup> (Penton et al. 2000b). The distribution function, $`d𝒩/dN_{\mathrm{HI}}N_{\mathrm{HI}}^{1.80\pm 0.05}`$, together with photoionization corrections suggest (Shull et al. 1999a) that the low-$`z`$ Ly$`\alpha `$ forest may contain a significant fraction of the baryons predicted by nucleosynthesis models of D/H (Burles & Tytler 1998).
However, a precise baryon census in the low-$`z`$ IGM requires measurements of the true doppler parameter, $`b`$, to obtain accurate column densities in saturated Ly$`\alpha `$ lines with log N$`{}_{\mathrm{HI}}{}^{}13.5`$. Measurements of Ly$`\beta `$ or higher Lyman lines are needed to constrain the degree of saturation through a curve-of-growth (COG) analysis. Indeed, recent ORFEUS studies (Hurwitz et al. 1998) of Ly$`\beta `$/Ly$`\alpha `$ ratios in two absorbers toward 3C 373 suggest that $`b`$ is less than the line width determined from Ly$`\alpha `$ profile fitting.
With the goal of characterizing the distribution of $`b`$-values and measuring accurate H I columns, we conducted a FUSE mini-survey of Ly$`\beta `$ absorbers toward 7 AGN with well-known Ly$`\alpha `$ lines. The FUSE mission and its capabilities are described by Moos et al. (2000) and Sahnow et al. (2000). In § 2 we describe our FUSE Ly$`\beta `$ observations. We also compare simple COG (Ly$`\beta `$/Ly$`\alpha `$) estimates of $`b`$ with single-component fits to (HST) Ly$`\alpha `$ lines to understand the kinematic structure of the low-$`z`$ absorbers. In § 3 we present our conclusions and give directions for future work on the IGM baryon content, line kinematics, and temperature.
## 2. FUSE OBSERVATIONS AND DATA ANALYSIS
Our observations were obtained with the FUSE satellite from 1999 September to November during the commissioning phase. Because the initial FUSE observations were taken with relatively long durations (16–49 ksec) we were able to obtain reasonable signal-to-noise ratios (S/N $``$ 15–20) in the LiF channels. The SiC channels had lower S/N and were not available for all targets. The spectral resolution was generally 20–25 km s<sup>-1</sup>, based on an analysis of narrow interstellar lines of H<sub>2</sub>, Ar I, and Fe II in our H<sub>2</sub> mini-survey (Shull et al. 2000).
Our FUSE Ly$`\beta `$ mini-survey consisted of 7 AGN chosen from targets for which the Colorado group had previously obtained HST spectra with GHRS/G160M or STIS/E140M. Our three GHRS targets and full line lists are: H1821+643 (Savage, Sembach, & Lu 1995; Tripp, Lu, & Savage 1998; Penton et al. 2000a), ESO 141-G55 (Sembach, Savage, & Hurwitz 1999; Penton et al. 2000a), and PKS 2155-304 (Shull et al. 1998; Penton et al. 2000a). Data on our four STIS targets (Mrk 876, PG 0804+761, VII Zw 118, and Ton S180) will be published separately.
From these 7 targets, we chose 12 strong Ly$`\alpha `$ absorbers with rest-frame equivalent widths $`W_\lambda >200`$ mÅ and relatively simple line profiles. Each of the Ly$`\alpha `$ lines had a detectable Ly$`\beta `$ counterpart, and in several cases we detected higher Lyman series lines. Figure 1 shows three examples: the 5512 km s<sup>-1</sup> Ly$`\beta `$ absorber toward Ton S180; the $`z=0.225`$ Ly$`\delta `$ absorber toward H1821+643; and the 10,463 km s<sup>-1</sup> Ly$`\beta `$ absorber toward ESO 141-G55. Although PKS 2155-304 contains several strong Ly$`\alpha `$ lines near $`cz=17,000`$ km s<sup>-1</sup>, their line blending is sufficiently complex that we did not include them in the survey. For the 12 chosen lines, we measured equivalent widths by gaussian profile fitting, using standard FUSE pipeline software. Assuming a single-component curve of growth, we used the concordance of the Ly$`\beta `$/Ly$`\alpha `$ ratios (or higher Lyman lines) to derive $`b`$ and N<sub>HI</sub>. After minimizing the $`\chi ^2`$ of our fits to $`N_{\mathrm{HI}}`$ and $`b`$, we constructed $`\mathrm{\Delta }\chi ^2`$ contours to define confidence regions for these inferred quantities. The error bars on $`N_{\mathrm{HI}}`$ and $`b`$ represent single-parameter 68% confidence intervals for each component and assume that the true line shape is well represented by a single-component, doppler-broadened line.
If Ly$`\alpha `$ and Ly$`\beta `$ are unsaturated, their equivalent widths are $`W_\lambda ^{Ly\alpha }=(54.5\mathrm{m}\mathrm{\AA })N_{13}`$ and $`W_\lambda ^{Ly\beta }=(7.37\mathrm{m}\mathrm{\AA })N_{13}`$, where N$`{}_{\mathrm{HI}}{}^{}=(10^{13}\mathrm{cm}^2)N_{13}`$. Thus, weak absorbers should have a ratio Ly$`\beta `$/Ly$`\alpha `$ = 0.135, and saturation gradually increases this ratio towards 1. Our typical FUSE Ly$`\beta `$ detection limit of 40 mÅ ($`4\sigma `$) corresponds to log N$`{}_{\mathrm{HI}}{}^{}13.73`$. Our success in detecting Ly$`\beta `$ for all Ly$`\alpha `$ lines with $`W_\lambda >200`$ mÅ is consistent with the relative line strengths.
Table 1 lists the 12 lines in our survey, together with Ly$`\alpha `$ velocities ($`cz`$) and line widths ($`b_{\mathrm{width}}`$) derived from instrumentally corrected gaussian fits to the GHRS or STIS Ly$`\alpha `$ profiles. Here, $`b_{\mathrm{width}}=2^{1/2}\sigma _{\mathrm{gauss}}`$, and $`b`$ and N<sub>HI</sub> are derived from the COG. Figures 2 and 3 show COG concordance plots for two absorbers in Fig. 1. Not including the broad, blended line at 16,203 km s<sup>-1</sup> toward PKS 2155-304, the FUSE distribution has mean $`b=31.4\pm 7.4`$ km s<sup>-1</sup> and median 28 km s<sup>-1</sup>, comparable to the values, $`b=27.5\pm 1.3`$ km s<sup>-1</sup> and median 26.4 km s<sup>-1</sup>, measured in the $`z=`$ 2.0–2.5 Ly$`\alpha `$ forest (Rauch et al. 1993). If thermal, these $`b`$-values correspond to $`T_{\mathrm{HI}}=m_Hb^2/2k50,000`$ K, a temperature higher than values predicted from models of low-metallicity Ly$`\alpha `$ clouds (Donahue & Shull 1991), heated and photoionized by the AGN metagalactic background. These models predict temperatures and doppler parameters, conveniently approximated by $`T=(24,300K)(U/0.01)^{0.152}`$ and $`b=(20\mathrm{km}\mathrm{s}^1)(U/0.01)^{0.076}`$. The doppler parameter scales weakly with the photoionization parameter $`U=n_\gamma /n_H(0.005)J_{23}(10^5\mathrm{cm}^3/n_H)`$. Here, $`U`$ is the ratio of ionizing photons to hydrogen nuclei for a specific intensity $`J_0=(10^{23}`$ ergs s<sup>-1</sup> cm<sup>-2</sup> sr<sup>-1</sup> Hz$`{}_{}{}^{1})J_{23}`$.
The doppler parameters inferred from Ly$`\beta `$/Ly$`\alpha `$ are considerably less than widths derived from Ly$`\alpha `$ profile-fitting. The statistical average, $`b/b_{\mathrm{width}}=0.52`$, suggests that the Ly$`\alpha `$ and Ly$`\beta `$ line profiles are more complex than single-component, doppler-broadened gaussians. The widths of individual components must be less than the COG-inferred $`b`$-values, a general principle exemplified by the special case where $`N`$ identical, well-separated gaussians behave as a single one with a dispersion equal to $`Nb`$. For each ensemble of components, the total column density derived from Ly$`\alpha `$ and Ly$`\beta `$ might be larger than the true value, but not by a large factor (Jenkins 1986).
The Ly$`\alpha `$ absorbers could be broadened by cosmological expansion, $`\mathrm{\Delta }v=(19.5\mathrm{km}\mathrm{s}^1)(\mathrm{\Delta }r/300\mathrm{kpc})h_{65}`$, across spatially extended H I absorbers, as seen in numerical simulations (Weinberg, Katz, & Hernquist 1998). The profiles could also be blends of velocity components arising from clumps of gas falling into dark-matter potential wells. If these clumps are sufficiently massive to be gravitationally bound, the velocity components may represent small-scale power at $`\mathrm{\Delta }v<100`$ km s<sup>-1</sup> in the cosmological spectrum of density fluctuations. We have seen spectral evidence for such velocity components in the low-$`z`$ Ly$`\alpha `$ absorbers studied by HST/GHRS and STIS at 19 km s<sup>-1</sup> resolution and in their two-point correlation function, $`\xi (\mathrm{\Delta }v)`$ (Penton et al. 2000b,c). We have also verified, through optical and 21-cm imaging, that some Ly$`\alpha `$ absorbers arise within small groups (van Gorkom et al. 1996; Shull et al. 1998).
With better statistics on the Ly$`\alpha `$ and Ly$`\beta `$ line widths and doppler parameters, we should be able to constrain the kinematics of the absorber profiles. A critical issue is whether cosmological expansion or velocity components are the dominant contributor to line broadening. Recent numerical simulations of the low-$`z`$ Ly$`\alpha `$ forest (Davé et al. 1999) suggest that both effects are important.
## 3. CONCLUSIONS AND DISCUSSION
First and foremost, FUSE has identified the low-redshift Ly$`\beta `$ forest. Lines between 1216 Å and (1216 Å)$`(1+z_{\mathrm{em}}`$) with no clear identification are often labeled as Ly$`\alpha `$. The FUSE detection of their Ly$`\beta `$ counterparts makes these identifications conclusive. We have used the COG concordance of Ly$`\beta `$/Ly$`\alpha `$ equivalent widths, and occasionally higher Lyman lines, to derive reliable values of $`b`$ and N<sub>HI</sub> for the stronger, saturated lines ($`W_\lambda >200`$ mÅ). Our major findings are: (1) The doppler parameters from single-component COG fits to Ly$`\beta `$/Ly$`\alpha `$ equivalent widths are considerably less than those derived from line profile fitting, with $`b/b_{\mathrm{width}}=0.52`$; (2) We find $`b=31.4\pm 7.4`$ km s<sup>-1</sup> (median 28 km s<sup>-1</sup>), similar to values at redshifts $`z=`$ 2.0–2.5; (3) Although these $`b`$-values correspond to $`T_{\mathrm{HI}}50,000`$ K, the low-$`z`$ absorbers may contain non-thermal motions or line broadening from cosmological expansion and infall.
Over its lifetime, FUSE will observe many AGN sightlines for the O VI, D/H, and AGN projects. This will produce a large survey of Ly$`\beta `$ absorbers at $`z<0.155`$ that we can use to characterize the distributions in $`b`$ and N<sub>HI</sub> in the low-redshift Ly$`\beta `$ forest. The distribution function, $`f(b)`$, can be used to derive the IGM temperature distribution and infer its equation of state (Schaye et al. 1999; Ricotti, Gnedin, & Shull 2000). With high-S/N data and a flux-limited sample, we can search for the hot baryons predicted by cosmological simulations (Cen & Ostriker 1999a). These absorbers should appear in the high-$`b`$ tail of the distribution as broad, shallow absorbers.
Because Ly$`\alpha `$ lines with $`W_\lambda >130`$ mÅ (log N$`{}_{\mathrm{HI}}{}^{}>13.5`$) appear to be saturated (Penton et al. 2000b), HST alone cannot provide accurate column densities for the strong Ly$`\alpha `$ absorbers that probably dominate the baryon content and opacity of the low-$`z`$ IGM. For lines in which $`b`$ is well determined, log $`N_{\mathrm{HI}}`$ typically increases by 0.3 dex, and up to 1 dex, compared to Ly$`\alpha `$ profile fitting. This increase, which is greatest in the most saturated Ly$`\alpha `$ lines, means that an even larger fraction of baryons may be found in the low-$`z`$ Ly$`\alpha `$ forest.
A full Ly$`\beta `$ survey will also provide a sample of high-$`N_{\mathrm{HI}}`$ absorbers which can be used to estimate the level of metallicity in the low-$`z`$ IGM (Shull et al. 1998; Cen & Ostriker 1999b). As an illustration, we constructed simple photoionization models using CLOUDY (Ferland 1996) that assume a specific ionizing intensity at 13.6 eV of $`J_\nu =10^{23}`$ ergs cm<sup>-2</sup> s<sup>-1</sup> sr<sup>-1</sup> Hz<sup>-1</sup> (Shull et al. 1999b) with $`J_\nu \nu ^{1.8}`$ and column density $`\mathrm{log}N_{\mathrm{HI}}=15`$. For total hydrogen densities (cm<sup>-3</sup>) of $`\mathrm{log}n_H=(3,4,5)`$, a 20 mÅ measurement of C III $`\lambda `$977 implies $`Z/Z_{}=(0.3,0.15,0.01)`$. For $`\mathrm{log}n_H=(5,6)`$, a 20 mÅ measurement of O VI $`\lambda `$1032 implies $`Z/Z_{}=(0.015,0.006)`$. Observations of C III $`\lambda 977`$ will be especially useful, because they can be compared with C II and/or C IV to obtain reliable ionization corrections.
This work is based on FUSE Team data obtained for the Guaranteed Time Team by the NASA-CNES-CSA FUSE mission operated by Johns Hopkins University. Financial support to U. S. participants has been provided by NASA contract NAS5-32985. The Colorado group also acknowledges support from astrophysical theory grants from NASA (NAG5-7262) and NSF (AST96-17073), the HST/COS project (NAS5-98043), and STScI grant GO-0653.01-95A which supported the Ly$`\alpha `$ data analysis.
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# Pion Cloud and Nucleon Mass in Finite Nuclei
## I Introduction
Until the pioneering work of Drukarev and Levin on the use of QCD sum rules in systems of finite nucleon density , descriptions of nucleons in nuclear matter had mostly been based on models of $`N`$-$`N`$ interactions or the direct use of meson exchange potentials in quantum hadrodynamics . These models have serious problems at both the short-distance and long-distance scales. At short distance one must treat the complex structure of hadrons, for which we now believe that QCD gives the correct description. At long-distance scales it seems that meson cloud effects must be included, and neither quantum hadrodynamics nor current treatments of QCD are adequate; while chiral perturbation theory is very successful in a quantitative treatment of pionic effects at low energies. Recently a model was developed for use in QCD sum rules with a nucleon current that has a Goldstone boson component, which has the main features of both microscopic QCD and chiral perturbation theory . In the present paper we extend this model for the study of nucleons in finite nuclei.
The strongest experimental evidence for the utility of introducing an explicit meson cloud component for the nucleon is the data on sea quark distributions obtained from Drell-Yan experiments. The NMC/CERN experiments have given evidence for violation of the Gottfried sum rule
$$_0^1\frac{dx}{x}\left[F_2^p\left(x\right)F_2^n\left(x\right)\right]=\frac{1}{3}_0^1\left\{\left[u_p^v\left(x\right)d_p^v\left(x\right)\right]+2\left[\overline{u}_p^v\left(x\right)\overline{d}_p^v\left(x\right)\right]\right\}𝑑x.$$
(1)
If the up-type and down-type sea quark distributions are equal, as would be expected from QCD, then the value of this integral is $`1/3`$. The NMC result is $`0.24\pm 0.016`$ at $`Q^2=5\text{GeV}^2`$ over the interval $`0.004x0.8`$. Similar effects were found in Drell-Yan measurements from Fermilab/E866 , where the ratio of down to up sea quark distributions is much larger than 1.0. Both of these results differ from the predictions of perturbative and nonperturbative QCD calculations. One approach to account for this discrepancy in sea quark distributions in hadronic models is to include a meson cloud for the nucleon . Our approach is to treat the pion as a basic Goldstone boson field at low energy, along with quarks and gluons, with separation of long-distance from medium- and short-distance effects.
Another strong motivation for introducing a correlator with a Goldstone boson component for use in QCD sum rules is given by chiral perturbaton theory. It was shown almost two decades ago that the leading nonanalytic contribution (LNAC) to the nucleon mass in chiral perturbation theory is a m$`{}_{}{}^{3}{}_{\pi }{}^{}`$ correction, with the lowest-order chiral logs not contributing. In the standard formulation of QCD sum rules one does not obtain such a m$`{}_{}{}^{3}{}_{\pi }{}^{}`$ term . In Ref it was shown not only that one obtains the m$`{}_{}{}^{3}{}_{\pi }{}^{}`$ LNAC, but that the known value of the mass shift (15 MeV) can be used to help determine the parameters of the model, in spite of the fact that the sum rule method is only accurate to about 10% for the nucleon mass. It is not possible to obtain the LNAC m$`{}_{}{}^{3}{}_{\pi }{}^{}`$ term with the usual quark fields correlator . We make use of this result in the present work.
It should also be noted that in the QCD sum rule method the microscopic QCD calculation gives $`m_\pi ^2\mathrm{ln}(m_\pi ^2)`$ terms for the nucleon in vacuum and also $`\rho m_\pi `$ terms for nucleons in the nuclear medium, both in contradiction to chiral perturbation theory. It was shown that both for the vacuum and for nuclear matter a careful treatment of the contunuum part of the phenomenological correlator enables one to cancel these terms, which are inconsistent with chiral perturbation theory. We made use of this cancellation of unwanted logs in the formulation of our model and use this mechanism in the present work.
We review Goldstone boson/QCD model of Ref. for a nucleon in the vacuum in Section 2. In Section 3 the surface effects on the pion propagator are quantified using a p-wave $`\pi A`$ optical potential where the $`\mathrm{\Delta }(1232)`$ resonance is dominant. We extend the model for the microscopic correlator from zero to finite density in Section 4, with the goal of calculating nucleon mass shifts due to the pion cloud effects in small symmetric finite nuclei. The sum rules are then constructed and evaluated in the final sections.
## II Review of Nucleon Correlator with Pion Cloud in Vacuum
In this section we review the model for the nucleon correlator with explicit Goldstone boson fields developed in Ref . The nucleon current has the form
$$\eta ^N(x)=c_1\eta ^{N,0}(x)+c_2\eta ^{N,\pi }(x)$$
(2)
where the constants $`c_i`$ are the amplitudes of the composite field operator without and with the Goldstone boson field. This operator is used to construct the nucleon correlator. We ignore possible contributions to the nucleon from the strange part of the meson cloud. The proton field operator without the pion cloud is chosen to be the current
$$\eta ^{p,0}(x)=ϵ^{abc}\left[u^a(x)^TC\gamma _\mu u^b(x)\right]\gamma _5\gamma ^\mu d_c(x),$$
(3)
where $`a,b,c`$ label the color indices, and $`u(x),d(x)`$ are the up and down quark fields. $`C`$ is the charge conjugation operator. The lowest-energy contribution to the phenomenological dispersion relation for the correlator, the pole term giving the proton intermediate state, depends on the transition matrix element
$$0\left|\eta ^{p,0}(p)\right|\text{proton}=\lambda _pv(p),$$
(4)
where $`\lambda _p`$ is a structure constant and the Dirac spinor $`v(p)`$ is normalized by
$$\overline{v}(p)v(p)=2m.$$
(5)
The neutron current $`\eta ^{n,0}(x)`$ can be obtained by the interchange of up and down quark fields in Equation (3). The nucleon current including the explicit Goldstone boson field is taken to be
$`\eta ^{p,\pi }(x)`$ $`=`$ $`{\displaystyle \frac{1}{\lambda _\pi ^2}}_\alpha \varphi _\pi (x)\tau \gamma ^\alpha \gamma ^\beta \eta ^{N,0}(x),`$ (6)
where $`\varphi _\pi (x)`$ is a massless pion field, $`\tau `$ is the I-spin operator and $`\lambda _\pi `$ is a $`D=1`$ scale parameter. The structure constant of this current is defined by the relation
$$0\left|\eta ^{p,\pi }(x)\right|\text{proton}=\lambda _p^{}v(p).$$
(7)
The p-wave coupling of the pion to the nucleon current and experience with hybrid mesons suggests that $`\lambda _p^2\lambda _p^2`$, so we expect little contribution to the phenomenological expression for our nucleon correlator from the resulting pole due to its very small residue.
The full correlator for the proton with the pion cloud is
$`\mathrm{\Pi }^p(p)`$ $`=`$ $`i{\displaystyle d^4x\text{e}^{ipx}0\left|T\left[\eta ^p(x)\overline{\eta }^p(0)\right]\right|0}`$ (8)
$`=`$ $`c_1^2i{\displaystyle d^4x\text{e}^{ipx}0\left|T\left[\eta ^{p,0}(x)\overline{\eta }^{p,0}(0)\right]\right|0}+(1c_1^2)i{\displaystyle d^4x\text{e}^{ipx}0\left|T\left[\eta ^{p,\pi }(x)\overline{\eta }^{p,\pi }(0)\right]\right|0}`$ (9)
$`=`$ $`c_1^2\mathrm{\Pi }^{(p,0)}(p)+(1c_1^2)\mathrm{\Pi }^{(p,\pi )}(p).`$ (10)
With the appropriate changes of current and isospin considerations, the neutron correlator is obtained. This model focuses on the long-range effects of the pion cloud and therefore does not include pion-quark interactions and the coupling between currents that include the pion cloud and currents that do not. It is assumed that these short-range contributions are already accounted for in QCD through the condensates, and their inclusion in the model would result in double counting some of the pionic effects on the nucleon correlator.
The part of the proton correlator without the pion cloud can be found in Ref.. The nucleon correlator with the meson cloud for the microscopic evaluation of the sum rule has the form
$$\mathrm{\Pi }^{(p,\pi )}(x)=D_{f,\alpha \beta }^\pi (x)\gamma ^\alpha \mathrm{\Pi }^{(p,0)}(x)\gamma ^\beta ,$$
(11)
where $`\mathrm{\Pi }^{(p,0)}(x)`$ is the coordinate space nucleon correlator without the pion cloud and in the chiral limit $`D_{f,\alpha \beta }^\pi (x)`$ has the form
$$D_{f,\alpha \beta }^\pi (x)=\frac{1}{\pi ^2x^4}\left(\frac{4x_\alpha x_\beta }{x^2}\delta _{\alpha \beta }\right).$$
(12)
In momentum space the nucleon correlator with the pion cloud is
$`\mathrm{\Pi }^{(p,\pi )}(p)`$ $`=`$ $`2{\displaystyle \frac{d^4k}{\left(2\pi \right)^4}\frac{(p/k/)\mathrm{\Pi }^{(p,0)}(k)(p/k/)}{\left(pk\right)^2}}`$ (14)
$`=`$ $`p/\mathrm{\Pi }_{\text{odd}}^{(p,\pi )}(p)+\mathrm{\Pi }_{\text{even}}^{(p,\pi )}(p).`$ (15)
The pion cloud contributions to the nucleon correlator are shown in the diagrams of Figure 1, and in momentum space are
$`\mathrm{\Pi }_{1a}(p)`$ $`=`$ $`{\displaystyle \frac{i}{2^{12}35\pi ^6\lambda _\pi ^4}}p^8\mathrm{ln}(p^2)p/`$ (17)
$`\mathrm{\Pi }_{1b}(p)`$ $`=`$ $`{\displaystyle \frac{i\overline{q}q^2}{2^33\pi ^2\lambda _\pi ^4}}p^2\mathrm{ln}(p^2)p/`$ (18)
$`\mathrm{\Pi }_{1d}(p)`$ $`=`$ $`{\displaystyle \frac{ig_s^2GG}{2^{11}3\pi ^6\lambda _\pi ^4}}p^4\mathrm{ln}(p^2)p/`$ (19)
$`\mathrm{\Pi }_{1e}(p)`$ $`=`$ $`{\displaystyle \frac{i\overline{q}q\overline{q}g_s^2\sigma Gq}{2^5\pi ^2\lambda _\pi ^4}}\mathrm{ln}(p^2)p/.`$ (20)
Details are given in Ref. . The usual QCD sum rule methods are used, as in Ref. , but with a modification of the treatment of the continuum given by the work of Ref. . In addition to the calculation of the nucleon mass the magnetic dipole moments of the nucleons were calculated. From the comparison of the results for the magnetic dipole moment to the result without the pion cloud component it was estimated that the pion cloud component has roughly the same probability as the cloudless component, or $`c_1^2c_2^20.5`$.
An important observation is that the diagrams shown in Figure 2 gives the main LNAC $`m_\pi ^3`$ terms. The largest contribution is from the first term shown in Figure 2a, which was shown to give the following contribution to the nucleon mass
$`\mathrm{\Delta }M_p`$ $`=`$ $`{\displaystyle \frac{em_\pi ^3g_A^2}{160\beta ^2f_\pi ^2}}`$ (21)
$``$ $`20\text{MeV},`$ (22)
with the choice of $`\lambda _\pi ^2`$ = M$`{}_{}{}^{2}{}_{p}{}^{}`$/g<sub>πN</sub>, which is quite reasonable.
The contribution to the nucleon mass from the pion cloud was found to be negligible, consistent with the assumption $`\lambda _p^2\lambda _p^2`$. We use this value for $`\lambda _\pi ^2`$ in the present work. This completes the model within the uncertainty of the probability of the meson cloud component.
## III Pion Propagator in Finite Nuclei
Extending the model for the nucleon correlator with the pion cloud from vacuum to finite nuclei requires an appropriate modification of the pseudo-Goldstone boson field in medium. In order to be able to able to apply the model at high momentum transfer, which is needed for various applications, we make use of conventional multiple scattering formalism \[see, e.g., Ref. for a review\] rather than chiral perturbation theory. This is a mean field method in which the pion propagator in the nucleon medium is given by the pion optical potential.
The propagator for a free massless pion has the form
$$D_f^\pi (k)=\frac{1}{k^2+iϵ}.$$
(23)
In the medium the pion acquires an effective mass given by the optical potential. The propagator for a pion with energy $`k_0`$ in a nucleus of mass number $`A`$ is
$$D_A^\pi (k)=\frac{1}{k^2+\mathrm{\Pi }(k_0,\stackrel{}{k};\rho ,A)+iϵ}$$
(24)
where $`\mathrm{\Pi }(k_0,\stackrel{}{k};\rho ,A)`$ is the pion self-energy in the medium and $`\rho `$ is the baryon density of the matter. The pion self energy is related to the pion-nucleus optical potential $`V_{\text{opt}}(k_0)`$ by
$$\mathrm{\Pi }(k_0,\stackrel{}{k},\stackrel{}{k}^{})=2k_0\stackrel{}{k}^{}\left|V_{\text{opt}}^{\pi A}(k_0)\right|\stackrel{}{k},$$
(25)
where $`\stackrel{}{k}`$ and $`\stackrel{}{k}^{}`$ are respectively the incoming and outgoing three-momenta of the pion. We begin constructing the $`\pi A`$ optical potential by looking at the partial wave expansion of the $`\pi N`$ scattering amplitude
$$(\stackrel{}{k}^{},\stackrel{}{k})=\underset{I}{}\widehat{P}_I\left(\underset{l}{}\left[\left(l+1\right)f_{I,l^+}(\omega )+lf_{I,l^{}}(\omega )\right]P_l(\mathrm{cos}\theta )\text{i}\stackrel{}{\sigma }\left(\widehat{k}^{}\widehat{k}\right)\underset{l}{}\left[f_{I,l^+}(\omega )f_{I,l^{}}(\omega )\right]P_l^{}(\mathrm{cos}\theta )\right).$$
(26)
The partial wave amplitude is related to the S-matrix by
$$f_\alpha =\frac{1}{2\text{i}|\stackrel{}{k}|}\left(S_\alpha (\omega )1\right)$$
(27)
where the scattering matrix is found using the experimentally obtained phase shifts $`\delta _\alpha (\omega )`$ as
$$S_\alpha (\omega )=\text{e}^{2\text{i}\delta _\alpha (\omega )}.$$
(28)
In this notation $`\alpha `$ indexes the isospin $`I`$, orbital angular momentum $`l`$ and total angular momentum $`J`$ of the partial wave. The phase shifts are expected to depend on momentum proportional to $`k^{2l+1}`$, and we define the constants $`a_\alpha `$ as
$$a_\alpha =\frac{\delta _\alpha }{k^{2l+1}}.$$
(29)
These constants are the so-called scattering lengths ($`a_{2I}`$) for s-waves and the scattering volumes ($`b_{2I,2J}`$) for p-waves. These parameters for d-wave and higher order partial waves are very small and have negligible effect on $`\pi N`$ physics, so we truncate the expansion of the scattering amplitude at $`l=1`$. Then the $`\pi N`$ scattering amplitude can take the form
$$(\stackrel{}{k}^{},\stackrel{}{k})=a(k_0)+a_I(k_0)\left(\stackrel{}{t}\stackrel{}{\tau }\right)+\left[b(k_0)+b_I(k_0)\left(\stackrel{}{t}\stackrel{}{\tau }\right)\right]\stackrel{}{k}^{}\stackrel{}{k}+\text{i}\left[b_{sf}(k_0)+b_{Isf}(k_0)\left(\stackrel{}{t}\stackrel{}{\tau }\right)\right]\stackrel{}{\sigma }\left(\stackrel{}{k}^{}\times \stackrel{}{k}\right),$$
(30)
where $`\stackrel{}{t}`$ and $`\stackrel{}{\tau }`$ are the isospin matrices for the pion and the nucleon. The energy-dependent coefficients can be related to the scattering volumes and lengths at threshold; for example, $`b=\frac{1}{3}\left(4b_{33}+2b_{13}+2b_{31}+b_{11}\right)`$. Because we are concerned in this work with spin/isospin symmetric nuclei, the isospin dependent and spin-flip parts of the scattering amplitude vanish. Furthermore, the p-wave contibution is considerably larger than the s-wave contribution at low and medium energies because of the dominance of the $`\mathrm{\Delta }(1232)`$ $`3,3`$ resonance in $`\pi N`$ scattering. We are then left with the simple scattering amplitude for the spin and isospin averaged scattering of a pion off of a nucleon,
$$(\stackrel{}{k}^{},\stackrel{}{k})=b(k_0)\stackrel{}{k}^{}\stackrel{}{k}.$$
(31)
In order to take into consideration the off-energy-shell effects in our optical potential, we modify the result of ( 31) by considering the p-wave t-matrix for the $`\pi N`$ scattering in the separable form
$$\stackrel{}{k}^{}\left|t_{\pi N}^{p\text{wave}}\left(k_0\right)\right|\stackrel{}{k}=b(k_0)g(\stackrel{}{k}^{})g(\stackrel{}{k})\stackrel{}{k}^{}\stackrel{}{k},$$
(32)
where $`g(\stackrel{}{k})`$ is a form-factor that roughly parametrizes the inelastic channels in the $`\pi N`$ scattering at high energies. $`g(\stackrel{}{k})`$ is chosen be unity when $`\overline{k}`$ is at the on-energy-shell value, and to approach zero as $`\overline{k}`$ gets large, where $`\overline{k}|\stackrel{}{k}|`$. We use a monopole parametrization of the form factor
$$g(\stackrel{}{k})=\frac{m_\pi ^2+\mathrm{\Lambda }^2}{|\stackrel{}{k}|^2+\mathrm{\Lambda }^2},$$
(33)
with vertex factor $`\mathrm{\Lambda }`$ taken to be 700 MeV, for the description of the off-shell properties of the amplitude.
The energy dependence of the function $`b(k_0)`$ is modelled by using the Breit-Wigner form for the $`\mathrm{\Delta }(1232)`$, giving
$$b(k_0,\overline{k})=\left(\frac{f^2\omega _\mathrm{\Delta }}{3\pi m_\pi ^3}\right)\frac{k_0^1}{\omega _\mathrm{\Delta }k_0\frac{i}{2}\mathrm{\Gamma }_\mathrm{\Delta }(\overline{k},k_0)},$$
(34)
where $`\omega _\mathrm{\Delta }`$ is the energy of the $`3,3`$ resonance and $`f`$ is the pseudovector $`\pi N`$ coupling constant. We model the resonance so that the three-momentum dependence of the width $`\mathrm{\Gamma }_\mathrm{\Delta }`$ is cut off at low momenta and at the momentum where the resonance is saturated:
$$\mathrm{\Gamma }_\mathrm{\Delta }(\overline{k},k_0)=\frac{8}{3}\frac{f^2}{4\pi }\frac{\omega _\mathrm{\Delta }}{m_\pi ^2}\frac{1}{k_0}\{\begin{array}{cc}\overline{k}_m^3\hfill & \overline{k}\overline{k}_m\hfill \\ & \\ \overline{k}^3\hfill & \overline{k}_m\overline{k}\overline{k}_\mathrm{\Delta }\hfill \\ & \\ \overline{k}_\mathrm{\Delta }^3\hfill & \overline{k}\overline{k}_\mathrm{\Delta }\hfill \end{array}.$$
(35)
The lower limit starting at $`\overline{k}_m=30`$ MeV ensures that the $`\mathrm{\Delta }(1232)`$ does not vanish even at very low pion momenta, and the upper limit at the three-momentum $`\stackrel{}{k}_\mathrm{\Delta }`$ where the Breit-Wigner curve peaks provides a mechanism for the resonance to disappear at high energies as expected.
To obtain the pion-nucleus optical potential for elastic scattering we make use of three approximations:
* low density
* impulse
* local density.
The low density and impulse approximations lead to the standard first-order in $`\rho `$ form for the optical potential, and the local density approximation is used for finite nuclei to treat the interaction at the nuclear surface as a function of $`\rho `$. The p-wave pion optical potential for elastic scattering by an $`A`$-nucleon nucleus can then be written to first order in density as
$$\stackrel{}{k}^{}\left|V_{\text{opt}}^{\pi A}\right|\stackrel{}{k}=A\stackrel{}{k}^{}\left|t_{\pi A}^{p\text{-wave}}\right|\stackrel{}{k}\stackrel{~}{\rho }\left(\stackrel{}{k}\stackrel{}{k}^{}\right).$$
(36)
$`\stackrel{~}{\rho }(\stackrel{}{k}\stackrel{}{k}^{})`$ is the Fourier transform of the ground state nuclear density $`\rho `$(r). The resulting optical potential in coordinate space, using Eq.(32), has the form
$$V(\stackrel{}{r})=Ac(E_\pi )\rho (r).$$
(37)
This is the so-called Kisslinger potential without the s-wave scattering term. It has been shown to give an accurate treatment of medium-energy pion-nucleus elastic scattring . Due to the p-wave nature of the interaction derived from the $`\mathrm{\Delta }(1232)`$ resonance the interaction takes place mainly on the surface, and indeed vanishes in the interior of the nucleus if a constant density is used. This will be seen to be an important aspect of our present work.
We model the spatial dependence of the nuclear density as a second order harmonic oscillator
$$\rho (r)=\rho _0\left(1+\alpha \left(\frac{r}{c}\right)^2\right)\mathrm{exp}\left(\left(\frac{r}{c}\right)^2\right),$$
(38)
where $`\rho _0`$ is the density of nuclear matter, $`c`$ is a size parameter, and $`\alpha `$ provides some control over the shape, allowing for modelling of heavier nuclei (like $`{}_{}{}^{40}\text{Ca}`$) and more complicated distributions (for example, $`{}_{}{}^{16}\text{O}`$ with a local minimum at the center).
We now turn our attention to the pion propagator in medium. The propagator (24) is expanded in density as
$$D_A^\pi (k)=\left[D_A^\pi (k)\right]_{\rho _0=0}+\underset{\stackrel{}{q}}{}\stackrel{~}{\rho }(\stackrel{}{q})\left[\frac{\delta }{\delta \stackrel{~}{\rho }(\stackrel{}{q})}D_A^\pi (k)\right]_{\rho _0=0}+\frac{1}{2!}\underset{\stackrel{}{q}_1}{}\underset{\stackrel{}{q}_2}{}\stackrel{~}{\rho }(\stackrel{}{q}_1)\stackrel{~}{\rho }(\stackrel{}{q}_2)\left[\frac{\delta }{\delta \stackrel{~}{\rho }(\stackrel{}{q}_1)}\frac{\delta }{\delta \stackrel{~}{\rho }(\stackrel{}{q}_2)}D_A^\pi (k)\right]_{\rho _0=0}+\mathrm{}$$
(39)
where the first term is just the free pion propagator (23) and $`\stackrel{}{q}`$ is the momentum transfer in the scattering. The scattering angles are averaged over, and the pion propagator in the nucleus takes the form
$$D_A^\pi (k)=D_f^\pi (k)+\frac{8\pi ^{3/2}(A1)\rho _0f^2\omega _\mathrm{\Delta }}{m_\pi ^2c}g\left(\overline{k}\right)^2F_1(\overline{k},\alpha ,c)\left(\omega _\mathrm{\Delta }k_0\frac{i}{2}\mathrm{\Gamma }_\mathrm{\Delta }(\overline{k},k_0)\right)^1D_f^\pi (k)^2+𝒪\left(\rho _0^2\right),$$
(40)
with the nuclear structure information contained in the function
$`F_1(\overline{k},\alpha ,c)`$ $`=`$ $`\left(1+{\displaystyle \frac{3\alpha }{2}}\right)\left[\left(1{\displaystyle \frac{2}{c^2\overline{k}^2}}\right)+\left(2c^2\overline{k}^2+3+{\displaystyle \frac{2}{c^2\overline{k}^2}}\right)\mathrm{exp}\left(c^2\overline{k}^2\right)\right]`$ (42)
$`\alpha \left[\left(1{\displaystyle \frac{3}{c^2\overline{k}^2}}\right)+\left(2c^4\overline{k}^4+4c^2\overline{k}^2+5+{\displaystyle \frac{3}{c^2\overline{k}^2}}\right)\mathrm{exp}\left(c^2\overline{k}^2\right)\right].`$
## IV The Nucleon Correlator with Pion Cloud in Medium
To construct the in-medium nucleon correlator with the meson field, we start by rewriting the vacuum matrix element of the product $`\eta (x)\overline{\eta }(x)`$ in Eq.( 8) as a nuclear matrix element so that the proton current correlator takes the form
$`\mathrm{\Pi }_A^p(p)`$ $`=`$ $`i{\displaystyle d^4x\text{e}^{ipx}A\left|T\left[\eta ^p(x)\overline{\eta }^p(0)\right]\right|A}`$ (43)
$`=`$ $`c_1^2i{\displaystyle d^4x\text{e}^{ipx}A\left|T\left[\eta ^{p,0}(x)\overline{\eta }^{p,0}(0)\right]\right|A}+c_2^2i{\displaystyle d^4x\text{e}^{ipx}A\left|T\left[\eta ^{p,\pi }(x)\overline{\eta }^{p,\pi }(0)\right]\right|A}`$ (44)
$`=`$ $`c_1^2\mathrm{\Pi }_A^{(p,0)}(p)+c_2^2\mathrm{\Pi }_A^{(p,\pi )}(p),`$ (45)
where $`|A`$ is the ground state of the finite nucleus. The part of the in-medium proton correlator with the pseudo-Goldstone field is
$`\mathrm{\Pi }_A^{(p,\pi )}(x)`$ $`=`$ $`ϵ^{abc}ϵ^{a^{}b^{}c^{}}{\displaystyle \frac{1}{\lambda _\pi ^4}}\{{\displaystyle \frac{2}{3}}\left(_\alpha _\beta D_A^{\pi ^0}(x)\right)\gamma _\alpha \gamma ^\mu S_{d,A}^{cc^{}}(x)\gamma ^\nu \gamma _\beta \text{Tr}\left[S_{u,A}^{bb^{}}(x)\gamma _\nu CS_{u,A}^{aa^{}T}(x)C\gamma _\mu \right]`$ (48)
$`+{\displaystyle \frac{4}{3}}\left(_\alpha _\beta D_A^{\pi ^+}(x)\right)\gamma _\alpha \gamma ^\mu S_{u,A}^{cc^{}}(x)\gamma ^\nu \gamma _\beta \text{Tr}\left[S_{d,A}^{bb^{}}(x)\gamma _\nu CS_{d,A}^{aa^{}T}(x)C\gamma _\mu \right]\}`$
$`+\left[\text{4-quark terms}\right]`$
$`=`$ $`{\displaystyle \frac{1}{\lambda _\pi ^4}}\left\{{\displaystyle \frac{2}{3}}D_{A,\alpha \beta }^{\pi ^0}(x)\gamma ^\alpha \mathrm{\Pi }_A^{(p,0)}(x)\gamma ^\beta +{\displaystyle \frac{4}{3}}D_{A,\alpha \beta }^{\pi ^+}(x)\gamma ^\alpha \mathrm{\Pi }_A^{(n,0)}(x)\gamma ^\beta \right\}.`$ (49)
This is simply the expression for the meson cloud part of the correlator, given in Ref. , where the quark propagator and pion propagator in vacuum, $`S_q^{ab}(x)`$ and $`D_f^\pi (x)`$, are replaced by $`S_{q,A}^{ab}(x)`$ and $`D_A^\pi (x)`$, the quark and pion propagators in the nucleus, respectively. For example,
$`S_{q,A}^{ab}(x)`$ $`=`$ $`A\left|T\left[q^a(x)\overline{q}^b(0)\right]\right|A.`$ (50)
$`D_{A,\alpha \beta }^\pi (x)`$ is the nuclear analogue to Eq.(12), replacing the free pion propagator $`D_f^\pi (x)`$ with $`D_A^\pi (x)`$. As with the case in vacuum, the neutron correlator $`\mathrm{\Pi }_A^{(n,0)}(x)`$ is found simply by an interchange of up and down quark fields in the current $`\eta (x)`$. With a spin/isospin symmetric nucleus in the chiral limit, the in-medium correlator with a pion cloud has the form analogous to Eq.(11),
$$\mathrm{\Pi }_A^{(p,\pi )}(x)=\frac{2}{\lambda _\pi ^4}D_{A,\alpha \beta }^\pi (x)\gamma ^\alpha \mathrm{\Pi }_A^{(p,0)}(x)\gamma ^\beta .$$
(51)
For simplicity we work in momentum space, where the in-medium proton correlator has the form
$$\mathrm{\Pi }_A^p(p)=c_1^2\mathrm{\Pi }_A^{(p,0)}(p)(1c_1^2)\frac{2}{\lambda _\pi ^4}\frac{d^4k}{\left(2\pi \right)^4}D_A^\pi (k)k/\mathrm{\Pi }_A^{(p,0)}(pk)k/$$
(52)
where we made use of the Fourier transform of $`D_{A,\alpha \beta }^\pi (x):`$
$$D_{A,\alpha \beta }^\pi (p)=p_\alpha p_\beta D_A^\pi (p).$$
(53)
Since the nucleus has a four-momentum $`u^\mu `$, breaking Lorentz invariance, there are additional Dirac structures in the nucleon correlator. Under the assumption that the nuclear ground state is invariant under parity and time-reversal, it was shown that the in-medium nucleon correlator is
$$\mathrm{\Pi }_A^p(p)=\mathrm{\Pi }_A^{(s)}(p^2,pu)+\mathrm{\Pi }_A^{(p)}(p^2,pu)p/+\mathrm{\Pi }_A^{(u)}(p^2,pu)u/.$$
(54)
The individual scalar functions $`\mathrm{\Pi }_A^{(n)}(p^2,pu)`$ depend only on the Lorentz invariants $`p^2`$ and $`pu`$ and can be projected out of the full nucleon correlator using the following formulae :
$`\mathrm{\Pi }_A^{(s)}(p^2,pu)`$ $`=`$ $`{\displaystyle \frac{1}{4}}\text{Tr}\left[\mathrm{\Pi }_A^p(p)\right]`$ (55)
$`\mathrm{\Pi }_A^{(p)}(p^2,pu)`$ $`=`$ $`{\displaystyle \frac{1}{p^2\left(pu\right)^2}}\left[{\displaystyle \frac{1}{4}}\text{Tr}\left[p/\mathrm{\Pi }_A^p(p)\right]{\displaystyle \frac{pu}{4}}\text{Tr}\left[u/\mathrm{\Pi }_A^p(p)\right]\right]`$ (56)
$`\mathrm{\Pi }_A^{(u)}(p^2,pu)`$ $`=`$ $`{\displaystyle \frac{1}{p^2\left(pu\right)^2}}\left[{\displaystyle \frac{1}{4}}\text{Tr}\left[u/\mathrm{\Pi }_A^p(p)\right]{\displaystyle \frac{pu}{4}}\text{Tr}\left[p/\mathrm{\Pi }_A^p(p)\right]\right].`$ (57)
We define in-medium matrix elements of an operator $`𝒪\left(x\right)`$ as an expansion in local nuclear density $`\rho `$
$`𝒪\left(x\right)_\rho `$ $`=`$ $`A\left|𝒪\left(x\right)\right|A0\left|𝒪\left(x\right)\right|0`$ (58)
$`=`$ $`\rho N\left|𝒪\left(x\right)\right|N+\mathrm{},`$ (59)
where $`N\left|𝒪\left(x\right)\right|N`$ is the single-nucleon matrix element of $`𝒪\left(x\right)`$. The in-medium quark condensate is related to the $`\pi `$-$`N`$ sigma commutator by
$$\overline{q}q_\rho =\rho \frac{\sigma _{\pi N}}{m_u+m_d}.$$
(60)
We take $`\sigma _{\pi N}47\text{MeV}`$ . The breaking of the Lorentz invariance of the medium by the introduction of the nuclear four-momentum $`u^\mu `$ leads to condensates with additional Dirac structures . The dimension three in-medium vector quark condensate, with the appropriate color and flavor weighting, is proportional to the nuclear density
$$q^{}q_\rho =\frac{3}{2}\rho .$$
(61)
The in-medium gluon condensate is found via the trace anomaly to be
$$g_s^2G_{\mu \nu }G^{\mu \nu }_\rho =\rho 2^639\pi ^2\text{MeV}.$$
(62)
The values of the five-dimensional mixed quark-gluon condensates in medium are not well known, although significant cancellation between diagrams containing these matrix elements occurs much like similar terms in the free nucleon correlator in the chiral limit. Furthermore, it was shown in that the nucleon sum rules are insensitive to the values of these condensates, so we ignore terms proportional to the mixed condensates in our sum rule analysis.
The nucleon correlator is analyzed in the rest frame of the nucleus, where the nucleus four-momentum is $`u=(M,0,0,0)`$ with the mass of the nucleus $`MAm`$ (neglecting the average binding energy per nucleon in the nucleus as small compared to the nucleon mass $`m`$). In this frame $`pu=2Mp_0`$, where $`p_0`$ is the energy of the interpolating field. We carry out Borel transforms with respect to $`p^2`$, using the relationship
$$\left(\frac{M}{A}+p\right)^2=4m^2.$$
(63)
In the nuclear rest frame, $`p_0=\frac{1}{2m}\left(3m^2p^2\right)`$, so when $`p^2\mathrm{}`$, $`pu\mathrm{}`$ while the ratio of the two invariants, analogous to the Bjorken scaling variable in DIS, remains fixed and the OPE is applicable.
We restrict our calculation of the nucleon correlator to first order in the nuclear density. We can then identify three distinct types of quark current diagrams that contribute to the correlator: diagrams that contribute to the part of the in-medium correlator without the pseudo-Goldstone field (Figure 3), diagrams that include the pion field where the quarks and gluons interact with the medium and the pion maintains its in-vacuum properties (Figure 4), and diagrams where the quarks and glue interact with the QCD vacuum and the pion scatters in the nucleus (Figure 5). The shaded boxes in these diagrams represent interaction of the field with the nuclear medium to first order in the density expansion. Diagrams where the quarks and pions both interact with the medium, or where the pion scattering includes two-body correlations would come in at second order in density and are not included in this work.
The first two sets of diagrams are most easily calculated in coordinate space using (51), after which they are Fourier transformed into momentum space. The contributions at first order in the density from diagrams without the pion cloud are found to be
$`\mathrm{\Pi }_{(2a,\rho )}(p)`$ $`=`$ $`{\displaystyle \frac{i}{3\pi ^2}}q^{}q_\rho \left[2u/p^2+p/up\right]\mathrm{ln}\left(p^2\right)`$ (64)
$`\mathrm{\Pi }_{(2b.\rho )}(p)`$ $`=`$ $`{\displaystyle \frac{i}{\left(2\pi \right)^2}}\overline{q}q_\rho p^2\mathrm{ln}\left(p^2\right)`$ (65)
$`\mathrm{\Pi }_{(2c,\rho )}(p)`$ $`=`$ $`{\displaystyle \frac{i}{2^7\pi ^4}}g_s^2GG_\rho p/\mathrm{ln}\left(p^2\right)`$ (66)
$`\mathrm{\Pi }_{(2d,\rho )}(p)`$ $`=`$ $`{\displaystyle \frac{2^2i}{3}}\overline{q}qq^{}q_\rho {\displaystyle \frac{up}{p^2}}`$ (67)
$`\mathrm{\Pi }_{(2e+2f,\rho )}(p)`$ $`=`$ $`{\displaystyle \frac{i}{3^22\pi ^2}}\left[\overline{q}qg_s^2GG_\rho +\overline{q}q_\rho g_s^2GG\right]{\displaystyle \frac{1}{p^2}}`$ (68)
$`\mathrm{\Pi }_{(2g,\rho )}(p)`$ $`=`$ $`{\displaystyle \frac{i}{2^53^2\pi ^2}}q^{}q_\rho g_s^2GG\left[5u/{\displaystyle \frac{1}{p^2}}2p/{\displaystyle \frac{up}{p^4}}\right]`$ (69)
$`\mathrm{\Pi }_{(2h,\rho )}(p)`$ $`=`$ $`{\displaystyle \frac{2^2i}{3}}\overline{q}q\overline{q}q_\rho p/{\displaystyle \frac{1}{p^2}}`$ (70)
$`\mathrm{\Pi }_{(2i,\rho )}(p)`$ $`=`$ $`{\displaystyle \frac{2^217i}{3^3}}g_s^2\overline{q}q^2\overline{q}q_\rho {\displaystyle \frac{1}{p^4}}.`$ (71)
Here we have used a factorization approximation to evaluate the effects of the four-quark and six-quark condensates, where to first order in density we take $`\overline{q}\mathrm{\Gamma }q\overline{q}\mathrm{\Gamma }q_\rho \overline{q}q\overline{q}q_\rho `$ and $`\overline{q}\mathrm{\Gamma }q\overline{q}\mathrm{\Gamma }q\overline{q}\mathrm{\Gamma }q_\rho \overline{q}q^2\overline{q}q_\rho `$, respectively. While the factorization of these condensates is unjustified in medium , our primary concern is a comparison of the nucleon mass shift due to the pion cloud terms and the shift due to terms without the pion cloud, so we find the approximation acceptable for the problem at hand.
The diagrams in Figure 4, where there are first-order interactions of the quarks and gluons with the medium but the pions interact only with the vacuum, provide the following contributions at first order in the density:
$`\mathrm{\Pi }_{(3a,\rho )}(p)`$ $`=`$ $`{\displaystyle \frac{i}{2^63^2\pi ^4\lambda _\pi ^4}}q^{}q_\rho \left[u/p^2+3p/up\right]p^4\mathrm{ln}\left(p^2\right)`$ (72)
$`\mathrm{\Pi }_{(3c,\rho )}(p)`$ $`=`$ $`{\displaystyle \frac{i}{2^{13}\pi ^6\lambda _\pi ^4}}g_s^2GG_\rho p/p^4\mathrm{ln}\left(p^2\right)`$ (73)
$`\mathrm{\Pi }_{(3g,\rho )}(p)`$ $`=`$ $`{\displaystyle \frac{i}{2^83^3\pi ^4\lambda _\pi ^4}}q^{}q_\rho g_s^2GG\left[25u/p^2+11p/up\right]\mathrm{ln}\left(p^2\right)`$ (74)
$`\mathrm{\Pi }_{(3h,\rho )}(p)`$ $`=`$ $`{\displaystyle \frac{i}{2^23\pi ^2\lambda _\pi ^4}}\overline{q}q\overline{q}q_\rho p/p^2\mathrm{ln}\left(p^2\right).`$ (75)
All diagrams which would be expected to contribute to scalar structure of the nucleon correlator vanish.
The calculation of the terms represented in Figure 5 where the pions interact with the nucleus and the quarks and glue interact within the QCD vacuum are more difficult to calculate. Instead of calculating these diagrams in coordinate space as we did with the other in-medium diagrams, which would require first a rather daunting Fourier transform of the in-medium pion propagator (Eq. 40) and then an inverse transform of the resulting expression, we perform the calculations in momentum space using Equation (52) as the starting point. For example, the term corresponding to Figure 5a gives the linear-density contribution
$`\mathrm{\Pi }_{(4a,\rho )}(p)`$ $`=`$ $`{\displaystyle \frac{\zeta i}{2^{10}\pi ^8\lambda _\pi ^4}}{\displaystyle d^3\stackrel{}{k}\left(g(\overline{k})\right)^2F_1(\overline{k},\alpha ,c)}`$
$`\times {\displaystyle }dk_0k_0{\displaystyle \frac{\left(\omega _\mathrm{\Delta }k_0\right)k_0+\frac{i}{2}\mathrm{\Gamma }_\mathrm{\Delta }\left(\overline{k}\right)}{\left(\omega _\mathrm{\Delta }k_0\right)^2k_0^2+\frac{1}{4}\left(\mathrm{\Gamma }_\mathrm{\Delta }\left(\overline{k}\right)\right)^2}}{\displaystyle \frac{\left(pk\right)^4\mathrm{ln}\left[\left(pk\right)^2\right]}{\left(k_0^2\overline{k}^2+iϵ\right)^2}}[(2pkk^2)k/k^2p/]`$
where $`\zeta =\frac{8\pi ^{\frac{3}{2}}\left(A1\right)\rho _0f_\pi ^2\omega _\mathrm{\Delta }}{m_\pi ^2c}`$. The Borel transform with respect to $`p^2`$, $`\widehat{B}`$, eliminates terms proportional to $`pk`$ in the expansion of $`\left(pk\right)^{2n}\mathrm{ln}[\left(pk\right)^2]`$, simplifying the integrand so that the term proportional to $`k/`$ vanishes from symmetry in the angle-averaging. We are left with integrals over the pion energy $`k_0`$ and the magnitude of the pion three-momentum $`\overline{k}`$. The integrand also picks up a Borel dependence of the form $`M_B^{2n}\mathrm{exp}(k^2/M_B^2)`$. The energy integration can be completed analytically, where we find that the factor $`\mathrm{exp}(k^2/M_B^2)1`$ and the dependence on the Borel mass can be taken out of the integral. This first diagram is found to be
$$\widehat{B}\mathrm{\Pi }_{(4a,\rho )}(p)=\frac{\zeta i}{\lambda _\pi ^4}I(\alpha ,c)\frac{1}{\left(2\pi \right)^7}p/M_B^6$$
(76)
where $`I(\alpha ,c)`$ is the remaining one-dimensional integral over the pion three-momentum. Defining $`\stackrel{~}{\mathrm{\Gamma }}(\overline{k})=k_0\mathrm{\Gamma }(\overline{k},k_0)`$, the integral takes the form
$$I(\alpha ,c)=\frac{i}{2\pi ^2c^2}_0^{\mathrm{}}𝑑\overline{k}F_1(\overline{k},\alpha ,c)\left(g\left(\overline{k}\right)\right)^2\frac{\omega _\mathrm{\Delta }\overline{k}\overline{k}^2+\frac{i}{2}\stackrel{~}{\mathrm{\Gamma }}\left(\overline{k}\right)}{\left(\omega _\mathrm{\Delta }\overline{k}\overline{k}^2\right)^2+\frac{1}{4}\left(\stackrel{~}{\mathrm{\Gamma }}\left(\overline{k}\right)\right)^2}.$$
(77)
We calculate this integral numerically.
The other diagrams shown in Figure 5 yield the following contributions to the nucleon correlator:
$`\widehat{B}\mathrm{\Pi }_{(4b,\rho )}(p)`$ $`=`$ $`{\displaystyle \frac{\zeta i}{\lambda _\pi ^4}}I(\alpha ,c){\displaystyle \frac{4\pi }{\left(2\pi \right)^6}}\overline{q}qM_B^4`$ (78)
$`\widehat{B}\mathrm{\Pi }_{(4c,\rho )}(p)`$ $`=`$ $`{\displaystyle \frac{\zeta i}{\lambda _\pi ^4}}I(\alpha ,c){\displaystyle \frac{2^3\pi }{\left(2\pi \right)^43}}p/\overline{q}q^2`$ (79)
$`\widehat{B}\mathrm{\Pi }_{(4e,\rho )}(p)`$ $`=`$ $`{\displaystyle \frac{\zeta i}{\lambda _\pi ^4}}I(\alpha ,c){\displaystyle \frac{1}{2^9\pi ^7}}p/g_s^2GGM_B^2`$ (80)
$`\widehat{B}\mathrm{\Pi }_{(4f,\rho )}(p)`$ $`=`$ $`{\displaystyle \frac{\zeta i}{\lambda _\pi ^4}}I(\alpha ,c){\displaystyle \frac{1}{2^43\pi ^7}}p/\overline{q}q\overline{q}g_s\sigma Gq{\displaystyle \frac{1}{M_B^2}}`$ (81)
$`\widehat{B}\mathrm{\Pi }_{(4g,\rho )}(p)`$ $`=`$ $`{\displaystyle \frac{\zeta i}{\lambda _\pi ^4}}I(\alpha ,c){\displaystyle \frac{1}{2^33^2\pi ^5}}\overline{q}qg_s^2GG`$ (82)
$`\widehat{B}\mathrm{\Pi }_{(4h,\rho )}(p)`$ $`=`$ $`{\displaystyle \frac{\zeta i}{\lambda _\pi ^4}}I(\alpha ,c){\displaystyle \frac{17}{3^4\pi ^3}}g_s^2\overline{q}q^3{\displaystyle \frac{1}{M_B^2}},`$ (83)
where $`I(\alpha ,c)`$ in each of these expressions is the same as Equation (77).
## V Construction of the In-Medium Sum Rules
### A In the Chiral Limit
The phenomenological (“right-hand”) side of the sum rules, is expressed as a dispersion relation with a lowest-lying hadronic pole of in-medium mass $`m^{}`$ and residue $`\lambda _p^2`$ clearly separated from a continuum of higher-energy states with threshold $`s_0^{}`$:
$$\mathrm{\Pi }_{A,\text{RHS}}^p(p)=\lambda _p^2c_1^2\frac{p/+m^{}}{p^2m^2}+\frac{1}{\pi }_{s_0^{}}^{\mathrm{}}𝑑s\frac{\text{Im}\mathrm{\Pi }_A^{\text{cont}}(s,uq)}{sp^2}$$
(84)
Guided by the vacuum calculations, we assume that the structure parameter associated with the current including the pion cloud $`\lambda _p^{}_{}{}^{}2\lambda _p^2`$, so there is no additional pole considered in our model of the phenomenological correlator. Rather than find the values of the parameters in medium, we follow the appoach of Drukarev and Levin and determine the shifts of the in-medium parameters from their vacuum values. We remove the contributions to the nucleon correlator due to the pure in-vacua diagrams by defining the quantity
$$\varphi ^{(i)}2i\left(2\pi \right)^4\widehat{B}\left[\mathrm{\Pi }_A^{(i)}\left(p\right)\mathrm{\Pi }^{(i)}\left(p\right)\right],$$
(85)
where $`i`$ labels the three structures in the sum rule ($`\mathrm{𝟏}`$, $`p/`$, and $`u/`$, respectively). On the phenomenological side of the sum rule we can express $`\varphi ^{(i)}`$ in terms of the parameter shifts and the Borel transform of the RHS in vacuum as
$$\varphi ^{(i)}(M_B^2,s)=\left[\mathrm{\Delta }m\frac{}{m}+\mathrm{\Delta }\stackrel{~}{\lambda }^2\frac{}{\stackrel{~}{\lambda }^2}+\mathrm{\Delta }s_0\frac{}{s_0}\right]\widehat{B}\mathrm{\Pi }_{0,\text{RHS}}^{(i)}\left(p\right)$$
(86)
for $`i=s,p`$ (since there is no $`u`$ structure in vacuum). The contributions to the microscopic side of the sum rule are given in the previous section. With the subscript labelling the contributions due to each figure, we find in the nuclear rest frame using the constraint in Equation (63):
$`\varphi _2^{(s)}\left(M_B^2\right)`$ $`=`$ $`2\rho a_NM_B^4E_1+\left[32^3\pi ^2\rho m^2a_0{\displaystyle \frac{2^2}{3^2}}\rho \left(b_Na_0+b_0a_N\right)\right]+{\displaystyle \frac{217}{3^2\pi ^2}}g_s^2\rho a_0^3{\displaystyle \frac{1}{M_B^2}}`$ (88)
$`\varphi _2^{(p)}\left(M_B^2\right)`$ $`=`$ $`2^4\pi ^2\rho M_B^4E_1L^{\frac{4}{9}}+\left[{\displaystyle \frac{1}{2^2}}\rho b_N{\displaystyle \frac{2^4\pi ^2}{3}}m^2\rho \right]M_B^2E_0L^{\frac{4}{9}}+\left[{\displaystyle \frac{2^3}{3}}\rho a_0a_NL^{\frac{4}{9}}{\displaystyle \frac{\pi ^2}{32}}\rho b_0L^{\frac{4}{9}}\right]`$ (90)
$`\pi ^2\rho b_0m^2L^{\frac{4}{9}}{\displaystyle \frac{1}{M_B^2}}`$
$`\varphi _2^{(u)}\left(M_B^2\right)`$ $`=`$ $`2^6\pi ^2\rho M_B^4E_1L^{\frac{4}{9}}{\displaystyle \frac{5\pi ^2}{32}}\rho b_0L^{\frac{4}{9}}`$ (91)
$`\varphi _3^{(s)}\left(M_B^2\right)`$ $`=`$ $`0`$ (92)
$`\varphi _3^{(p)}\left(M_B^2\right)`$ $`=`$ $`{\displaystyle \frac{1}{\lambda _\pi ^4}}({\displaystyle \frac{3}{2^2}}\rho M_B^8E_3L^{\frac{4}{9}}+[{\displaystyle \frac{3}{2^2}}m^2\rho {\displaystyle \frac{1}{2^7\pi ^2}}\rho b_N]M_B^6E_2L^{\frac{4}{9}}`$ (94)
$`+[{\displaystyle \frac{1}{23\pi ^2}}\rho a_0a_NL^{\frac{4}{9}}{\displaystyle \frac{11\pi ^4}{2^53^2}}\rho b_0L^{\frac{4}{9}}]M_B^4E_1+{\displaystyle \frac{11\pi ^4}{2^53^2}}\rho b_0M_B^2E_0L^{\frac{4}{9}})`$
$`\varphi _3^{(u)}\left(M_B^2\right)`$ $`=`$ $`{\displaystyle \frac{1}{\lambda _\pi ^4}}\left({\displaystyle \frac{3}{2}}\rho M_B^8E_3L^{\frac{4}{9}}+{\displaystyle \frac{5}{2^33^2}}\rho b_0M_B^4E_1L^{\frac{4}{9}}\right)`$ (95)
$`\varphi _4^{(s)}\left(M_B^2\right)`$ $`=`$ $`{\displaystyle \frac{\zeta }{\lambda _\pi ^4}}\text{Re}\left[I(\alpha ,c)\right]\left[{\displaystyle \frac{a_0}{\pi ^3}}M_B^4E_1+{\displaystyle \frac{2a_0b_0}{3^2\pi ^3}}{\displaystyle \frac{17}{3^4\pi ^5}}g_s^2a_0^3{\displaystyle \frac{1}{M_B^2}}\right]`$ (96)
$`\varphi _4^{(p)}\left(M_B^2\right)`$ $`=`$ $`{\displaystyle \frac{\zeta }{\lambda _\pi ^4}}\text{Re}\left[I(\alpha ,c)\right]\left[{\displaystyle \frac{1}{2\pi ^3}}M_B^6E_2L^{\frac{4}{9}}{\displaystyle \frac{b_0}{2^3\pi ^3}}M_B^2E_0L^{\frac{4}{9}}+{\displaystyle \frac{2a_0^2}{3\pi ^3}}L^{\frac{4}{9}}{\displaystyle \frac{a_0^2m_0^2}{32^2\pi ^7}}{\displaystyle \frac{1}{M_B^2}}\right]`$ (97)
$`\varphi _4^{(u)}\left(M_B^2\right)`$ $`=`$ $`0.`$ (98)
Here we have used for the condensates
$`a_0`$ $`=`$ $`\left(2\pi \right)^2\overline{q}q=0.55\text{GeV}^3,`$
$`a_N`$ $`=`$ $`\left(2\pi \right)^2N\left|\overline{q}q\right|N,`$
$`b_0`$ $`=`$ $`g_s^2GG=0.474\text{GeV}^4,`$
$`b_N`$ $`=`$ $`N\left|g_s^2GG\right|N,\text{ and}`$
$`m_0^2`$ $`=`$ $`{\displaystyle \frac{\overline{q}g_s\sigma Gq}{a_0}}=0.8\text{GeV}^2.`$
The left-hand side of the sum rule is then $`\varphi ^{(i)}=c_1^2\varphi _2^{(i)}c_2^2\varphi _3^{(i)}c_2^2\varphi _4^{(i)}`$. We have taken into account perturbative corrections using the renormalization group anomalous dimensions of the interpolating field and the local operators used in constructing the LHS. Each term is multiplied by the appropriate power of the anomalous dimension factor
$$L\frac{\mathrm{ln}\left(M_B/\mathrm{\Lambda }_{\text{QCD}}\right)}{\mathrm{ln}\left(\mu /\mathrm{\Lambda }_{\text{QCD}}\right)}.$$
(99)
Since it arises from a conserved current, the anomalous dimension of the pseudo-Goldstone boson is zero. We set the renormalization point of the OPE $`\mu =500\text{MeV}`$ and the QCD scale parameter $`\mathrm{\Lambda }_{\text{QCD}}=100\text{MeV}`$. In the continuum model, the contributions from higher-energy states are included by using the prescription whereby terms proportional to $`M_B^{2n}`$ are multiplied by factors of $`E_{n1}`$ defined as
$$E_n1\mathrm{exp}\left(s_0^{}/M_B^2\right)\underset{k=0}{\overset{n}{}}\frac{1}{k!}\left(\frac{s_0^{}}{M_B^2}\right)^k.$$
(100)
We introduce a pair of functions $`\mathrm{\Phi }^{(i)}`$ as
$`\mathrm{\Phi }^{(s)}(M_B^2,s)`$ $`=`$ $`\varphi ^{(s)}(M_B^2,s)`$ (101)
$`=`$ $`c_1^2\left\{\stackrel{~}{\lambda }^2\left(1{\displaystyle \frac{2m^2}{M_B^2}}\right)e^{\frac{m^2}{M_B^2}}\mathrm{\Delta }m+me^{\frac{m^2}{M_B^2}}\mathrm{\Delta }\stackrel{~}{\lambda }^22as_0e^{\frac{s_0}{M_B^2}}\mathrm{\Delta }s_0\right\}`$ (102)
$`\mathrm{\Phi }^{(p)}(M_B^2,s)`$ $`=`$ $`m\varphi ^{(p)}(M_B^2,s),`$ (103)
$`=`$ $`c_1^2\{{\displaystyle \frac{2m^2}{M_B^2}}\stackrel{~}{\lambda }^2e^{\frac{m^2}{M_B^2}}\mathrm{\Delta }mme^{\frac{m^2}{M_B^2}}\mathrm{\Delta }\stackrel{~}{\lambda }^2`$ (105)
$`+{\displaystyle \frac{m}{2}}L^{\frac{4}{9}}[{\displaystyle \frac{s_0^4}{\lambda _\pi ^432^55}}+s_0^2(1{\displaystyle \frac{b}{\lambda _\pi ^432^5}})+{\displaystyle \frac{a^2L^{\frac{8}{9}}s_0}{\lambda _\pi ^432^2}}+({\displaystyle \frac{b}{2}}+{\displaystyle \frac{m_0^2a^2}{\lambda _\pi ^432^2}})]e^{\frac{s_0}{M_B^2}}\mathrm{\Delta }s_0\}.`$
The sum of these functions has no dependence on the residue shift $`\mathrm{\Delta }\stackrel{~}{\lambda }^2`$, and sufficient cancellation occurs between the terms proportional to the threshold shift $`\mathrm{\Delta }s_0`$ that we can ignore these terms. The in-nucleus mass shift of the nucleon can then be determined using
$$\mathrm{\Delta }m=\frac{1}{c_1^2\stackrel{~}{\lambda }^2}\mathrm{exp}\left(\frac{m^2}{M_B^2}\right)\left[\mathrm{\Phi }^{(s)}(M_B^2,s)+\mathrm{\Phi }^{(p)}(M_B^2,s)\right].$$
(106)
We adopt a derivative improvement to the sum rules suggested by Drukarev and Ryskin where
$$\psi _j^{(i)}\frac{M_B^4}{m^2}\frac{d\varphi _j^{(i)}}{dM_B^2}.$$
(107)
This improvement minimizes the effect of the factorization approximation of the four quark condensates by eliminating the terms proportional to $`p^2`$ from the dominant contributing diagrams, therefore aiding the convergence of the sum rules. Applying the derivative operator to Equations (86) we find
$`\psi _2^{(s)}\left(M_B^2\right)`$ $`=`$ $`2^2\rho a_N{\displaystyle \frac{1}{m^2}}M_B^6E_2{\displaystyle \frac{1}{m^2}}{\displaystyle \frac{217}{3^3\pi ^2}}g_s^2\rho a_0^3`$ (109)
$`\psi _2^{(p)}\left(M_B^2\right)`$ $`=`$ $`2^5\pi ^2\rho {\displaystyle \frac{1}{m^2}}M_B^6E_2L^{\frac{4}{9}}+\left[{\displaystyle \frac{1}{2^2}}\rho b_N{\displaystyle \frac{1}{m^2}}{\displaystyle \frac{2^4\pi ^2}{3}}\rho \right]M_B^4E_1L^{\frac{4}{9}}+\pi ^2\rho b_0L^{\frac{4}{9}}`$ (110)
$`\psi _2^{(u)}\left(M_B^2\right)`$ $`=`$ $`2^7\pi ^2\rho {\displaystyle \frac{1}{m^2}}M_B^6E_2L^{\frac{4}{9}}`$ (111)
$`\psi _3^{(s)}\left(M_B^2\right)`$ $`=`$ $`0`$ (112)
$`\psi _3^{(p)}\left(M_B^2\right)`$ $`=`$ $`{\displaystyle \frac{1}{\lambda _\pi ^4}}({\displaystyle \frac{3}{m^2}}\rho M_B^{10}E_4L^{\frac{4}{9}}+[{\displaystyle \frac{3^2}{2^2}}\rho {\displaystyle \frac{3}{2^7\pi ^2}}{\displaystyle \frac{\rho b_N}{m^2}}]M_B^8E_3L^{\frac{4}{9}}`$ (114)
$`+[{\displaystyle \frac{1}{3\pi ^2}}\rho {\displaystyle \frac{a_0a_N}{m^2}}L^{\frac{4}{9}}{\displaystyle \frac{11\pi ^4}{2^43^2}}{\displaystyle \frac{\rho b_0}{m^2}}L^{\frac{4}{9}}]M_B^6E_2+{\displaystyle \frac{11\pi ^4}{2^53^2}}{\displaystyle \frac{\rho b_0}{m^2}}M_B^4E_1L^{\frac{4}{9}})`$
$`\psi _3^{(u)}\left(M_B^2\right)`$ $`=`$ $`{\displaystyle \frac{1}{\lambda _\pi ^4}}\left({\displaystyle \frac{32}{m^2}}\rho M_B^{10}E_4L^{\frac{4}{9}}+{\displaystyle \frac{5}{2^33^2}}\rho {\displaystyle \frac{b_0}{m^2}}M_B^6E_2L^{\frac{4}{9}}\right)`$ (115)
$`\psi _4^{(s)}\left(M_B^2\right)`$ $`=`$ $`{\displaystyle \frac{\zeta }{\lambda _\pi ^4}}\text{Re}\left[I(\alpha ,c)\right]\left[{\displaystyle \frac{2a_0}{m^2\pi ^3}}M_B^6E_2+{\displaystyle \frac{17}{3^4\pi ^5}}{\displaystyle \frac{g_s^2}{m^2}}a_0^3\right]`$ (116)
$`\psi _4^{(p)}\left(M_B^2\right)`$ $`=`$ $`{\displaystyle \frac{\zeta }{\lambda _\pi ^4}}\text{Re}\left[I(\alpha ,c)\right]\left[{\displaystyle \frac{3}{2\pi ^3}}{\displaystyle \frac{1}{m^2}}M_B^8E_3L^{\frac{4}{9}}{\displaystyle \frac{b_0}{2^3\pi ^3}}{\displaystyle \frac{1}{m^2}}M_B^4E_1L^{\frac{4}{9}}+{\displaystyle \frac{a_0^2m_0^2}{32^2\pi ^7}}{\displaystyle \frac{1}{m^2}}\right]`$ (117)
$`\psi _4^{(u)}\left(M_B^2\right)`$ $`=`$ $`0`$ (118)
The improvement does not remove all dependence of the LHS on the four-quark condensate, but these remaining four-quark terms are largely suppressed in comparison to terms with similar Borel weighings. As with the unimproved sum rules, a linear combination of these expressions can be built that has no dependence on the residue shift and negligible dependence on the threshold shift, giving the mass shift of the nucleon as
$$\mathrm{\Delta }m=\frac{1}{c_1^2\stackrel{~}{\lambda }^2}\mathrm{exp}\left(\frac{m^2}{M_B^2}\right)\left[\psi ^{(s)}(M_B^2,s)m\psi ^{(p)}(M_B^2,s)\right].$$
(119)
The improved sum rule is used in the analysis of the nucleon mass shifts in this work.
### B Beyond the Chiral Limit
As discussed above, in Ref. it was shown that processes illustrated in Figure 2 produce terms proportional to $`m_\pi ^3/f_\pi ^2`$. Also the $`m_\pi ^2ln(m_\pi ^2)`$ terms were shown to cancel using the analysis of Ref. . It was thereby shown that our pion cloud model is consistent with chiral perturbation theory to leading order in the nonanalytic contributions to the nucleon mass in vacuum, and that we could use the chiral perturbation theory result to constrain the parameters of the model
In dense nuclear matter (and at finite temperature, Refs. ) another problem with chiral symmetry breaking seems to occur in the QCD sum rule method . Equation (60) can be rewritten as
$$\frac{\overline{q}q_\rho }{\overline{q}q_0}=\rho \frac{\sigma _{\pi N}}{f_\pi ^2m_\pi ^2}.$$
(120)
However, the LNAC to the $`\pi N`$ sigma commutator is proportional to $`m_\pi ^3`$ ). The leading nonanalytic contribution to the in-medium quark condensate is thereby seen to be of order $`\rho m_\pi `$, which is forbidden by chiral counting rules . Birse and Krippa accounted for this by adding contributions to the RHS where the nucleon current interacts with soft pions from the pion cloud of another nucleon in the nucleus . Using the soft pion theorem of PCAC they find a term proportional to $`\rho m_\pi `$ on the phenomenological side which cancels the offending term on the microscopic side arising from the in-medium quark condensate. As we go away from the chiral limit in our pion cloud model, we adopt the same modification to the phenomenological nucleon correlator, and in doing so we need to make sure that there are no contributions to the OPE from the pion cloud that would violate the chiral counting rules to leading nonanalytic order. This is indeed the case, since the chiral expansion of the pion propagator contributes lowest order terms proportional to $`m_\pi ^2`$, maintaining the leading order behavior. Likewise, similar to the vacuum contribution illustrated in Figure 2a lowest dimension in-medium contributions shown in Figures 6a and 6b are possible. However, their effects are on higher order nonanalytic terms as they provide corrections proportional to $`\rho m_\pi ^3`$ and $`\rho m_\pi ^4\mathrm{ln}(m_\pi ^2)`$. Thus, to first order in baryon density the pion cloud model for the nucleon correlator in a nucleus is consistent with chiral perturbation theory in the leading nonanalytic contributions to the nucleon mass.
## VI Results for the Nucleon Mass Shift in Spin-Isospin Symmetric Nuclei
This model is best applied to light nuclei where there are large surface volume regions and therefore the density gradients extend over a large fraction of the nucleus. In infinite nuclear matter these denisty gradients are zero, so that the s-wave part of the $`\pi A`$ optical potential, which we neglect dominates the in-medium pion cloud contribution to the nucleon correlator. In matter with $`Z=N`$, the strength of this potential is very small and the contributions of the diagrams in Figure 5 to the nucleon mass are negligible. As a result, the pion cloud has almost no effect on the nucleon mass in infinite nuclear matter, similar to the pion cloud effect on the nucleon mass in vacuum. It is evident that in large nuclei the pion cloud contribution to the nucleon mass becomes non-negligible only near the nuclear surface, a well-known property of the Kisslinger potential. So long as the nucleon remains a distance away from the surface comparable to about twice the $`\pi N`$ interaction range (with $`r_{\pi N}0.7\text{fm}`$) our mass shift result for this nucleon reproduce the work of previous authors . Also application to large nuclei would require an extension of the method to account for $`NZ`$.
In this paper, we work with models of two small symmetric nuclei that are familiar from $`\pi A`$ scattering experiments, $`{}_{}{}^{16}\text{O}`$ and $`{}_{}{}^{40}\text{Ca}`$. The nucleon density distribution of the ground states of both nuclei are given by Eq (38)). For $`{}_{}{}^{16}\text{O}`$ we choose the parameters $`c=1.75\text{fm}`$ and $`\alpha =2.0`$ . This produces a distribtuion of nucleons with a local minimum at $`r=0`$ and an RMS radius of $`2.6\text{fm}`$. The parameters for the $`{}_{}{}^{40}\text{Ca}`$ density distribution are taken as $`c=2.80\text{fm}`$ and $`\alpha =1.1`$.
We find the mass shifts in the chiral limit, ignoring in our analysis the effects of the small leading nonanalytic terms discussed above. The mass shift of a nucleon in each nucleus is found by evaluating Eq. (119) in a Borel window $`0.8\text{GeV}M_B^21.4\text{GeV}`$. If the plot of the mass shift as a function of the Borel mass is relatively flat in this window, then the sum rules are sufficiently convergent and the mass shift is taken at this flat value. To test the stability of the sum rules we look at the mass shift at two points in the nucleus: first in the middle of the surface where the local density is half the nuclear matter density $`\rho _0`$, then at the center of the nucleus, where the central density is equivalent to $`\rho _0`$. As can be seen in Figures 7-10, the sum rules in both nuclei provide sufficiently stable solutions in the Borel window.
The convergence is amplified in the plots of the individual contributions of the vector and scalar correlators (Fig. 11). We find that the nucleon mass shift in an $`{}_{}{}^{16}\text{O}`$ nucleus is negligible at the central baryon density, and $`150\text{MeV}`$ at one-half central density. The mass shifts in a $`{}_{}{}^{40}\text{Ca}`$ nucleus are $`260\text{MeV}`$ at $`\rho =\rho _0`$ and $`400\text{MeV}`$ at $`\rho =\rho _0/2`$.
Some pion cloud effects might be related to the spontaneous breaking of chiral symmetry. Since the reduction in the value the quark condensate is a signal of partial restoration of chiral symmtery in the nucleus, the pion cloud amplitude $`c_2`$ might decrease accordingly. We explore this idea by modifying our model, giving the pion cloud probability $`c_2^2`$ a linear dependence on density such that it has the vacuum value at zero density, and it vanishes at the density where the quark condensate $`\overline{q}q_\rho `$ goes to zero. Similarly, the probability of no pion cloud $`c_1^2`$ increases from the vacuum value of $`0.5`$ to the phase transition point value of $`1.0`$. Figure 12 shows the mass shift of a nucleon in a $`{}_{}{}^{40}\text{Ca}`$ nucleus at central density utilizing this modification to the model. The shift in $`{}_{}{}^{16}\text{O}`$ shows a similar, if less dramatic, reduction. As in the constant pion cloud probability calculation, the sum rule calculation is sufficiently flat in the Borel window to provide a value for the mass shift. The mass shifts in $`{}_{}{}^{16}\text{O}`$ are $`0\text{MeV}`$ at central and $`100\text{MeV}`$ at surface density. In $`{}_{}{}^{40}\text{Ca}`$ we find $`100\text{MeV}`$ and $`275\text{MeV}`$, respectively.
We compare these two methods to the sum rules without an explicit pion cloud throughout the nucleus in Figures 13 and 14. In these calculations, the density used at each point in the sum rules is the local density obtained from the density distribution (Eq. (38)) and the nucleon mass is plotted as a function of the distance from the center of the nucleus. Each mass shift is evaluated at the Borel mass $`M_B=1.0\text{GeV}`$.
In an $`{}_{}{}^{16}\text{O}`$ nucleus, the mass shift due to the pion cloud is small compared to the mass shift from the non-pionic terms but positive out to a distance of about $`4.5\text{fm}`$ from the nuclear center. As the nucleon appraoches the nuclear surface the pion cloud contribution to the mass shift becomes positive and larger than the mass shift without pion cloud. Even though the local density in the interval $`3.0\text{fm}r4.0\text{fm}`$ is small (hence the small mass shift due to the correlator contributions without the pion cloud), the pion cloud is still sampling a significant range of density gradients leading to a sizable optical potential strength, and therefore a large mass shift. The two different approaches to the pion cloud amplitudes in $`{}_{}{}^{16}\text{O}`$ do not yield significantly different results throughout the nucleus. In the $`{}_{}{}^{40}\text{Ca}`$ nucleus, the mass shift due to the pion cloud is always negative and comparable to the non-pionic mass shift. As with the oxygen nucleus, the largest mass shift from the pionic terms occurs near the surface, approaching $`600\text{MeV}`$. The constant and linearly decreasing pion cloud probability calculations converge to the same value in this region. The nuclear interior shows a significant difference between the three methods. If $`c_1`$ and $`c_2`$ remain constant over the range of densities, then the pion cloud provides a rather large mass shift ($`275\text{MeV}`$ at central density). Once $`c_1`$ and $`c_2`$ are dependent on density, the pion cloud contribution is more than halved at the center of the nucleus, with a mass shift closer to $`100\text{MeV}`$. Clearly the choice of pion cloud method has a significant effect in larger nuclei.
## VII Conclusions
We have described a model based on QCD for the nucleon in a finite spin and isospin symmetric nucleus including a pseudo-Goldstone boson component of the correlator This approach is based on the model proposed by one of us in Ref., with the extension of the pion free propagator to the in-medium propagator using an optical potential which is first order in the nuclear density. A fit to chiral perturbation theory restricts the parameters of the pion cloud component. Our main conclusion is that there are significant contributions to the mass of the nucleon in the nucleus arising from the pion cloud component of the nucleon.
In the deep interior of the oxygen nucleus we find a small positive mass shift which becomes large and negative near the surface. In the calcium nucleus, however, the mass shift is everywhere large and negative. We introduce a parametrization of the pion cloud probability $`c_2^2`$ which decreases linearly from the vacuum value to zero at the chiral symmetry restoration as a possible test of partial restoration of chiral symmetry. At high densities, the effect of the pion cloud is then reduced leading to smaller mass shifts.
One obvious extension to this model is the consideration of the strange sector. By allowing for the inclusion of kaon cloud terms in the baryon current, we can construct similar sum rules for the mass shift of hyperons. The primary obstacles are the uncertainties of several in-medium matrix elements on the microscopic side, the complicated structure of the dispersion relation on the phenomenological side , and the difficulty in finding an appropriate form for the $`K`$-$`A`$ optical potential to determine the kaon self-enegry. The model can also be modified to handle N $``$ Z nuclei. The nucleon correlator in such nuclei without the pion cloud has been discussed in Ref., and the optical potential is easily treated in asymmetric nuclei by including the isospin-dependent and spin flip terms in the $`\pi N`$ scattering amplitude (Eq. (30)). A third venue of further research concerns the applications of the model to sea quark distributions in nuclei and the development of an improvement to QCD calculations of nuclear deep inelastic scattering in the interval $`0.1x0.3`$.
We would like to thank Mountaga Aw and Otto Linsuain for helpful discussions.
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# 1 Introduction
## 1 Introduction
Physical consequences of extra space-like dimensions have been the subject of intensive theoretical investigations over the past two years. In general, the presence of new dimensions can change our interpretation of the fundamental parameters in the theory. Four-dimensional constants, such as Plank mass, electroweak scale, coupling constants and so on, are, in fact, functions of the volume of the extra dimensions. This dependence is particularly interesting in theories where gravity and the rest of the Standard Model fields “live” in different numbers of dimensions. This can be realized in the brane-world scenario, in which Standard Model particles are confined to a 3+1 space-time dimensional manifold, embedded into a higher-dimensional space. A very well known stringy prototype of these constructions , uses a relatively large size of the transverse 11-th dimension in supergravity to reconcile effective Plank scale and GUT/string scale. In such a construction the 4D-Planck scale, $`M_P`$, is derived from the fundamental scale, $`M_{11}`$, using $`M_P^2=M_{11}^3R_{11}`$ having assumed that the 6-space volume, $`V_6M_{11}^6`$. For $`R_{11}M_{GUT}`$, we can obtain, $`M_{GUT}M_{11}M_P`$, thus relaxing the hierarchy problem of the GUT and string scales.
Interest in such models, driven by particle phenomenology and cosmology, was further amplified when the brane world idea was used to lower the fundamental gravitational scale, bringing it close to the weak scale in an attempt to resolve the hierarchy problem between the weak and Planck scales. Indeed, as it was shown in , due to a possibly very large volume of extra dimensions felt only by gravity, the fundamental scale could be much lower than the one inferred from measuring the Newtonian force in four dimensions. This opens up a number of interesting phenomenological consequences, including the possibility of observing higher-dimensional gravitational interactions in collider experiments. An alternative to solve the weak scale hierarchy problem was proposed by Randall and Sundrum . Rather than requiring a large extra dimension, one simply constructs an exponentially large ratio of the two length scales on two different branes, which results from an “exponentially decaying” scale factor (or warp factor) in AdS geometry. For other attempts at solving the hierarchy problem with extra dimensions see .
The study of the cosmological evolution in this type of model has undergone several interesting turns. It was shown in that the Hubble parameter on the brane is proportional to the brane tension, rather than to the square root of it, as one would expect from pure low-energy four-dimensional reasoning. Two-brane models of this type have additional problems since the matter densities $`\rho _1`$, $`\rho _2`$ and pressures $`p_1`$ and $`p_2`$ need to be correlated in order to achieve a consistent solution (in fact one would derive that $`\rho _2\rho _1`$ which lacks a consistent interpretation). This also leads to the correlation between the states of matter on the two different branes. These problems with cosmology are indicative that some crucial ingredients are missing in the brane models which would allow a smooth transition to “normal” four-dimensional physics. It was conjectured in and further proven in that the missing element in brane models is the stabilization of the extra dimension.
It was shown in Refs. that single-brane models with a compactified extra dimension lead to standard cosmology, irrespective of the nature of the precise mechanism responsible for the stabilization of the transverse direction. In technical terms, stabilization provides a bulk value for the transverse pressure-like component of the stress-energy tensor, proportional to the trace of the stress-energy tensor of the brane. By including the bulk component of energy-momentum tensor ($`T_{55}`$), one not only recovers the standard form of the Friedmann equation, but also relieves the necessity of the (anti) correlation of the energy density on the brane allowing for the construction of realistic two- and one-brane models. In , it was shown that $`T_{55}`$ is automatically generated as a back-reaction to matter on the brane and induces a shift (proportional to $`T_\mu ^\mu `$) in the minimum of the dilaton potential. However, in such models with matter, the warp factor is insignificant (it is proportional to $`\rho /M_P^4`$) and we lose the potential solution to the hierarchy problem.
In our subsequent paper , we generalized this construction and found a static solution for the case of two positive tension branes, assuming the same stabilization mechanism. An interesting feature of this construction, noted in , is the possibility of the exponentially large ratio between the two scales on two different branes, which allows to address the gauge hierarchy problem a là Randall and Sundrum (RS), and at the same time avoid the presence of the negative tension brane. Two positive tension branes are allowed in this context because the solution to the warp factor is $`cosh`$-like rather than a pure exponential. Even though the negative tension branes are plausible in string theory, their existence in the real physical world is questionable. Indeed, these branes may turn out to be unstable, once the supersymmetry is broken, and to the best of our knowledge this question has not been properly studied so far. In contrast, positive tension branes are ubiquitous, and occur not only in string theory, but also in an ordinary field theories where they can simply be interpreted as domain walls, domain wall junctions, etc.
We believe that this idea deserves further consideration, and in this paper we take it one step further, introducing matter density and pressure on the two branes in order to obtain consistent cosmological solutions. In accordance with what we learned in the case of the single-brane models, this expansion is exactly of the Friedmann form, provided that the transverse dimension is stabilized. At the same time, there is no correlation between $`\rho _1`$, $`\rho _2`$, $`p_1`$ and $`p_2`$, and no restriction on the equations of state. We obtain general solutions for arbitrary brane tensions and we further demonstrate that it is possible to recover a solution to the hierarchy problem for negative or positive values of the observable brane tension. As a result, the general system of two branes discussed in this paper gives rise to cosmological solutions that include, as special cases, both the positive-negative brane combination of the RS model and the more natural case of the two positive tension branes . We find that the system is characterized by the presence of three parameters: the ratio of the self-energy over the bulk cosmological constant of each brane and the inter-brane distance. Demanding the cancellation of the effective cosmological constant, we place one constraint on a combination of these parameters while the resolution of the hierarchy problem fixes one more parameter. We keep the inter-brane distance, $`L`$, as the free parameter and we demonstrate that suitably small values of the warp factor, which we show to be directly related to the ratio $`M_P/M_W`$, is achieved for any value of $`L`$. In other words, no other fine tuning is necessary to obtain the necessary warp factor which we show to be directly related to the ratio of the lapse functions on the two branes.
As indicated above, we find a normal Friedmann expansion on the visible brane and show that in fact both the visible and hidden branes expand with the same expansion rate. Upon rescaling the matter density on the visible brane, we further show that the deduced Plank scale is the same on both branes. Although the scale factors and lapse functions change in time leading to a time dependent ratio of weak scale and four-dimensional Plank constant, we show that the residual change in $`M_W/M_P`$ is minute, suppressed by the ratios of $`\rho _i/M_P^4`$, and thus representing no danger to this model.
This paper is organized as follows. In section 2, we introduce the ansatz for the metric, stress-energy tensor and present the five-dimensional Einstein’s equations. In the next section, we derive the solution for the spatial scale factor, in the bulk, consistent with the brane boundary conditions, and show that the cosmological expansion rate on the two branes is of the form of the standard four-dimensional Friedmann equation, once the inter-brane distance is stabilized. In section 4, we consider the solution to the hierarchy problem with branes with arbitrary tensions, and define the necessary conditions for this to happen. We demonstrate that the hierarchy problem can be resolved for arbitrary values of the inter-brane separation, by appropriately choosing the value of the observable brane tension, with the RS choice of parameters arising only as a special case. We draw our conclusions in section 5.
## 2 The Field Equations
The starting point of our analysis will be the assumption of the existence of an extra dimension, denoted by the coordinate $`y`$, in addition to the usual four coordinates, $`\{t,x^i\}`$, of the 4-dimensional spacetime. We consider the following ansatz for the line-element of the 5-dimensional manifold
$$ds^2=n^2(t,y)dt^2+a^2(t,y)\delta _{ij}dx^idx^j+b^2(t,y)dy^2,$$
(2.1)
where $`a(t,y)`$ and $`b(t,y)`$ are the scale factors of the 4-dimensional spacetime (which we assume is spatially flat, i.e., $`{}_{}{}^{3}k=0`$) and the extra dimension, respectively, and $`n(t,y)`$ the lapse function that defines the time variable. Motivated by M-theory , we assume that the extra dimension has the topology of an orbifold, i.e. is a circle where the discrete $`Z_2`$ symmetry, $`yy`$, has been imposed. Two 3-branes of zero thickness located at the orbifold’s fixed points, $`y=0`$ and $`y=L`$, set the size of the extra dimension to be equal to $`2bL`$. The inter-brane distance is assumed to remain fixed due to a stabilization mechanism which ensures that $`\dot{b}=b^{}=0`$. The coordinate $`y`$ is scaled so that we may set, for simplicity, $`b=1`$. Note that, throughout the paper, dots and primes denote differentiation with respect to time and $`y`$, respectively. For other approaches to brane-world cosmologies, see .
We now turn to the 5-dimensional action that describes the coupling of the matter content of the universe with gravity. We assume that all the usual matter fields are localized on the two 3-branes, at $`y=0`$ and $`y=L`$, that play the role of a hidden and observable universe, respectively. The two branes are characterized by self-energies $`\mathrm{\Lambda }_1`$ and $`\mathrm{\Lambda }_2`$, while a non-vanishing cosmological constant $`\mathrm{\Lambda }_B`$ exists in the bulk. Under the above assumption, the action takes the form
$$S=d^4x𝑑y\sqrt{\widehat{g}}\left\{\frac{M_5^3}{16\pi }\widehat{R}+\mathrm{\Lambda }_B+\mathrm{\Lambda }_1\delta (y)+\mathrm{\Lambda }_2\delta (yL)+\widehat{}_o\right\}.$$
(2.2)
In the above, $`\widehat{}_o`$ represents all possible contributions to the action which are not strictly gravitational, $`M_5`$ is the fundamental 5-dimensional Planck mass and the hat denotes 5-dimensional quantities.
Einstein’s equations, with the above spacetime background (2.1) and the assumption that $`\dot{b}=b^{}=0`$, take the simplified form (for the full version of these equations see e.g. Refs. )
$`\widehat{G}_{00}=3\{({\displaystyle \frac{\dot{a}}{a}})^2n^2[{\displaystyle \frac{a^{\prime \prime }}{a}}+\left({\displaystyle \frac{a^{}}{a}}\right)^2]\}=\widehat{\kappa }^2\widehat{T}_{00},`$ (2.3)
$`\widehat{G}_{ii}=a^2\left\{{\displaystyle \frac{a^{}}{a}}\left({\displaystyle \frac{a^{}}{a}}+2{\displaystyle \frac{n^{}}{n}}\right)+2{\displaystyle \frac{a^{\prime \prime }}{a}}+{\displaystyle \frac{n^{\prime \prime }}{n}}\right\}+{\displaystyle \frac{a^2}{n^2}}\left\{{\displaystyle \frac{\dot{a}}{a}}\left({\displaystyle \frac{\dot{a}}{a}}+2{\displaystyle \frac{\dot{n}}{n}}\right)2{\displaystyle \frac{\ddot{a}}{a}}\right\}=\widehat{\kappa }^2\widehat{T}_{ii},`$ (2.4)
$`\widehat{G}_{05}=3\left({\displaystyle \frac{n^{}}{n}}{\displaystyle \frac{\dot{a}}{a}}{\displaystyle \frac{\dot{a}^{}}{a}}\right)=0,`$ (2.5)
$`\widehat{G}_{55}=3\left\{{\displaystyle \frac{a^{}}{a}}\left({\displaystyle \frac{a^{}}{a}}+{\displaystyle \frac{n^{}}{n}}\right){\displaystyle \frac{1}{n^2}}\left[{\displaystyle \frac{\dot{a}}{a}}\left({\displaystyle \frac{\dot{a}}{a}}{\displaystyle \frac{\dot{n}}{n}}\right)+{\displaystyle \frac{\ddot{a}}{a}}\right]\right\}=\widehat{\kappa }^2\widehat{T}_{55},`$ (2.6)
where $`\widehat{\kappa }^2=8\pi \widehat{G}=8\pi /M_5^3`$. In the above, $`\widehat{T}_{MN}`$ is the total 5-dimensional energy-momentum tensor which can be decomposed in terms of the bulk and brane contributions as follows
$$\widehat{T}_{MN}=\widehat{T}_{MN}^{(B)}+\widehat{T}_{MN}^{(1)}+\widehat{T}_{MN}^{(2)},$$
(2.7)
where
$`\widehat{T}_{(B)N}^M=\mathrm{diag}(\mathrm{\Lambda }_B,\mathrm{\Lambda }_B,\mathrm{\Lambda }_B,\mathrm{\Lambda }_B,\widehat{T}_{(B)\mathrm{\hspace{0.17em}5}}^5),`$
$`\widehat{T}_{(1)N}^M=\mathrm{diag}[{\displaystyle \frac{\delta (y)}{b}}(\rho _1\mathrm{\Lambda }_1,p_1\mathrm{\Lambda }_1,p_1\mathrm{\Lambda }_1,p_1\mathrm{\Lambda }_1),0],`$
$`\widehat{T}_{(2)N}^M=\mathrm{diag}[{\displaystyle \frac{\delta (yL)}{b}}(\rho _2\mathrm{\Lambda }_2,p_2\mathrm{\Lambda }_2,p_2\mathrm{\Lambda }_2,p_2\mathrm{\Lambda }_2),0].`$
Note that we allow for a non-vanishing value of the (55)-component of the energy-momentum tensor in the bulk that includes, in addition to the contribution due to the bulk cosmological constant, a part proportional to the trace of the energy-momentum tensor on the brane. As it has been shown , this component results from a 5-dimensional stabilization mechanism which is responsible for keeping the inter-brane distance, and thus the size of the extra dimension, fixed. The presence of a stabilizing potential for the radion field in the framework of the 5-dimensional, fundamental theory leads to a $`(t,y)`$-dependent value for $`\widehat{T}_{(B)\mathrm{\hspace{0.17em}5}}^5`$ distinctly different from $`\mathrm{\Lambda }_B`$. In particular, it can be shown that $`\widehat{T}_{(B)\mathrm{\hspace{0.17em}5}}^5+\mathrm{\Lambda }_B`$ is (to leading order in $`\rho `$) proportional to the trace $`T_\mu ^\mu `$ on the brane. For later use, we display here the relations that follow from the conservation of the energy-momentum tensor, $`D_M\widehat{T}_N^M=0`$, on the branes
$$\frac{d\rho _1}{dt}+3(\rho _1+p_1)\frac{\dot{a}_0}{a_0}=0,\frac{d\rho _2}{dt}+3(\rho _2+p_2)\frac{\dot{a}_L}{a_L}=0,$$
(2.9)
and in the bulk
$$\left(\widehat{T}_{(B)\mathrm{\hspace{0.17em}5}}^5\right)^{^{}}+\widehat{T}_{(B)\mathrm{\hspace{0.17em}5}}^5\left(\frac{n^{}}{n}+3\frac{a^{}}{a}\right)+\mathrm{\Lambda }_B\left(\frac{n^{}}{n}+3\frac{a^{}}{a}\right)=0.$$
(2.10)
Here, we closely follow the analysis of Ref. which we extend to the case of two branes with arbitrary self-energies $`\mathrm{\Lambda }_i`$’s. Our subsequent analysis will be greatly simplified by the fact that the (05)-component of Einstein’s equations (2.5) can be easily integrated to give the result
$$n(t,y)=\lambda (t)\dot{a}(t,y).$$
(2.11)
Using the normalization $`n(t,y=0)=1`$, the arbitrary function of time $`\lambda (t)`$ turns out to be $`\lambda (t)=1/\dot{a}_0`$ in terms of which, as we will shortly see, the Hubble parameter on both branes can be easily expressed.
## 3 Cosmological Evolution
In this section, we focus on the general solution for the scale factor $`a(t,y)`$ which is consistent with the aforementioned energy distribution in the 5-dimensional universe and the boundary conditions that the two, infinitely thin, brane-universes introduce in the theory. We will also determine the form of the generalized Friedmann equations that the boundary quantities $`a_0`$ and $`a_L`$ satisfy and we will investigate the conditions necessary for the restoration of the usual Friedmann equation on the two branes, through the vanishing of the effective cosmological constant, as well as for the successful stabilization of the extra dimension.
We note that by using the relation (2.11), the $`(00)`$-component of Einstein’s equations (2.3), in the bulk, reduces to an ordinary second-order differential equation for $`a(t,y)`$ with respect to $`y`$. As a result, we can easily find the general solution for the scale factor, in the bulk, which, for $`\mathrm{\Lambda }_B<0`$, has the form
$$a^2(t,y)=d_1(t)\mathrm{cosh}(A|y|)+d_2(t)\mathrm{sinh}(A|y|)\frac{B^2(t)}{A^2},$$
(3.12)
where
$$A^2=\frac{2\widehat{\kappa }^2}{3}|\mathrm{\Lambda }_B|,B^2(t)=\frac{2}{\lambda ^2(t)},$$
(3.13)
and where $`d_1`$ and $`d_2`$ are two functions of time which will be shortly determined. The above solution, derived in the bulk, needs to be smoothly connected to its boundary values $`a_0`$ and $`a_L`$. We can write
$$(𝐀):a^2(t,0)=a_0^2(t),a^2(t,L)=a_L^2(t).$$
(3.14)
The first of these conditions leads to the determination of the unknown function $`d_1(t)`$,
$$d_1(t)=a_0^2(t)+\frac{B^2(t)}{A^2}$$
(3.15)
while the second leads to the relation between the two boundary values $`a_0`$ and $`a_L`$
$$a_L^2(t)=d_1(t)\mathrm{cosh}(AL)+d_2(t)\mathrm{sinh}(AL)\frac{B^2(t)}{A^2}.$$
(3.16)
Two more sets of boundary conditions need to be satisfied by the general solution (3.12). These conditions follow from the inhomogeneity in the distribution of matter in the universe and involve the jumps in the first derivatives of $`a`$ and $`n`$ across the two branes. They have the form
$`(𝐁):{\displaystyle \frac{[a^{}]_0}{a_0}}={\displaystyle \frac{\widehat{\kappa }^2}{3}}(\rho _1+\mathrm{\Lambda }_1),{\displaystyle \frac{[a^{}]_L}{a_L}}={\displaystyle \frac{\widehat{\kappa }^2}{3}}(\rho _2+\mathrm{\Lambda }_2),`$ (3.17)
$`(𝐂):{\displaystyle \frac{[n^{}]_0}{n_0}}={\displaystyle \frac{\widehat{\kappa }^2}{3}}(3p_1+2\rho _1\mathrm{\Lambda }_1),{\displaystyle \frac{[n^{}]_L}{n_L}}={\displaystyle \frac{\widehat{\kappa }^2}{3}}(3p_2+2\rho _2\mathrm{\Lambda }_2).`$ (3.18)
We first focus on set (B) : the condition at $`y=0`$ determines the remaining arbitrary function $`d_2(t)`$ to be
$$d_2(t)=\frac{\widehat{\kappa }^2}{3A}a_0^2(\rho _1+\mathrm{\Lambda }_1),$$
(3.19)
while the same condition at $`y=L`$, combined with the relation (3.16), leads to two extremely interesting results: the Friedmann equation at $`y=0`$,
$$\left(\frac{\dot{a}_0}{a_0}\right)^2=\frac{\widehat{\kappa }^2|\mathrm{\Lambda }_B|}{3}\frac{\left[\mathrm{\Lambda }_{RS}+\frac{\left(\rho _1+\mathrm{\Lambda }_1+\rho _2+\mathrm{\Lambda }_2\right)}{\mathrm{tanh}\left(AL\right)}\frac{\left(\rho _1+\mathrm{\Lambda }_1\right)\left(\rho _2+\mathrm{\Lambda }_2\right)}{\mathrm{\Lambda }_{RS}}\right]}{\mathrm{\Lambda }_{RS}(\rho _2+\mathrm{\Lambda }_2)\mathrm{tanh}(\frac{AL}{2})},$$
(3.20)
and the ratio of the scale factors on the two branes
$$\frac{a_L^2}{a_0^2}=\frac{\mathrm{\Lambda }_{RS}(\rho _1+\mathrm{\Lambda }_1)\mathrm{tanh}(\frac{AL}{2})}{\mathrm{\Lambda }_{RS}(\rho _2+\mathrm{\Lambda }_2)\mathrm{tanh}(\frac{AL}{2})},$$
(3.21)
expressed in terms of $`\mathrm{\Lambda }_{RS}\sqrt{6|\mathrm{\Lambda }_B|/\widehat{\kappa }^2}`$, $`\mathrm{\Lambda }_i`$, $`\rho _i`$ and $`L`$. The above form of the Friedmann equation has also been derived in Ref. while some of the expressions presented in this section have been determined, in the special case $`\mathrm{\Lambda }_1=\mathrm{\Lambda }_2=\mathrm{\Lambda }_{RS}`$, in Ref. . As in the case of static solutions , the ratio of the scale factors on the two branes is given in terms of the ‘detuning’ of the total energy densities of the branes from the limiting value $`\mathrm{\Lambda }_{RS}`$. Differentiating the above expression with respect to time, we are able to derive the form of the Hubble parameter on the second brane, which is
$$\left(\frac{\dot{a}_L}{a_L}\right)=\left(\frac{\dot{a}_0}{a_0}\right)\left[\frac{\mathrm{\Lambda }_{RS}(\rho _2+\mathrm{\Lambda }_2)\mathrm{tanh}(\frac{AL}{2})}{\mathrm{\Lambda }_{RS}(\rho _1+\mathrm{\Lambda }_1)\mathrm{tanh}(\frac{AL}{2})}\right]\left[\frac{2\mathrm{\Lambda }_{RS}(2\mathrm{\Lambda }_1\rho _13p_1)\mathrm{tanh}(\frac{AL}{2})}{2\mathrm{\Lambda }_{RS}(2\mathrm{\Lambda }_2\rho _23p_2)\mathrm{tanh}(\frac{AL}{2})}\right].$$
(3.22)
Finally, let us note that the set (C) of boundary conditions involving the jumps of the lapse function $`n`$ across the two branes must also be satisfied. By using the fact that $`n(t,y)=\dot{a}(t,y)/\dot{a}_0(t)`$, differentiating with respect to $`y`$ and substituting in eqs. (3.18), we may see, after some algebra, that the jump conditions on both boundaries are identically satisfied. The above result reveals the absence of any correlation between the equations of state on the two brane-universes. This feature was also pointed out in Ref. , in the special case $`\mathrm{\Lambda }_1=\mathrm{\Lambda }_2=\mathrm{\Lambda }_{RS}`$, however, it turns out to be valid even in the case where the self-energies of the branes and the bulk cosmological constant are completely arbitrary and uncorrelated.
We now turn to the conditions that we need to impose in order to recover the usual Friedmann equation on both branes. The most obvious is the vanishing of the effective cosmological constant. From eq. (3.20), this translates to
$$\mathrm{\Lambda }_{eff}\mathrm{\Lambda }_{RS}\left(1\alpha \beta +\frac{\alpha +\beta }{\mathrm{tanh}(AL)}\right)=0,$$
(3.23)
where we have defined
$$\alpha \frac{\mathrm{\Lambda }_1}{\mathrm{\Lambda }_{RS}},\beta \frac{\mathrm{\Lambda }_2}{\mathrm{\Lambda }_{RS}}.$$
(3.24)
The condition (3.23) involves three parameters, $`\alpha `$, $`\beta `$ and $`L`$ which, until now, were completely uncorrelated. The condition of the vanishing of $`\mathrm{\Lambda }_{eff}`$ renders one of them a dependent parameter, thus, reducing the number of independent parameters to two. Note that there is a special choice for the parameters $`\alpha `$ and $`\beta `$ that eliminates $`\mathrm{\Lambda }_{eff}`$, for every value of the inter-brane distance $`L`$. This is $`\alpha =\beta =\pm 1`$ and corresponds to the Randall-Sundrum choice . For every other combination of $`\alpha `$ and $`\beta `$, the distance between the two branes has to be carefully chosen in order to ensure that $`\mathrm{\Lambda }_{eff}=0`$.
Although the vanishing of the effective cosmological constant has simplified the Friedmann-like equation at $`y=0`$, eq. (3.20), it has not completely restored its usual form. The existence of two branes with non-vanishing energy densities, $`\rho _i`$, leads to the appearance of both $`\rho _1`$ and $`\rho _2`$, in the equations that govern the cosmological evolution of each brane. Besides terms linear in $`\rho _1`$ and $`\rho _2`$, the resulting cosmological equations will be corrected by higher powers of energy densities, i.e. $`\rho _1^2`$, $`\rho _2^2`$, $`\rho _1\rho _2`$, and so on. Exact numerical coefficients in front of these terms are not known, as they depend on subtle details of radius stabilization . However, for a natural choice of parameters, these terms can easily be shown to be of secondary importance due to their extremely small magnitude. If we assume that $`\mathrm{\Lambda }_{RS}M_P^4`$ while $`\rho _i<\rho _c10^{123}M_P^4`$ (as dictated by the observable expansion rate of our 4D universe), and expand eq. (3.20) keeping only terms linear in $`\rho _i/\mathrm{\Lambda }_{RS}`$, we obtain the result
$$\left(\frac{\dot{a}_0}{a_0}\right)^2=\frac{\widehat{\kappa }^2|\mathrm{\Lambda }_B|}{3\mathrm{\Lambda }_{RS}\mathrm{tanh}(AL)}\left\{\frac{1\beta \mathrm{tanh}(AL)}{1\beta \mathrm{tanh}(\frac{AL}{2})}\rho _1+\frac{1\alpha \mathrm{tanh}(AL)}{1\beta \mathrm{tanh}(\frac{AL}{2})}\rho _2\right\}.$$
(3.25)
We could also use the condition of the vanishing of $`\mathrm{\Lambda }_{eff}`$ in order to eliminate one of the three parameters $`\alpha `$, $`\beta `$ and $`L`$. We have chosen to eliminate $`\alpha `$, using (3.23). The final form of the Friedmann equation, at $`y=0`$, then takes the form
$$\left(\frac{\dot{a}_0}{a_0}\right)^2=\frac{\widehat{\kappa }^2|\mathrm{\Lambda }_B|}{3\mathrm{\Lambda }_{RS}\mathrm{tanh}(AL)}\frac{1\beta \mathrm{tanh}(AL)}{1\beta \mathrm{tanh}(\frac{AL}{2})}\left\{\rho _1+\frac{\rho _2}{[\mathrm{cosh}(AL)\beta \mathrm{sinh}(AL)]^2}\right\}.$$
(3.26)
Before commenting on the above result, we should point out that the Friedmann equation on the observable brane, at $`y=L`$, has exactly the above form, in the same linear approximation. We can easily see, from eq. (3.22), that by ignoring terms of $`𝒪(\rho _i/\mathrm{\Lambda }_{RS})^2`$ or $`𝒪(\rho _i\rho _j/\mathrm{\Lambda }_{RS}^2)`$, the two Hubble parameters become identical. In eq. (3.26) we would restore the normal 4-dimensional Hubble expansion law if we define Newton’s constant as
$$\kappa ^2=\frac{\widehat{\kappa }^2|\mathrm{\Lambda }_B|}{\mathrm{\Lambda }_{RS}\mathrm{tanh}(AL)}\frac{1\beta \mathrm{tanh}(AL)}{1\beta \mathrm{tanh}(\frac{AL}{2})}.$$
(3.27)
The square of the Hubble parameter is now proportional to the total energy density of the universe, which contains contributions from both branes. Note that the constant coefficient in front of the energy density of our brane-universe, $`\rho _2`$, is not unity (this was first pointed out in Ref. ). As was discussed in , the observable energy density on our universe must be redefined in such a way as to absorb this coefficient. In the next section we will see that this redefinition is a result of the conformal transformation necessary to solve the hierarchy problem.
Finally, we need to address the problem of the stabilization of the extra dimension and determine the constraints on the energy distribution of the system that makes this stabilization possible. As we mentioned before and demonstrated in earlier works , the stabilization mechanism responsible for keeping fixed the inter-brane distance manifests itself through the existence of a non-vanishing (55)-component of the energy-momentum tensor. The (55)-component of Einstein’s equations, that we have ignored so far, will serve to determine the expression of this extra component. For the general solution (3.12), eq. (2.6) takes the form
$$\frac{d}{dt}\left[\frac{A^2}{4}(d_2^2d_1^2)+\frac{1}{A^2\lambda ^4}\right]=\frac{2\widehat{\kappa }^2}{3}a^3\dot{a}(\widehat{T}_{(B)\mathrm{\hspace{0.17em}5}}^5|\mathrm{\Lambda }_B|).$$
(3.28)
Substituting $`d_1`$, $`d_2`$ and $`\lambda `$ from eqs. (3.15), (3.19) and (3.20), respectively, we obtain the result
$`\widehat{T}_{(B)\mathrm{\hspace{0.17em}5}}^5`$ $`=`$ $`|\mathrm{\Lambda }_B|+{\displaystyle \frac{a_0^3\widehat{\kappa }^2}{12na^3}}\{(\rho _1+\mathrm{\Lambda }_1)(2\mathrm{\Lambda }_1\rho _13p_1)+\mathrm{\Lambda }_{RS}[2\mathrm{\Lambda }_{RS}{\displaystyle \frac{(4\mathrm{\Lambda }_1+\rho _13p_1)}{\mathrm{tanh}(AL)}}]`$ (3.29)
$``$ $`{\displaystyle \frac{\mathrm{\Lambda }_{RS}}{\mathrm{sinh}(AL)}}{\displaystyle \frac{a_L}{a_0}}[\mathrm{\hspace{0.17em}2}{\displaystyle \frac{a_L}{a_0}}(\rho _2+\mathrm{\Lambda }_2)+n_L(2\mathrm{\Lambda }_23p_2\rho _2)]\},`$
and to linear order in $`\rho _i`$
$`\widehat{T}_{(B)\mathrm{\hspace{0.17em}5}}^5`$ $`=`$ $`|\mathrm{\Lambda }_B|+{\displaystyle \frac{a_0^3|\mathrm{\Lambda }_B|}{2na^3}}\{{\displaystyle \frac{2(\beta ^21)}{\mathrm{cosh}^2(AL)[1\beta \mathrm{tanh}(AL)]^2}}`$ (3.30)
$``$ $`{\displaystyle \frac{(\rho _13p_1)[1\beta \mathrm{tanh}(\frac{AL}{2})]+(\rho _23p_2)[1+\beta \mathrm{tanh}(\frac{AL}{2})]}{\mathrm{\Lambda }_{RS}\mathrm{sinh}(AL)\mathrm{cosh}(AL)[1\beta \mathrm{tanh}(AL)]}}\}.`$
Note that this expression already assumes the vanishing of the cosmological constant (3.23). We may easily check that the expression (3.29), as well as the linearized version above, with $`\widehat{T}_{(B)\mathrm{\hspace{0.17em}5}}^5w(t)/n(t,y)a^3(t,y)`$, where $`w(t)`$ an arbitrary function of time, is consistent with the fifth component of the equation for the conservation of energy (2.10). In addition, as we expected, the coefficient $`w`$ is proportional to the trace of the energy momentum tensor on the branes . Note also that we neglect any small changes in $`b`$, due to changes in $`\widehat{T}_{(B)\mathrm{\hspace{0.17em}5}}^5`$ or $`\rho _i`$, since these are proportional to $`\rho _i^{3/2}`$ and we consistently ignore terms of this order throughout our analysis.
## 4 Resolution of the Hierarchy Problem
After having stabilized the extra dimension and derived the Friedmann equations that govern the cosmological evolution of each brane, we now focus our attention on the observable brane-universe. Our goal will be to determine the correct definition of mass scales, $`\stackrel{~}{m}`$, and energy density, $`\stackrel{~}{\rho }_2`$, as measured by a 4D observer in a FRW background, in terms of the corresponding quantities, $`m`$ and $`\rho _2`$, defined on our 4-brane which, however, is part of a curved, non-FRW spacetime. (In this context, by FRW we refer to a metric with lapse function $`n=1`$, with only a time dependent spatial scale factor.) To this end, we consider the theory of a scalar field, $`\mathrm{\Phi }`$, confined on our brane-universe which is embedded in a 5D spacetime background. Such a toy model was also considered in Ref. , however, here, we extend it to the case of a non-static spacetime background of the form (2.1). We consider the following action
$`S`$ $`=`$ $`{\displaystyle d^4x𝑑y\sqrt{\widehat{g}}\left[\widehat{g}^{\mu \nu }_\mu \mathrm{\Phi }_\nu \mathrm{\Phi }+\lambda (\mathrm{\Phi }^2v_0^2)^2\right]\delta (yL)}`$ (4.1)
$`=`$ $`{\displaystyle d^4x\sqrt{g_L}\left[g_L^{\mu \nu }_\mu \mathrm{\Phi }_\nu \mathrm{\Phi }+\lambda (\mathrm{\Phi }^2v_0^2)^2\right]},`$
where $`g_{L\mu \nu }=\widehat{g}_{\mu \nu }(y=L)`$, is the induced metric tensor on our brane. Using the ansatz (2.1), the above action takes the form
$`S`$ $`=`$ $`{\displaystyle }d^4x(a_L^3n_L)\{[{\displaystyle \frac{1}{n_L^2}}\dot{\mathrm{\Phi }}^2+{\displaystyle \frac{1}{a_L^2}}(\mathrm{\Phi })^2]+\lambda (\mathrm{\Phi }^2v_0^2)^2\}`$ (4.2)
$`=`$ $`{\displaystyle d^4x\sqrt{\stackrel{~}{g}}n_L^4\left[\frac{1}{n_L^2}\stackrel{~}{g}^{\mu \nu }_\mu \mathrm{\Phi }_\nu \mathrm{\Phi }+\lambda (\mathrm{\Phi }^2v_0^2)^2\right]},`$
where a conformal transformation
$$g_{L\mu \nu }\stackrel{~}{g}_{\mu \nu }=\frac{1}{n_L^2}g_{L\mu \nu }$$
(4.3)
has been performed that restores the FRW character of the 4D spacetime and brings the corresponding line-element to the form
$$d\stackrel{~}{s}^2=dt^2+\stackrel{~}{a}^2(t)\delta _{ij}dx^idx^j,$$
(4.4)
where $`\stackrel{~}{a}(t)a_L(t)/n_L(t)`$ is the observed scale factor of our universe. The last step we need to make involves the rescaling of the scalar field in order to absorb the gravitational coefficients that remain in the action. By setting $`\varphi \stackrel{~}{\varphi }=n_L\varphi `$, we arrive at
$$S=d^4x\sqrt{\stackrel{~}{g}}\left[\stackrel{~}{g}^{\mu \nu }_\mu \stackrel{~}{\mathrm{\Phi }}_\nu \stackrel{~}{\mathrm{\Phi }}+\lambda (\stackrel{~}{\mathrm{\Phi }}^2\stackrel{~}{v}_0^2)^2\right].$$
(4.5)
The $`vev`$ $`v_0`$ of the scalar field has also been rescaled, $`\stackrel{~}{v}_0=n_Lv_0`$, and as a result, the relation between any mass scale, $`\stackrel{~}{m}`$, defined in terms of the theory (4.5) and the corresponding one, $`m`$, defined in terms of (4.1), has the form
$$\stackrel{~}{m}=n_L(t)m.$$
(4.6)
We may, thus, conclude that the rescaling coefficient that relates any effective mass scale with the corresponding fundamental one, in a general spacetime background, is the lapse function of the original 5D line-element evaluated at the location of the brane. In the case of the RS model, $`n_L=e^{bkL}`$ and thus any mass scale $`\stackrel{~}{m}`$ is exponentially suppressed compared to $`m`$. Due to the fact that $`n_L=a_L`$, the induced line-element, in their case, was static since $`\stackrel{~}{a}(t)=a_L/n_L=1`$. In our case, however, the presence of the energy density $`\rho _2`$ inevitably leads to the expansion of our brane and to the time-dependence of the rescaling factor $`n_L`$ . In Ref. , it was argued that such a rescaling is not acceptable since the time-dependence of $`n_L`$ would modify the equations of motion of the scalar field $`\mathrm{\Phi }`$. Since $`\dot{\stackrel{~}{\mathrm{\Phi }}}=\dot{n}_L\mathrm{\Phi }`$, the above problem is resolved only if
$$\frac{\dot{\stackrel{~}{\mathrm{\Phi }}}}{\stackrel{~}{\mathrm{\Phi }}}=\frac{\dot{n}_L}{n_L}1.$$
(4.7)
Moreover, any significant time evolution of the conformal factor $`n_L`$ would lead to the time-dependence of either the fundamental mass scale, $`m`$, or the effective one, $`\stackrel{~}{m}`$.
It is straightforward to check that the above condition is indeed satisfied in our case. The expression for the lapse function at $`y=L`$ can be obtained using the relation $`n_L(t)=\lambda (t)\dot{a}_L(t)`$ and eq. (3.22), and is found to be
$$n_L(t)=\frac{a_0}{a_L}\frac{2\mathrm{\Lambda }_{RS}(2\mathrm{\Lambda }_13p_1\rho _1)\mathrm{tanh}(\frac{AL}{2})}{2\mathrm{\Lambda }_{RS}(2\mathrm{\Lambda }_23p_2\rho _2)\mathrm{tanh}(\frac{AL}{2})}.$$
(4.8)
Differentiating the above expression with respect to time, we find
$$\frac{\dot{n}_L}{n_L}=\frac{\dot{a}_0}{a_0}\frac{\dot{a}_L}{a_L}+\frac{(3w_1+1)\dot{\rho }_1\mathrm{tanh}(\frac{AL}{2})}{2\mathrm{\Lambda }_{RS}[1\alpha \mathrm{tanh}(\frac{AL}{2})]}\frac{(3w_2+1)\dot{\rho }_2\mathrm{tanh}(\frac{AL}{2})}{2\mathrm{\Lambda }_{RS}[1\beta \mathrm{tanh}(\frac{AL}{2})]},$$
(4.9)
where we have set $`p_i=w_i\rho _i`$ and ignored higher-order terms of order $`𝒪(\rho _i\dot{\rho }_i)`$. As we have shown in the previous section, the first two terms are identical if we ignore corrections of order $`𝒪(\rho _i/\mathrm{\Lambda }_{RS})^2`$, which are extremely small if we recall that $`\rho _i/\mathrm{\Lambda }_{RS}10^{123}`$. The remaining two terms can be easily shown to be proportional to $`\dot{\rho }_i/\mathrm{\Lambda }_{RS}(\rho _i/\mathrm{\Lambda }_{RS})^{3/2}M_P`$ which undoubtly proves the validity of the condition (4.7). It may also be worth checking the validity of this argument at the time of the Big-Bang Nucleosynthesis. The corresponding term in the expression of $`\dot{n}_L`$ would be $`10^{148}M_P(\rho _N/\rho _c)^{3/2}`$, where $`\rho _N`$ is the energy density of the universe during the BBN. If we assume that $`T_{BBN}10`$ MeV, we may easily find that $`\rho _N/\rho _c10^{39}`$. This, in turn, leads to the final result that the magnitude of the term $`\dot{\rho }_i/\mathrm{\Lambda }_{RS}`$ at the same period was of the order of $`10^{90}M_P`$ rendering the time-dependence of the lapse function $`n_L`$ completely unobservable even in the early universe.
If we assume that the fundamental energy scale of the 5D theory in eq. (4.6), $`m`$, is of order $`M_P`$, while $`\stackrel{~}{m}1`$ TeV, the rescaling coefficient should be of order $`10^{16}`$. Expanding the result for the lapse function (4.8) and keeping only terms linear in $`\rho _i/\mathrm{\Lambda }_{RS}`$, we obtain
$$n_L=\sqrt{\frac{1\alpha \mathrm{tanh}(\frac{AL}{2})}{1\beta \mathrm{tanh}(\frac{AL}{2})}}\left\{1+\frac{(3w_1+2)\rho _1\mathrm{tanh}(\frac{AL}{2})}{2\mathrm{\Lambda }_{RS}[1\alpha \mathrm{tanh}(\frac{AL}{2})]}\frac{(3w_2+2)\rho _2\mathrm{tanh}(\frac{AL}{2})}{2\mathrm{\Lambda }_{RS}[1\beta \mathrm{tanh}(\frac{AL}{2})]}\right\}.$$
(4.10)
Since, the last two terms inside the brackets are negligible compared to unity, the condition that $`n_L10^{16}`$ should be imposed on the prefactor appearing in (4.10) involving $`\alpha `$, $`\beta `$ and $`L`$. Note that this prefactor is nothing more than the ratio of the scale factors on the two branes (3.21), once the small energy density contributions are ignored. By making use of eq. (3.23), this condition can be conveniently rewritten as
$$n_L^2=\frac{1}{\mathrm{cosh}(AL)\beta \mathrm{sinh}(AL)}10^{32},$$
(4.11)
Once eq. (4.11) is imposed, we are left with only one independent parameter in the theory which we choose to be $`L`$, the inter-brane distance. The parameter $`\beta `$ can be determined by the condition (4.11), which together with the tuning of the cosmological constant (3.23) give the following expressions for $`\alpha `$ and $`\beta `$ in terms of $`L`$
$$\alpha =\frac{\mathrm{cosh}(AL)10^{32}}{\mathrm{sinh}(AL)},$$
(4.12)
and
$$\beta =\frac{\mathrm{cosh}(AL)10^{32}}{\mathrm{sinh}(AL)}.$$
(4.13)
According to our notation,
$$AL=\sqrt{\frac{2}{3}\frac{|\mathrm{\Lambda }_B|}{M_5^3}}LM_PL,$$
(4.14)
if we assume that $`|\mathrm{\Lambda }_B|^{1/5}M_5M_P`$. In that case, the quantity $`AL`$ is nothing more than the distance between the two branes measured in units $`M_P^1`$. Since the inter-brane distance is the free parameter of the theory, it would be interesting to see how the dependent parameters $`\alpha `$ and $`\beta `$ change as we vary $`AL`$. Their values will be dictated by eqs. (4.12) and (4.13), respectively, if we demand that both of the cosmological and the hierarchy problems are simultaneously solved. Numerical results for $`\alpha `$ and $`\beta `$ are plotted in Figure 1 as a function of $`AL`$ and some specific numerical examples are given in Table I.
Table I : The parameters $`\beta `$ and $`1\alpha `$ for various values of the inter-brane distance $`L`$.
$`\begin{array}{ccccc}& & & & \\ & & & & \\ \mathrm{𝐀𝐋}& \mathrm{𝟏}& \mathrm{𝟐𝟎}& \mathrm{𝟒𝟎}& \mathrm{𝟔𝟎}\\ & & & & \\ \beta & 8.5092\times 10^{31}& 9.0799\times 10^{27}& 8.4967\times 10^{14}& 1.7513\times 10^6\\ & & & & \\ 1\alpha & 0.31304& 4.1223\times 10^9& 3.6097\times 10^{35}& 1.5335\times 10^{52}\\ & & & & \\ & & & & \\ \mathrm{𝐀𝐋}& \mathrm{𝟕𝟎}& \mathrm{𝟕𝟓}& \mathrm{𝟖𝟎}& \mathrm{𝟖𝟓}\\ & & & & \\ \beta & 78.509& 0.46427& 0.99639& 0.99998\\ & & & & \\ 1\alpha & 3.0814\times 10^{61}& 3.9223\times 10^{65}& 3.6032\times 10^{67}& 1.6388\times 10^{71}\\ & & & & \end{array}`$
We note that, for small values of the free parameter $`L`$, the brane parameter $`\alpha `$ is positive and of $`𝒪(1)`$ while $`\beta `$ is forced to acquire a value which is, not only negative, but also many orders of magnitude larger than that of $`\alpha `$. Such a situation is rather unphysical as it leads to an observable brane-universe with negative, total energy density and to an unexplained hierarchy between the sizes of the brane tensions. As $`L`$ increases, however, both $`\alpha `$ and $`\beta `$ start approaching unity: the combination $`1\alpha `$ becomes extremely small while $`\beta `$ is moving towards less negative values. At $`L73.682722\mathrm{}`$, the critical point $`\beta =1`$ and $`\alpha 1`$, the choice of brane parameters in the Randall-Sundrum<sup>1</sup><sup>1</sup>1Note, that our conventions differ from those in Ref. by a factor of 2. Thus, their warp factor $`e^{kr_c\pi }`$ corresponds to $`e^{AL/2}`$, in our case, if we ignore the matter on the brane and set $`\alpha =\beta =1`$. model , is reached. If we further increase $`L`$, the combination $`1\alpha `$ changes sign but it keeps approaching zero. On the other hand, $`\beta `$ becomes positive at a slightly larger value of the inter-brane distance, at $`AL75`$. Any further increase in the value of $`AL`$ forces both of the brane parameters, $`\alpha `$ and $`\beta `$, to unity, which is the largest value that these parameters may take on. This is clear from eq. (3.23), for $`\alpha =\beta =1`$, where the hyperbolic tangent takes its asymptotic value, i.e. $`\mathrm{tanh}(AL)=1`$, and corresponds to an infinite distance between the two branes. Once the two branes become isolated, the only solution to the cosmological constant problem follows for $`\alpha =\beta =1`$. This is obvious from eq. (3.20), which in the limit $`\mathrm{tanh}(AL)1`$, takes the form
$$\left(\frac{\dot{a}_0}{a_0}\right)^2=\frac{\widehat{\kappa }^2|\mathrm{\Lambda }_B|}{3}\left[1+\frac{(\rho _1+\mathrm{\Lambda }_1)}{\mathrm{\Lambda }_{RS}}\right].$$
(4.15)
In this limit, the expansion of the brane at $`y=0`$ becomes independent of the presence of the second brane at $`y=L`$ and the flatness of the corresponding 4D spacetime is guaranteed only for $`\alpha =1`$ (similar results hold for the observable universe at $`y=L`$). Thus, for any value of the inter-brane distance, larger than, approximately, $`75M_P^1`$, the brane parameters are always positive, with their upper limit being unity, a result which leads to two brane-universes with positive, total energy densities. Clearly, if $`AL75`$, then, in addition to tuning $`1\alpha `$ to zero, $`\beta `$ must also be tuned to unity. We should stress here, once again, that all of the above values of $`\alpha `$ and $`\beta `$, either positive or negative, are in accordance with both the vanishing of the effective cosmological constant on both branes and with the resolution of the hierarchy problem on the observable universe.
Let us finally address the problem of the definition of the energy density, $`\stackrel{~}{\rho }_2`$, on our brane as measured by an observer living in a FRW spacetime in terms of the original energy density, $`\rho _2`$, defined in a 4D sub-space with lapse function $`n_L1`$. The energy-momentum tensor associated with the matter content of our universe is defined as
$$T_{\mu \nu }^M=\frac{2}{\sqrt{g_L}}\frac{\delta S^M}{\delta g_L^{\mu \nu }},$$
(4.16)
where $`S_M`$ is the part of the original action (2.2) defined only on our brane and containing time-dependent functions that describe the distribution of matter sources in our universe. After the conformal transformation (4.3), $`T_{\mu \nu }^M`$ changes as
$$T_{\mu \nu }^M\stackrel{~}{T}_{\mu \nu }^M=\frac{2}{\sqrt{\stackrel{~}{g}}}\frac{\delta S^M}{\delta \stackrel{~}{g}^{\mu \nu }}=n_L^2T_{\mu \nu }^M.$$
(4.17)
By using the fact that $`\stackrel{~}{T}_\nu ^\mu =(\stackrel{~}{\rho }_2,\stackrel{~}{p}_2,\stackrel{~}{p}_2,\stackrel{~}{p}_2)`$ and $`\stackrel{~}{T}_\nu ^\mu =\stackrel{~}{g}^{\mu \rho }\stackrel{~}{T}_{\rho \nu }=n_L^4T_\nu ^\mu `$, we finally obtain the results
$$\stackrel{~}{\rho }_2=n_L^4\rho _2,\stackrel{~}{p}_2=n_L^4p_2.$$
(4.18)
The above definition of the observed energy density $`\stackrel{~}{\rho }_2`$ is in perfect agreement with the form of the Friedmann equation (3.26) for the cosmological expansion of our brane-universe. Recall that, to linear order in $`\rho _i`$, the expansion rate is
$$\left(\frac{\dot{a}_0}{a_0}\right)^2(\rho _1+n_L^4\rho _2)=(\rho _1+\stackrel{~}{\rho }_2),$$
(4.19)
when one uses the expression for $`n_L`$ (4.11) in terms of the brane parameter $`\beta `$ and the inter-brane distance $`AL`$ (here, we consider only the lowest-order expression of $`n_L`$ since any $`\rho _i`$-dependent terms would create higher-order corrections, of $`𝒪(\rho _i^2,\rho _i\rho _j)`$, in the Friedmann equation). Note that no redefinition is necessary for the components of the 4D energy-momentum tensor on the brane at $`y=0`$ since the lapse function, and thus the conformal factor, $`n_0`$, is unity. Similarly, we have $`\stackrel{~}{a}_0=a_0`$, and, as a result, eq. (4.19) still governs the expansion of the 4D FRW spacetime at $`y=0`$. Although the rescaling of the scale factor at $`y=L`$ is not trivial, we can easily see that
$$\frac{\dot{\stackrel{~}{a}}_L}{\stackrel{~}{a}_L}=\frac{\dot{a}_L}{a_L}\frac{\dot{n}_L}{n_L}=\frac{\dot{a}_0}{a_0}+𝒪(\frac{\rho _i^2}{\mathrm{\Lambda }_{RS}^2},\frac{\rho _i\rho _j}{\mathrm{\Lambda }_{RS}^2}).$$
(4.20)
Therefore, in the linear-order approximation in the energy densities, the same equation (3.26) describes the expansion of the 4D FRW spacetime at $`y=L`$. After performing the redefinition (4.18), the unique Friedmann equation (4.19) reveals the fact that the cosmological evolution of both branes is governed by the sum of the energy densities of the two branes as measured by observers living in flat, FRW 4D spacetimes. The same conclusion regarding the Friedmann expansion of both branes was also derived in Ref. in the special case of the RS choice of the brane parameters $`\alpha `$ and $`\beta `$.
## 5 Conclusions
Before summarizing the work in this paper, let us emphasize briefly the main features of the two-brane models which may naturally explain the large ratio $`M_P/M_W`$ via the hierarchy of the two length scales on different branes. In the RS model , the solution for the warp factor in the bulk, $`a(y)`$, was given by an exponential function of $`y`$ that interpolated between a positive and a negative self-energy brane. In the two-positive-brane model , the corresponding solution behaved like a $`\mathrm{cosh}`$-function characterized by the existence of a minimum. The solution to the hierarchy problem may be explained by the position of these two positive branes with respect to the minimum. Indeed, as we have shown, the ratio of the two lapse functions can be large, when the branes are located at non-equal distances from the minimum.
In this work, we have generalized these static models by introducing energy densities on the branes. Moreover, we have allowed for the possibility of arbitrary (positive or negative) brane tensions. We have shown that models of this type exhibit “normal” cosmology after demanding the vanishing of the effective cosmological constant on both branes and redefining the energy density on the observable brane in order to restore the FRW character of our 4D universe. It is also worth noting that no correlation between states of matter on the respective branes is required. In addition, we have shown that despite the time-dependence of the lapse function on the observable brane, that defines the ratio $`M_P/M_W`$, its time evolution does not lead to any measurable change of the effective four-dimensional parameters since such changes would be proportional to $`\rho _i/M_P^4`$.
A crucial element in our model is the mechanism, responsible for the stabilization of the inter-brane distance. Here we simply assume that this mechanism exists without specifying physical causes, leading to the radion/dilaton stabilization. An extended discussion of this issue can be found in Ref. . Of course, any model with a string like dilaton must ensure such a stabilization, since the dilaton expectation value fixes the gauge coupling and mass scales in the standard model.
We can not claim that two-brane models of this type are devoid of fine-tunings. As in all other known models, the cosmological constant requires a fine-tuning of different parameters in the model to ensure $`\mathrm{\Lambda }_{eff}=0`$. This reduces the initial set of three parameters (the two brane tensions relative to the bulk cosmological constant, and the inter-brane separation) to two. Note that the bulk cosmological constant relates the 4D and 5D (fundamental) Planck scales. When $`\mathrm{\Lambda }_{eff}`$ is set to zero, the weak scale/Plank scale hierarchy problem may be explained by fixing a relation between the remaining brane tension and the inter-brane separation. For every value of the inter-brane distance, a solution to both the cosmological constant problems and the hierarchy problem can be derived by appropriately choosing the tension of the observable universe. While for small values of $`AL`$, the brane parameter $`\beta `$ is forced to acquire a large negative value, as the inter-brane distance increases, both of the parameters $`\alpha `$ and $`\beta `$ start decreasing (in absolute value) towards unity. The RS choice of parameters, i.e. $`\alpha =\beta =1`$, arises as a special solution when we reach the value $`AL73`$ while, for all separations $`AL>75`$, more natural solutions (with positive brane tensions for both branes) arise. The parameters $`\alpha `$ and $`\beta `$ remain positive no matter how large the inter-brane separation becomes. The asymptotic value $`\alpha =\beta =1`$ is reached in the limit of infinite $`AL`$, when the two branes become isolated. In this limit, the problem of the cosmological constant is solved by imposing independent conditions on each brane and the framework for the resolution of the hierarchy problem ceases to exist.
Acknowledgments This work was supported in part by the Department of Energy under Grant No. DE-FG-02-94-ER-40823 at the University of Minnesota.
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# 1 Introduction
## 1 Introduction
While a brane breaks half of the space-time supersymmetry, the anti-brane breaks precisely the other half of the supersymmetry. Thus, a system of a brane and anti-brane breaks together all the space-time supersymmetry. The system is not stable, however, since the brane and anti-brane attract each other. This can be understood as the appearance of a tachyon on the world-volume of the branes. It arises from the open string stretched between the brane and the anti-brane and it is charged under the world-volume gauge groups. The decay of the system can be seen by the tachyon rolling down to the minimum of its potential . The phenomenon of tachyon condensation is fairly well studied by now in the open string description . It would be interesting to ask how the phenomenon appears from the closed string viewpoint. One of the aims of this paper is to construct supergravity solutions that correspond to N D$`p`$-branes coinciding with $`\overline{N}`$ $`\overline{\mathrm{D}p}`$-branes (anti D-branes) and analyse the supergravity description of tachyon condensation.
While Type IIA (Type IIB) string theory has BPS D-branes of even (odd) dimensions, they also admit non-BPS D-branes of odd (even) dimensions. These branes are not stable. They have been interpreted as the string theoretical analogues of sphalerons in field theory . The families of supergravity solutions that we will discuss contain also backgrounds that correspond to these branes. Stable non-BPS brane configurations are much studied too . However, we will not discuss supergravity backgrounds that correspond to these objects.
Another motivation that we have for studying brane-antibrane solutions is to understand the relation between these solutions and the Schwarzschild black hole solution (see, e.g. for an early indication of such a connection in the context of five-dimensional black holes of Type IIB theory), which may have possible applications in the study of non-supersymmetric black holes.
This paper is organised as follows. In Section 2 we describe the supergravity solution that corresponds to N D$`p`$-branes coinciding with $`\overline{N}`$ $`\overline{\mathrm{D}p}`$-branes and its physical properties. In Section 3 we analyse the supergravity description of tachyon condensation. We will also discuss the issue of decoupling and open-closed string duality. In section 4 we describe a general family of supergravity solutions that includes non-Poincare-invariant world-volumes. In particular it contains an interpolation between the brane-antibrane solution and the Schwarzschild solution. We discuss the possible application to the study of non-supersymmetric black holes. Section 5 contains a short discussion of the results.
We note that supergravity descriptions of smeared brane-antibrane configurations have been presented in . We will discuss in this paper the localized ones. Unstable branes on AdS have been analysed in .
## 2 The Supergravity Description
In this section we will describe Type II supergravity solutions that correspond to $`N`$ D$`p`$-branes coincident with $`\overline{N}`$ $`\overline{\mathrm{D}p}`$-branes and their physical properties.
### 2.1 The Supergravity Solution
The strategy for constructing such solutions will be the following. We know that a brane-antibrane configuration must have the full world-volume Poincare symmetry $`ISO(p,1)`$<sup>1</sup><sup>1</sup>1By contrast, a non-extremal Dp-brane breaks $`ISO(p,1)ISO(p)`$, which is expected of a finite temperature world-volume field theory (see Section 4). Here $`I`$ stands for “inhomogeneous”, referring to the translational symmetries.. Furthermore, it should have rotational symmetry $`SO(9p)`$ in the $`9p`$ transverse directions. For $`N\overline{N}`$, the system will also carry an appropriate RR charge.
We therefore look for the most general solution of Type II A/B supergravity which possess the symmetry
$$𝒮=ISO(p,1)\times SO(9p),$$
(1)
and carries charge under a RR field<sup>2</sup><sup>2</sup>2Our convention for the RR field and potentials is as follows. For electric $`p`$-branes (i.e. for $`p=0,1,2`$), the RR field strength is $`F_{p+2}dC^{(p+1)}`$. For magnetic $`p`$-branes i.e. for $`p=4,5,6`$, we interpret $`C^{(p+1)}`$ as the dual potential, and the RR field-strength will be given by $`F_{8p}e^{\frac{(3p)\varphi }{2}}\left(dC^{(p+1)}\right)`$. For 3-branes ($`p=3`$) the self-dual field strength is given by $`F_5=\frac{1}{\sqrt{2}}(dC^{(4)}+dC^{(4)})`$..
The most general form of the metric, dilaton and RR-field consistent with the symmetry (1) is
$`ds^2`$ $`=`$ $`e^{2A(r)}dx_\mu dx^\mu +e^{2B(r)}\left(dr^2+r^2d\mathrm{\Omega }_{8p}^2\right),`$
$`\varphi `$ $`=`$ $`\varphi (r),`$
$`C^{(p+1)}`$ $`=`$ $`e^{\mathrm{\Lambda }(r)}dx^0dx^1\mathrm{}dx^p.`$ (2)
We look for solutions of the form (2), of Type II A/B supergravity Lagrangian, whose relevant part is given (in the Einstein frame) by
$$S=\frac{1}{16\pi G_N^{10}}d^{10}x\sqrt{g}\left(R\frac{1}{2}_M\varphi ^M\varphi \frac{1}{2n!}e^{a\varphi }F_n^2\right)$$
(3)
where $`a=\frac{5n}{2}`$. The relation between the rank $`n`$ of the RR field strength $`F_n`$ and the dimensionality $`p`$ of the brane has been explained in the footnote 2.
In (2) and in the rest of the paper we represent ten-dimensional coordinates by $`x^M,M=0,\mathrm{},9`$ and brane world-volume coordinates (including time) by $`x^\mu ,\mu =0,1,\mathrm{},p`$. We will denote the transverse coordinates by $`x^i,i=1,\mathrm{},9p`$ or, alternatively, by the polar coordinates $`r,\theta _1,\mathrm{},\theta _{8p}(r^2x^ix^i`$).
The equations of motion that follow from (3) for the ansatz (2) are (see, e.g.,)
$`A^{\prime \prime }+(p+1)(A^{})^2+(7p)A^{}B^{}+{\displaystyle \frac{8p}{r}}A^{}={\displaystyle \frac{7p}{16}}S^2,`$
$`B^{\prime \prime }+(p+1)A^{}B^{}+{\displaystyle \frac{p+1}{r}}A^{}+(7p)(B^{})^2+{\displaystyle \frac{152p}{r}}B^{}={\displaystyle \frac{1}{2}}{\displaystyle \frac{p+1}{8}}S^2,`$
$`(p+1)A^{\prime \prime }+(8p)B^{\prime \prime }+(p+1)(A^{})^2+{\displaystyle \frac{8p}{r}}B^{}(p+1)A^{}B^{}+{\displaystyle \frac{1}{2}}(\varphi ^{})^2={\displaystyle \frac{1}{2}}{\displaystyle \frac{7p}{8}}S^2,`$
$`\varphi ^{\prime \prime }+\left((p+1)A^{}+(7p)B^{}+{\displaystyle \frac{8p}{r}}\right)\varphi ^{}={\displaystyle \frac{a}{2}}S^2,`$
$`\left(\mathrm{\Lambda }^{}e^{\mathrm{\Lambda }+a\varphi (p+1)A+(7p)B}r^{8p}\right)^{}=0,`$ (4)
where
$$S=\mathrm{\Lambda }^{}e^{\frac{1}{2}a\varphi +\mathrm{\Lambda }(p+1)A}.$$
(5)
The mathematical solution to this system of differential equations has already been presented in . The solutions depend on three<sup>3</sup><sup>3</sup>3One would naively predict two parameters, corresponding to the mass and the charge of the system, based on a suitable generalisation of Birkhoff’s theorem. However, the proofs of such uniqueness theorems assume regular manifolds and therefore do not apply here. parameters $`r_0,c_1,c_2`$ (we have relabelled $`c_3`$ of as $`c_2`$, and $`k`$ as $`k`$) and are given by
$`A(r)`$ $`=`$ $`{\displaystyle \frac{(7p)(3p)c_1}{64}}h(r)`$
$`{\displaystyle \frac{7p}{16}}\mathrm{ln}\left[\mathrm{cosh}(kh(r))c_2\mathrm{sinh}(kh(r))\right],`$
$`B(r)`$ $`=`$ $`{\displaystyle \frac{1}{7p}}\mathrm{ln}\left[f_{}(r)f_+(r)\right]+{\displaystyle \frac{(p3)(p+1)c_1}{64}}h(r)`$
$`+{\displaystyle \frac{p+1}{16}}\mathrm{ln}\left[\mathrm{cosh}(kh(r))c_2\mathrm{sinh}(kh(r))\right],`$
$`\varphi (r)`$ $`=`$ $`{\displaystyle \frac{(7p)(p+1)c_1}{16}}h(r)`$
$`+{\displaystyle \frac{3p}{4}}\mathrm{ln}\left[\mathrm{cosh}(kh(r))c_2\mathrm{sinh}(kh(r))\right],`$
$`e^{\mathrm{\Lambda }(r)}`$ $`=`$ $`\eta (c_2^21)^{1/2}{\displaystyle \frac{\mathrm{sinh}(kh(r))}{\mathrm{cosh}(kh(r))c_2\mathrm{sinh}(kh(r))}},`$ (6)
where
$`f_\pm (r)`$ $``$ $`1\pm \left({\displaystyle \frac{r_0}{r}}\right)^{7p},`$
$`h(r)`$ $`=`$ $`\mathrm{ln}\left[{\displaystyle \frac{f_{}(r)}{f_+(r)}}\right],`$
$`k`$ $`=`$ $`\pm \sqrt{{\displaystyle \frac{2(8p)}{7p}}{\displaystyle \frac{(p+1)(7p)}{16}}c_1^2},`$
$`\eta `$ $`=`$ $`\pm 1.`$ (7)
The parameter $`\eta `$ describes whether we are measuring the “brane” charge or the “antibrane” charge of the system.
The parameters $`(r_0,c_1,c_2)`$ appear as integration constants and as such they could be complex, describing a six-dimensional space. However, the reality of the supergravity fields singles out three distinct three-dimensional subspaces I, II and III, as discussed in appendix A. For the rest of our paper, we will concentrate on the physical properties of the solution I where the above three parameters are all real; we will comment on II and III in appendix A. We also note that besides the three continuous parameters $`r_0,c_1`$ and $`c_2`$, our solution has two additional discrete parameters: sign$`(k),\eta `$.
The solution is invariant under three independent $`Z_2`$ transformations which act on the space of the parameters
$$\begin{array}{c}(\mu ,c_1,c_2,\mathrm{sign}(k),\eta )(\mu ,c_1,c_2,\mathrm{sign}(k),\eta )\hfill \\ (\mu ,c_1,c_2,\mathrm{sign}(k),\eta )(\mu ,c_1,c_2,\mathrm{sign}(k),\eta )\hfill \\ (\mu ,c_1,c_2,\mathrm{sign}(k),\eta )(\mu ,c_1,c_2,\mathrm{sign}(k),\eta )\hfill \\ \mu r_0^{7p}.\hfill \end{array}$$
(8)
For convenience we will fix the above $`Z_2`$’s by choosing
(a) the positive branch of the square root for $`k`$, namely
$$k=\sqrt{\frac{2(8p)}{7p}\frac{(p+1)(7p)}{16}c_1^2},$$
(9)
(b)
$$c_10.$$
(10)
The case of the instanton ($`p=1`$):
The solutions mentioned above also include $`p=1`$. In this case there is no $`A(r)`$; the metric, dilaton and the RR potential are explicitly given by
$`ds^2`$ $`=`$ $`\left({\displaystyle \frac{f_{}(r)}{f_+(r)}}\right)^{1/4}\left(dr^2+r^2d\mathrm{\Omega }_{8p}^2\right),`$
$`\varphi `$ $`=`$ $`\mathrm{ln}\left[\mathrm{cosh}({\displaystyle \frac{3}{2}}h(r))c_2\mathrm{sinh}({\displaystyle \frac{3}{2}}h(r))\right],`$
$`C^{(0)}`$ $`=`$ $`e^{\mathrm{\Lambda }(r)}=\eta (c_2^21)^{1/2}{\displaystyle \frac{\mathrm{sinh}(\frac{3}{2}h(r))}{\mathrm{cosh}(\frac{3}{2}h(r))c_2\mathrm{sinh}(\frac{3}{2}h(r))}},`$ (11)
where
$`f_\pm (r)`$ $`=`$ $`\left(1({\displaystyle \frac{r_0}{r}})^8\right),`$
$`h(r)`$ $`=`$ $`\mathrm{ln}\left[{\displaystyle \frac{f_{}(r)}{f_+(r)}}\right].`$ (12)
An interesting point to note is that in this case the solution depends only on two parameters $`r_0,c_2`$ (which are functions of mass and charge), consistent with Birkhoff’s theorem. The extra parameter $`c_1`$ does not appear. According to the interpretation in the next section it implies that there is no tachyon associated with this solution.
The neutral case (taken as $`c_2=1`$) is described by
$`ds^2`$ $`=`$ $`\left({\displaystyle \frac{f_{}(r)}{f_+(r)}}\right)^{1/4}\left(dr^2+r^2d\mathrm{\Omega }_{8p}^2\right),`$
$`\mathrm{exp}[\varphi ]`$ $`=`$ $`\left({\displaystyle \frac{f_{}(r)}{f_+(r)}}\right)^{3/2}.`$ (13)
Regarded as a IIB solution, this should be interpreted as a $`D(1)`$-$`\overline{D}(1)`$ pair. On the other hand, the same solution can alternatively be regarded as a IIA solution; in that case it has a natural lift to eleven dimensions, given by the formula
$$ds_{11}^2=\mathrm{exp}[4\varphi /3]dx_{10}^2+\mathrm{exp}[\varphi /6]ds^2.$$
(14)
It is easy to see that the eleven-dimensional metric becomes the Euclidean Schwarzschild metric
$$ds_{11}^2=\left(1\frac{M}{\stackrel{~}{r}^8}\right)dx_{10}^2+\frac{dr^2}{1\frac{M}{\stackrel{~}{r}^8}}+\stackrel{~}{r}^2d\mathrm{\Omega }_9^2$$
(15)
where $`M=4r_0^8`$ and $`\stackrel{~}{r}=rf_+^{1/4}`$. It has been pointed out in that this metric describes the non-BPS D-instanton of type IIA <sup>4</sup><sup>4</sup>4We thank Y. Lozano for a comment on this case.. Thus, we see that (13), regarded as a IIA solution, describes the non-BPS D-instanton. This is in keeping with our later observations about non-BPS D-branes. The interesting point here is that in the absence of the extra parameter $`c_1`$, the same neutral supergravity solution describes both the $`D(1)\overline{D}(1)`$ pair as well as the non-BPS $`D(1)`$ brane. This is presumably a consequence of our earlier observation that there is no tachyon associated with this solution.
### 2.2 Physical Properties
In the physical interpretation of the above three-parameter solution (6),(7) was not presented. We will see that it corresponds to brane-antibrane systems along with condensates.
In a brane-antibrane system, there are two obvious physical parameters $`N`$ and $`\overline{N}`$ which are the numbers of branes and antibranes respectively. In the above supergravity solution too, there are two obvious physical parameters: the RR charge $`Q`$ and the ADM mass $`M_{ADM}`$, which clearly depend on $`N`$ and $`\overline{N}`$. We will discuss in Section 3 the brane interpretation of the third parameter. Before that, however, it will be useful to discuss $`Q`$ and $`M_{ADM}`$ in greater detail.
For convenience, we consider wrapping the spatial world-volume directions on a torus $`T^p`$ of volume $`V_p`$ (this is always possible, since the metric and other fields do not depend on these directions). The RR charge $`Q`$, defined by an appropriate surface integral over the sphere-at-infinity in the transverse directions (see, e.g. ), is given by
$$Q=2\eta N_pr_0^{7p}k\sqrt{c_2^21},$$
(16)
where
$$N_p\frac{(8p)(7p)\omega _{8p}V_p}{128\pi G_N^{10}},$$
(17)
and $`\omega _d=\frac{2\pi ^{(d+1)/2}}{\mathrm{\Gamma }((d+1)/2)}`$ is the volume of the unit sphere $`S^d`$. We have normalized the charge $`Q`$ such that the BPS relation becomes $`M_{BPS}=Q`$.
The ADM mass $`M`$ is defined, in terms of the Einstein-frame metric, by <sup>5</sup><sup>5</sup>5This definition differs from the one presented, e.g. in , by an overall factor.
$$g_{00}=1+\frac{16\pi G_N^{10p}M}{(8p)\omega _{8p}r^{7p}}+\mathrm{higher}\mathrm{order}\mathrm{terms}$$
(18)
where $`G_N^{10p}=G_N^{10}/V_p`$.
This gives us
$$M=N_pr_0^{7p}\left[\frac{3p}{2}c_1+2c_2k\right].$$
(19)
Since the solution is generically non-BPS, $`M`$ is different from $`M_{BPS}Q`$. The mass difference is given by
$$\mathrm{\Delta }MMM_{BPS}=N_pr_0^{7p}\left[\frac{3p}{2}c_1+2k(c_2\sqrt{c_2^21})\right].$$
(20)
In order to have a better understanding of the space of solutions represented by (6),(7), we now consider some special limiting cases.
The BPS case $`(\overline{N}=0)`$
Since the BPS D$`p`$-brane clearly respects the symmetry (1), it should be part of our solution space.
We recall that the D$`p`$-brane solution is given by
$`ds^2`$ $`=`$ $`f_p^{\frac{p7}{8}}dx_\mu dx^\mu +f_p^{\frac{p+1}{8}}(dr^2+r^2d\mathrm{\Omega }_{8p}^2),`$
$`e^\varphi `$ $`=`$ $`f_p^{\frac{3p}{4}},`$
$`C_{01\mathrm{}p}^{(p+1)}`$ $`=`$ $`\eta {\displaystyle \frac{1}{2}}(f_p^11),`$
$`f_p`$ $`=`$ $`1+{\displaystyle \frac{\mu _0}{r^{7p}}},`$ (21)
with ADM-mass $`M_{Dp}`$ and charge $`Q`$ given by
$$M_{Dp}=Q=\mu _0N_p,$$
(22)
This solution indeed exists in a “scaled neighbourhood” of the point $`(r_0,c_1,c_2)=(0,c_m,\mathrm{})`$, defined by
$`r_0^{7p}`$ $`=`$ $`ϵ^{\frac{1}{2}}\overline{r}_0^{7p},`$
$`c_1`$ $`=`$ $`c_mϵ{\displaystyle \frac{8\overline{k}^2}{(p+1)(7p)c_m}},`$
$`c_2`$ $`=`$ $`{\displaystyle \frac{\overline{c}_2}{ϵ}},`$ (23)
where $`c_m=(\frac{32(8p)}{(p+1)(7p)^2})^{1/2}`$ denotes the point where $`k=0`$. The second condition is better stated as
$$k=ϵ^{\frac{1}{2}}\overline{k}.$$
(24)
The scaling is defined by the limit $`ϵ0`$ such that $`\overline{r}_0,\overline{c}_2`$ and $`\overline{k}`$ are fixed.
It is easy to check that the solution (6) reduces to (21) with
$$\mu _0=2c_2kr_0^{7p}=2\overline{c}_2\overline{k}\overline{r}_0^{7p}.$$
(25)
It is useful to consider the three-parameter space of solutions as parameterised by $`M,Q,c_1`$. Figure 1 depicts the $`M,c_1`$ plane for a given fixed $`Q`$. The BPS solution corresponds to the scaled neighbourhood represented by the shaded circle. Other parts of the figure will be explained later.
The D$`p`$-$`\overline{\mathrm{D}p}`$System ($`N=\overline{N}`$)
In this case the RR charge $`Q(N\overline{N})`$ must vanish. According to (16) this corresponds to the subspace
$$|c_2|=1.$$
(26)
We represent this subspace in Fig 2.
Now (26) implies $`c_2=\pm 1`$. As remarked in Section 3 below, the physically relevant choice for $`p>3`$ is $`c_2=1`$, while for $`p<3`$ it is $`c_2=1`$ (for $`p=3`$ the two choices are physically equivalent). To simplify the discussion we will present the formulae in the rest of this section for $`p>3`$; it is straightforward to write down the formulae in the other cases.
The solution now reads
$`e^{2A}`$ $`=`$ $`\left({\displaystyle \frac{f_{}}{f_+}}\right)^\alpha ,`$
$`e^{2B}`$ $`=`$ $`f_{}^\beta _{}f_+^{\beta _+},`$
$`e^\varphi `$ $`=`$ $`(f_{}/f_+)^\gamma ,`$
$`e^\mathrm{\Lambda }`$ $`=`$ $`0,`$ (27)
where
$`\alpha `$ $`=`$ $`(7p)\left({\displaystyle \frac{(3p)c_1+4k}{32}}\right),`$
$`\beta _\pm `$ $`=`$ $`{\displaystyle \frac{2}{7p}}\left({\displaystyle \frac{(p+1)((p3)c_14k}{32}}\right),`$
$`\gamma `$ $`=`$ $`{\displaystyle \frac{1}{16}}\left((7p)(p+1)c_14(3p)k\right).`$ (28)
These represent the most general 2-parameter ($`r_0,c_1`$) solution of Type II supergravity with no gauge field and SO(p,1) $`\times `$ SO(9-p) symmetry.
Consider for instance the case $`p=6`$. The solution reads
$$\begin{array}{c}e^{2A}=\left(\frac{1r_0/r}{1+r_0/r}\right)^{(4k3c_1)/32},\hfill \\ e^{2B}=\left(1r_0/r\right)^{2+7(3c_14k)/32}\left(1+r_0/r\right)^{27(3c_14k)/32},\hfill \\ e^\varphi =\left(\frac{1r_0/r}{1+r_0/r}\right)^{(7c_1+12k)/16}\hfill \end{array}$$
(29)
where $`k=\sqrt{47c_1^2/16}`$.
The Einstein metric has a curvature singularity at $`r=r_0`$. The scalar curvature in (29), e.g., goes as
$$\frac{1}{(rr_0)^{2+\beta _{}}}.$$
(30)
The physical regime is $`rr_0`$. In the case of a single D$`p`$-brane the curvature singularity is resolved by the appropriate inclusion of the brane degrees of freedom. We will discuss this issue in our case later on.
For the specific value
$$c_1=0,$$
(31)
we get
$$\begin{array}{c}e^{2A}=\left(\frac{1\frac{r_0}{r}}{1+\frac{r_0}{r}}\right)^{1/4},\hfill \\ e^{2B}=(1\frac{r_0}{r})^{1/4}(1+\frac{r_0}{r})^{15/4},\hfill \\ e^\varphi =\left(\frac{1\frac{r_0}{r}}{1+\frac{r_0}{r}}\right)^{3/2},\hfill \end{array}$$
(32)
which is the coincident D6-$`\overline{\mathrm{D6}}`$ solution in isotropic coordinates. In Fig 2, this corresponds to the point $`(M,c_1)=(M_0,0)`$.
The above observation implies that for $`c_10`$ we get a generalisation of the coincident D6-$`\overline{\mathrm{D6}}`$ solution. We will argue in the next section that the parameter $`c_1`$ is related to the “vev” <sup>6</sup><sup>6</sup>6We actually consider generically off-shell values of the tachyon. The issue of why we may have supergravity solutions corresponding to an off-shell tachyon is discussed in Sec 3.1. of (the zero momentum mode of the) the open string tachyon arising from open strings stretched between the D6 and $`\overline{\mathrm{D6}}`$ (and more generally between D$`p`$and $`\overline{\mathrm{D}p}`$) branes. The Sen solution corresponds to the particular case where the tachyon vev is zero.
Other cases of $`\mathrm{\Delta }M=0`$
Clearly, from (20) we can have
$$M=Q$$
(33)
if we have
$$(3p)/2c_1+2k(c_2\sqrt{c_2^21})=0.$$
(34)
This solution (represented by $`c_1=c_e`$ in Figs 1,2) is nonsupersymmetric. Indeed, there is a range of the parameters (see Figs 1,2) in which
$$M<Q,$$
(35)
These solutions cannot correspond to physical states of string theory (for $`Q=0`$, these correspond to negative ADM mass).
This implies that we expect additional contribution to the ADM mass formula, coming perhaps from a better understanding of the curvature singularity at $`r=r_0`$. In the case of BPS D-branes or the fundamental string the ADM mass formula as found by the asymptotic behaviour of the Einstein metric does represent the energy-momentum of the source sitting at the curvature singularity. The reason our case is different may have to do with the fact that we have a naked singularity at $`r=r_0`$; a computation of the Euclidean action similar to that in indeed shows that the action receives contribution not only from $`r=\mathrm{}`$, but also from $`r=r_0`$.
## 3 Tachyon Condensation
In the following, we will interpret the 3-parameter family of supergravity solutions as a bound state of $`N`$ D$`p`$-branes coincident with $`\overline{N}`$ $`\overline{\mathrm{D}p}`$-branes, together with a “vev” $`v`$ of the tachyon condensate. The three parameters $`r_0,c_1,c_2`$ will be argued to correspond to various combinations of the three parameters $`N,\overline{N},v`$.
### 3.1 $`\mathrm{T}`$ in Supergravity
A system of $`N`$ D$`p`$-branes on top of $`\overline{N}`$ $`\overline{\mathrm{D}p}`$-branes has a tachyon arising from the open string stretched between the D$`p`$-branes and the $`\overline{\mathrm{D}p}`$-branes. The tachyon $`T`$ transforms in the $`(N,\overline{N})`$ (and $`T^{}`$ in $`(\overline{N},N)`$) representation of the $`U(N)\times U(\overline{N})`$ gauge group. Consider first the case $`N=\overline{N}`$ (the neutral case). The cases that are studied most are $`N=\overline{N}=1`$. In this case the tachyon is a complex field ($`T,T^{}`$) that transforms in the $`(1,1)(1,1)`$ representation of the $`U(1)\times U(1)`$ gauge group of the world-volume theory. The brane system is unstable due to the tachyon. The tachyon has a potential $`V(T)`$ which is a function of $`|T|^2`$. The D$`p`$-$`\overline{\mathrm{D}p}`$-branes configuration is expected to decay into the closed string (Type II) vacuum. Such a decay into the vacuum is conjectured to happen through the process of tachyon condensation in which the zero-momentum mode of the tachyon gets a specific vev. In particular, it is conjectured that at the minimum of the tachyon potential, denoted by $`|T|=T_0`$, the total energy of the system actually vanishes:
$$E=V(T_0)+2M_{Dp}=0,$$
(36)
where $`M_{Dp}`$ is the mass of a D$`p`$-brane. Equation (36) has been established numerically to a very high accuracy via open string field theory . When $`N>1`$ it was argued in that at the minimum of the potential all the eigenvalues of $`T_0`$ are equal. In the following we will denote $`\frac{1}{N}`$Tr$`(TT^{})`$ by $`|T|^2`$.
Let us ask ourselves how the above phenomenon appears from the viewpoint of closed string theory. We concentrate on the neutral case first ($`Q=0`$) and on the charged case later. There are two ways of looking at the problem:
(a) Real-time: The physical decay process in terms of the brane (open string) variables in which the tachyon rolls down to its minimum is time-dependent. The supergravity background of such a time-dependent brane configuration is naively expected to be time-dependent<sup>7</sup><sup>7</sup>7We remark, though, that the exterior geometry of a pulsating spherically symmetric star is given by the static Schwarzschild solution, thanks to Birkhoff’s theorem. It is not inconceivable, therefore, to have a time-dependent brane configuration with a static supergravity background for $`r>r_0`$. In such a case the time-dependence could presumably be discerned at the level of higher mass modes of the closed string (see for a similar analysis where the supergravity background of a BPS state does not see the “polarisation” of the state, although the higher closed string modes see it.).
(b) Path-in-configuration-space: One can alternatively view the decay as a one-parameter path in the open string configuration space, which for our purposes here is the space of values of $`|T|`$. Except at the two extremities of the path ($`|T|=0,T_0`$), the other values of $`|T|`$ are not at an extremum of $`V(T)`$ and is therefore off-shell. Let us ask how such a path would appear in the closed string description. Let us imagine doing an experiment in which gravitons and other massless closed string probes are scattered off the brane-antibrane system for various values of $`|T|`$ as $`|T|`$ is varied from $`0`$ to $`T_0`$. We will assume here that such an experiment makes sense with off-shell values of the tachyon <sup>8</sup><sup>8</sup>8Coupling on-shell bulk degrees of freedom to off-shell brane degrees of freedom is also familiar from AdS/CFT.. In principle one can imagine coupling closed string degrees of freedom to the off-shell tachyon through, e.g., the modified DBI action appropriate to brane-antibrane systems. The supergravity solution away from the brane will have the same symmetry as the brane-antibrane system, namely (1). However, the metric and other fields must reflect the extra parameter $`|T|`$. We will try to argue that the one-parameter deformation represented by $`c_1`$ in our solution corresponds to this $`|T|`$.
We begin by asking whether we see in the supergravity description an analogue of the tachyon potential. The obvious supergravity counterpart of the total energy $`E`$ (Eqn. (36)) of the brane-antibrane system is the ADM mass (19). For the suggested identification to be correct we should have
$$M=V(T)+2NM_{Dp}^{(1)}.$$
(37)
where by $`M^{(1)}`$ we mean the ADM mass for a single D$`p`$ brane. The supergravity solution in question here is the 2-parameter family (27) of solutions parameterised by $`(r_0,c_1)`$. Since the left hand side of (37) is the ADM mass (19), viz.
$$M=N_pr_0^{7p}\left[\frac{3p}{2}c_1+\left(\frac{2(8p)}{7p}\frac{(p+1)(7p)}{16}c_1^2\right)^{\frac{1}{2}}\right],$$
(38)
let us ask whether the the qualitative behaviour of $`M`$ as a function of $`c_1`$ in (38) agrees with the right hand side of (37) for some appropriate identification between $`c_1`$ and $`T`$.
Comment on branches: As explained in Appendix A, the dependence of the ADM mass on $`c_1`$ depends on the specific branch of the solution. In the following we will find that it is for the branch $`I_{++}`$ for $`p>3`$ (and $`I_{}`$ for $`p<3`$)<sup>9</sup><sup>9</sup>9For $`p=3`$ and $`Q=0`$ $`I_{++}`$ and $`I_{}`$ are physically indistinguishable. which lends to a tachyon interpretation. Later on we will briefly comment on the possible interpretation of the other branches.
Once we choose the appropriate branch of the supergravity solution, the qualitative behaviour of $`M`$ as a function of $`c_1`$ (at a fixed $`r_0`$) is given by Fig 3.
Consider first the case $`p=6`$. When $`c_1=0`$ we have the coincident D6-D$`\overline{6}`$ solution . The ADM mass (38) for $`p=6,c_1=0`$ is $`M=4N_pr_0`$. We will argue in Sec. 3.2 that this mass coincides with
$$M=2NM_{Dp}^{(1)}.$$
(39)
This implies that $`V(T)=0`$ at $`c_1=0`$; since the tachyon potential vanishes only at $`T=0`$ , we conclude that
$$T=0\mathrm{at}c_1=0.$$
(40)
As we will see, the last equation is valid for all $`p`$. This will imply that the subspace of our three-parameter solution defined by $`c_1=0`$ represents D$`p`$-$`\overline{\mathrm{D}p}`$branes with zero value of the tachyon $`|T|`$, that is, brane-antibrane configurations which sit at the maximum of the tachyon potential.
Let us now consider small deformations away from $`c_1=0`$. Since $`V(T)`$ is known to be a function only of $`|T|^2`$, we expect the ADM mass, and hence $`c_1`$, to be a function of $`|T|^2`$ too. For small deformations, we can write
$$c_1=a|T|^2+b|T|^4+\mathrm{}$$
(41)
Clearly $`a>0`$. It is easy to see that the behaviour of the ADM mass $`M`$ as a function of $`|T|`$ (Fig. 3(b)) qualitatively matches the behaviour of $`V(|T|)`$ near $`T=0`$.
Tachyon condensation
In Fig 3(b) we have not plotted the ADM mass in the whole range of $`|T|`$ because (41) is valid only near $`T=0`$. The question then is whether our solution can describe the full double-well potential $`V(T)`$. In other words, can we describe the process of tachyon condensation all the way to the vacuum?
In Fig 2, vacuum is represented by any point in the line $`M=0`$. Any path connecting the point $`(M_0,c_1=0)`$ to this line (e.g. path I or path II) therefore in principle represents a family of supergravity solutions corresponding to a flow of $`|T|`$ from $`|T|=0`$ to $`|T|=T_0`$.
To know what the actual path is, we need to have a more precise knowledge of mapping (more detailed than (41)) between the open string variables $`(N,|T|)`$ to the supergravity variables $`(r_0,c_1)`$. Assuming that such maps exist and are smooth and invertible, the generic form will be
$$\begin{array}{c}r_0=\stackrel{~}{f}_1(N,|T|^2),c_1=\stackrel{~}{f}_2(N,|T|^2)\hfill \\ N=\stackrel{~}{g}_1(r_0,c_1),|T|=\stackrel{~}{g}_2(r_0,c_1)\hfill \end{array}$$
(42)
These can alternatively be stated as a map $`(N,|T|)(M,c_1)`$:
$$\begin{array}{c}M=f_1(N,|T|^2),c_1=f_2(N,|T|^2)\hfill \\ N=g_1(M,c_1),|T|=g_2(M,c_1)\hfill \end{array}$$
(43)
Of course (42),(43) should be consistent with (41) near $`T=0`$ (we need to consider the coefficients $`a,b,\mathrm{}`$ to be functions of $`r_0`$ or $`N`$).
The path I in Fig 2 corresponds, in terms of (42), to $`r_0=\stackrel{~}{f}_1(N)`$ and $`c_1=\stackrel{~}{f}_2(|T|^2)`$. This path corresponds to the plot Fig 3(a) of $`M`$ as a function of $`c_1`$ at fixed $`r_0`$. It has the unphysical feature that it does not stop at $`M=0`$ and goes down to the domain of $`M<0`$.
Path II in Fig 2 requires the functions $`\stackrel{~}{f}_{1,2}`$ (or the functions $`f_{1,2}`$) to be necessarily a function of two variables. In other words, the flow of $`|T|`$ from $`0`$ to $`T_0`$ should mean here that both $`r_0`$ and $`c_1`$ should change appropriately to take the solution to the point $`(M,c_1)=(0,c_m)`$. The nice feature of this path is that it automatically ends at the flat space solution, since $`c_1`$ cannot go beyond $`c_m`$ (actually there is another branch of solution (Branch II, appendix A) for $`c_1>c_m`$, but it can be shown that the ADM mass increases for $`c_1>c_m`$).
In the absence of a decoupling limit (as we will discuss in Section 3.3) it may not be possible to determine the exact functions mentioned in (42) or (43) and therefore to know any more about the nature of $`V(T)`$ than what we have already presented here. In any case, if an analysis of brane degrees of freedom is expected to remove the $`M<0`$ region, presumably the formulae for the mass will change.
In summary, we see that a path exists (path II in Fig 2) in our space of solutions which describes the flow of $`|T|`$ from $`0`$ to $`T_0`$ and the behaviour of the ADM mass $`M`$ along this path matches the qualitative features of $`V(T)`$.
The other branches
In the above we have discussed only the branch $`I_{++}`$ (see Appendix A for notation) for $`p3`$ and $`I_{}`$ for $`p<3`$. It is easy to see that the behaviour of the branches $`I_+,I_+`$ are outright unphysical. This leaves $`I_{}`$ for $`p3`$ and $`I_{++}`$ for $`p<3`$. In this branch (except for $`p=3`$) for small deformations of $`c_1`$ away from zero, $`M`$ initially rises beyond the combined rest mass of the brane-antibrane system and then falls again. This seems puzzling since equation (37) does not allow such an increase in the energy of the system. We should recall however that when the vev of the tachyon field is zero the world-volume gauge group is not broken. That means that we are allowed to have other condensates such as a gluon condensate. This can increase the energy of the system. An estimate of such an increase can be obtained from the modified DBI action
$$S=T_pd^{p+1}\sigma e^\varphi V(T)\sqrt{\mathrm{det}\left[G_{ij}+2\pi \alpha ^{}\left(F_{ij}+_iT_jT\right)\right]}.$$
(44)
The interpretation of the $`c_1`$ deformation (for $`p3`$) in these branches could therefore be in terms of a gluon condensate. However,it remains a mystery in that case why (a) there is no such phenomenon for $`p=3`$ (since the branches $`I_{++}`$ and $`I_{}`$ appear to be identical), and (b) why the ADM mass starts to decrease after a while.
Non-BPS D-branes
Since we are only discussing the tachyon condensate in terms of a real quantity $`|T|`$ we are left with the possibility that our supergravity solution may represent a real tachyon as well. Recall that a real tachyon characterizes the non-BPS Dp branes, i.e. $`p`$ odd for IIA and $`p`$ even for IIB, which are obtained from the D$`p`$-$`\overline{\mathrm{D}p}`$-brane system by a $`(1)^{F_L}`$ projection. So the natural question arises: which brane system does the supergravity solution describe. It is plausible that in the neutral case the solution describes both. In both cases the background has no RR charge, and one expects the full $`SO(p,1)\times SO(9p)`$ symmetry. The solution (27) is the most general one that satisfies these conditions. The question is whether the ADM mass of a non-BPS brane (with or without tachyon) occurs in these solutions. We recall that the tension of non-BPS Dp branes (for $`N=1`$) is related to the tension of the D$`p`$-$`\overline{\mathrm{D}p}`$-brane system by $`M_{\mathrm{non}\text{-}\mathrm{BPS}}=\frac{1}{\sqrt{2}}M_{Dp\text{-}\overline{Dp}}`$, reflecting a bound system. For $`N>1`$ too, the tension of the non-BPS Dp brane system $`M_{\mathrm{non}\text{-}\mathrm{BPS}}^{(N)}`$ should be less than that of the combined rest mass $`2NM_{Dp}^{(1)}`$ of the brane-antibrane system. Since the values of ADM mass discussed in the context of (37) range all the way from $`2NM_{Dp}^{(1)}`$ to $`0`$, we see that in a suitable range of parameters the solution (27) does have ADM masses that can be fitted to $`M=M_{\mathrm{non}\text{-}\mathrm{BPS}}^{(N)}+\stackrel{~}{V}(T)`$ where $`\stackrel{~}{V}(T)`$ is the potential for the real tachyon in this case.
The charged case: $`Q0`$
In this case we expect the relation
$$M=(N+\overline{N})M_{Dp}^{(1)}+V(T).$$
(45)
where $`M_{Dp}^{(1)}`$ denotes the ADM mass for a single D$`p`$ brane. The analysis of the binding energy in the next section once again suggests that $`c_1=0`$ corresponds to the point where the tachyon potential vanishes, which we expect to be for vanishing tachyon field. The discussion of tachyon condensation is similar to the neutral case. Again path II in Fig 1 is more physical than path I because the former ends at the BPS point and does not go to the region $`M<Q`$. The qualitative behaviour of $`M`$ along this path again matches the qualitative features of a tachyon potential which has a local maximum at $`|T|=0`$ and a minimum at $`|T|=T_0`$ where we denote $`\frac{1}{N}`$Tr$`(TT^{})`$ by $`|T|^2`$ (we assume that all the eigenvalues of $`TT^{}`$ are the same, namely $`T_0^2`$, at the minimum). We expect that at the minimum $`V(T)=(|N\overline{N}|(N+\overline{N}))M_{Dp}^{(1)}`$.
### 3.2 Dp-brane probes and Binding energy
In the last section we mentioned that $`V(T)=0`$ corresponds to $`c_1=0`$. We derive this in the present section.
We will consider the general 3-parameter solution parametrized by $`(r_0,c_1,c_2)`$. Let us define the binding energy of the D$`p`$-$`\overline{\mathrm{D}p}`$-branes solution to be
$$E_B=(N+\overline{N})M_{Dp}^{(1)}M,$$
(46)
where $`M`$ is given by (19) and $`M_{Dp}^{(1)}`$ represents the rest mass of a single D$`p`$-brane (or $`\overline{\mathrm{D}p}`$-brane), given by (22) with the scale parameter $`\mu _0=\mu _0^{(1)}`$, which depends on $`g_{str}`$ and $`p`$, the dimensionality of the brane.
In view of the equation (37),
$$E_B=V(T).$$
(47)
A straightforward comparison between $`(N+\overline{N})M_{Dp}`$ and $`M`$ of (19) is hampered by the fact that we do not know a priori the relation between the two parameters $`r_0`$ and $`\mu _0`$ that characterise the respective solutions (6) and (21). We will find this relation by the following strategy.
We consider the static force between a D$`p`$-$`\overline{\mathrm{D}p}`$-branes system and a D$`p`$-brane probe (respectively a $`\overline{Dp}`$-brane probe) at a distance $`r`$. This can be computed in two ways:
(a) From supergravity:
$$S_{\mathrm{probe}}=\frac{1}{g_sl_s^{p+1}}d^{p+1}\sigma (e^\varphi \sqrt{\widehat{G}}\pm C_{p+1})$$
(48)
where $`G_{MN}=e^{\varphi /2}g_{MN}`$ represents the string frame metric corresponding to the solution (6) and $`\widehat{G}`$ is its pull-back to the world-volume. For a D$`p`$(resp. $`\overline{\mathrm{D}p}`$) probe, we use the upper (resp. lower) sign.
Subtracting the flat space DBI part, and keeping only the leading term in the $`1/r`$ expansion we get
$$S_{\mathrm{probe}}=2k\frac{V_p}{g_sl_s^{p+1}}(\frac{r_0}{r})^{7p}(c_2\sqrt{c_2^21}).$$
(49)
(b) By a string theory computation:
$$DpD\overline{p}|\mathrm{exp}(\beta H)|Dp$$
(50)
where the states are regarded as boundary states constructed out of closed-string oscillators. (We consider here the case of the D$`p`$-probe first.) At weak coupling and for $`\mathrm{T}=0`$, the boundary state on the left is given by
$$DpD\overline{p}|=Dp|D\overline{p}|.$$
(51)
We will assume that (51) can be used for computation of the leading term in the $`1/r`$ expansion for large distances $`r`$, when $`\mathrm{T}=0`$ (see for earlier work on connection between boundary states and classical solutions). Since the static force between two Dp-branes vanishes, the computation (b) then reduces, at $`\mathrm{T}=0`$, to
$$D\overline{p}|\mathrm{exp}(\beta H)|Dp.$$
(52)
This latter can be computed at large distances from supergravity, by the DBI action of a D$`p`$-brane probe in the background of a $`\overline{\mathrm{D}p}`$\- brane:
$$S_{\mathrm{probe}}^{}\frac{1}{g_sl_s^{p+1}}d^{p+1}\sigma [e^\varphi \sqrt{\widehat{G}}+C^{(p+1)}],$$
(53)
where the metric, dilaton and the RR potential are now obtained from (21), with $`\mu _0=\overline{N}\mu _0^{(1)}`$. We get, again after subtracting the flat space DBI part, and keeping only the leading term in the $`1/r`$ expansion,
$$S_{\mathrm{probe}}^{}=\frac{V_p}{g_sl_s^{p+1}}\left(2\frac{\overline{N}\mu _0^{(1)}}{r^{7p}}\right).$$
(54)
This result holds for the D$`p`$-probe. For the $`\overline{\mathrm{D}p}`$-probe we need to replace $`\overline{N}N`$ in the above expression.
Matching (49) and (54) leads to
$`N\mu _0^{(1)}`$ $`=`$ $`kr_0^{7p}(c_2+\sqrt{c_2^21}),`$
$`\overline{N}\mu _0^{(1)}`$ $`=`$ $`kr_0^{7p}(c_2\sqrt{c_2^21}).`$
from this we deduce that
$$Q=N_p\mu _0^{(1)}(N\overline{N})$$
(56)
and
$$\mu _0^{(1)}(N+\overline{N})=2kr_0^{7p}c_2.$$
(57)
Using (19), (22) and (LABEL:mu-vs-r0) we can find the zero of the binding energy (46) of the D$`p`$-$`\overline{\mathrm{D}p}`$bound state. We get
$$E_B=0=N_pr_0^{7p}\left[\frac{3p}{2}c_1\right].$$
(58)
Clearly $`E_B`$ vanishes at $`c_1=0`$<sup>10</sup><sup>10</sup>10 The case $`p=3`$ is subtle and we extrapolated the result to this value of $`p`$ from the other values. An alternative way would presumably be to use some other probe.. In view of the identification (47), this implies that
$$c_1=0V(T)=0,$$
(59)
as promised in the last section. Note that
(a) If we put $`c_1=0`$ in (19) we indeed get $`M=M_{Dp}+M_{\overline{Dp}}`$, consistent with the vanishing of the binding energy,
(b) Equations (56) and (57) give us essentially $`N\overline{N}`$ and $`N+\overline{N}`$ in terms of the supergravity parameters in the subspace $`c_1=0`$,
(c) The expression for the total mass (57) matches exactly with the BPS mass (25) (recall that at the BPS point $`\overline{N}=0`$.
### 3.3 Open-closed String Duality
In the spirit of the AdS/CFT correspondence (for a review see ), it is natural to ask whether we can apply a decoupling limit of the brane modes from the bulk modes to the supergravity description of the D$`p`$-$`\overline{\mathrm{D}p}`$-branes system. Typically for Dp-branes this is a low energy limit with the resulting background being the near-horizon metric. In the present case, the closest analogue of the near horizon metric is some suitably scaled neighbourhood of $`r=r_0`$. However, it is easy to see that for the neutral solution (27) there is no such region which by itself is a solution of the supergravity field equations. Also, we cannot find an appropriate rescaling that keeps a metric finite in $`l_s`$ units as $`l_s0`$. This means that the interactions between the open and closed strings remain relevant.
Another manifestation of this issue is the form of the potential $`V(r)`$ for a graviton scattered on the D$`p`$-$`\overline{\mathrm{D}p}`$-branes. The potential is depicted in figure 4. Near $`r=r_0`$ it goes like $`\frac{1}{(rr_0)^2}`$ while at infinity it approaches $`\omega ^2`$ where $`\omega `$ is the frequency of the scattered graviton. The potential poses no barrier for the gravitons sent from infinity to reach the $`r=r_0`$ and their absorption cross section does not vanish <sup>11</sup><sup>11</sup>11For a similar but detailed analysis see .. The absence of a decoupling of the closed strings from the open strings prevents us from making a precise correspondence between the field theory on the D$`p`$-$`\overline{\mathrm{D}p}`$-branes world-volume and the supergravity background. This suggests that there is also a limitation on the quantitative understanding of the tachyon condensation process by using only the open string description.
The singularity of the supergravity solution at $`r=r_0`$ is time-like. Having such a singularity of the classical geometry which we can reach at a finite proper time, there is the natural question whether it is resolved quantum mechanically. One criterion is the existence of a self-adjoint Laplacian. This can still be the case even if the metric is geodesically incomplete. The requirement is the existence of a non-normalizable solution of the wave equation. This criterion is satisfied by our geometry. To see that we consider the Laplace equation in the form
$$\frac{^2\varphi }{t^2}=A\varphi .$$
(60)
The equation $`A\varphi =\lambda \varphi `$ takes the form
$$\rho ^\beta _\rho (\rho _\rho \varphi )=\lambda \varphi ,$$
(61)
where $`\beta =2p15+2(7p)(\frac{3p}{8}c_1\frac{k}{2})`$ and $`\rho =rr_0`$. Defining $`z=\sqrt{\lambda }\rho ^{(1\beta )/2}`$ we get
$$\varphi ^{\prime \prime }+\frac{\varphi ^{}}{z}+\frac{4}{(1\beta )^2}\varphi =0.$$
(62)
This has Bessel function solutions behaving like $`1`$ and $`\mathrm{ln}\rho `$. The norm of the latter $`𝑑\rho \rho \rho ^2`$ diverges.
## 4 The Four-Parameter Solution
In this section we will briefly describe the most general $`p`$-brane solution of Type II string theory in which we relax the requirement of Poincare invariance in the $`(p+1)`$ dimensional world-volume. In other words, we ask ourselves about the most general solution which respects the symmetry
$$𝒮^{}=SO(p)\times SO(9p),$$
(63)
Clearly the previous 3-parameter solution already respects this symmetry and hence should be part of this most general family of solutions. The modified ansatz for the Einstein metric is
$$ds^2=\mathrm{exp}(2A(r))(f(r)dt^2+dx_mdx^m)+\mathrm{exp}(2B(r))\left(dr^2+r^2d\mathrm{\Omega }_{8p}^2\right),$$
(64)
where we split the world-volume index $`\mu `$ as $`0,m=1,\mathrm{},p`$. The ansatz for the dilaton and the gauge potential remain the same as in (2).
The equations of motion for this ansatz have been written down in appendix B. Once again the mathematical solution of the differential equations has been worked out in . We write the explicit solution in Appendix B for completeness and discuss here some salient physical features (see Fig 5)
$``$ The general solution has 4 independent parameters $`(r_0,c_1,c_2,c_3)`$. The Poincare-invariant 3-parameter subspace discussed in the previous sections corresponds to $`c_3=0`$. In Fig 5 this is schematically represented by the arm AC of the triangle ABC. $`c_30`$ breaks world-volume Poincare invariance.
$``$ The two-parameter subspace $`(c_1,c_3)=(\frac{3p}{2(7p)},2)`$ corresponds to the black $`p`$-brane solutions of . This has already been identified in . In Fig 5, this is represented by the arm AB of the triangle. Recall that the black $`p`$-branes are parameterised by their charge and mass (equivalently $`r_+,r_{}`$, the outer and inner horizons). Note that the BPS D-brane can be reached as a limit along the arm BA, like it can be reached along CA, although the $`c_3`$ values characterising these two arms are different. It is likely that there are continuous families of solutions between BA and CA (corresponding to different $`c_3`$ values) which can reach the BPS solution under a limiting procedure.
$``$ The three-parameter subspace defined by $`|c_2|=1`$ describes the most general spherically symmetric solution with no gauge fields. This is represented by the arm BC of the triangle. It is well-known that the neutral limit of the black $`p`$-brane (point B) corresponds to the Schwarzschild black hole in $`10p`$ dimensions ($`\times T^p`$, assuming a wrapped $`p`$-brane). On the other hand, as discussed at great length in this paper, the neutral limit of the arm AC corresponds to the coincident brane-antibrane solutions. The arm BC therefore provides interpolating solutions which connect the brane-antibrane solution to the Schwarzschild solution.
It is clear that there is a rather rich phase structure in Fig 5. Parts of this diagram have obvious decoupling limits and dual field theory descriptions. It would be interesting to chart out these parts completely .
Interpolations similar to the arm BC are of paramount importance to the study of the D1/D5 system and the five-dimensional black hole . It has been found that CFT descriptions seem to work in some contexts for non-rotating BTZ black holes which are the analogues of Schwarzschild black holes in AdS<sub>3</sub>. An interpolation of such a solution to a brane-antibrane solution of the D1/D5 system would shed light on both brane-antibrane dynamics and nonsupersymmetric black holes.
It has been pointed out by that the equations of motion of the above system are identical to those of a Toda molecule. It is tempting to construct a “mini-superspace” kind of model for this space based on Toda dynamics.
## 5 Discussion
In this paper we constructed localised supergravity solutions corresponding to bound states of $`N`$ D$`p`$-branes coinciding with $`\overline{N}`$ $`\overline{\mathrm{D}p}`$-branes for $`p=0,1,\mathrm{},6`$ (and non-BPS D-branes of odd (even) dimensions of Type IIA (Type IIB) string theory)<sup>12</sup><sup>12</sup>12The case $`p=1`$ has been mentioned separately in Section 2. We constructed these by looking for the most general solution of Type II A/B supergravity (in the presence of a single RR gauge field) which respect world-volume Poincare invariance and rotational invariance in the transverse directions. Contrary to the naive expectation that the solution should have only two parameters corresponding to the charge and the mass, we found that the most general solution has one extra parameter. We found that in the physically relevant branch there are two special values of the extra parameter at which the ADM mass respectively coincides with (a) the combined rest mass of the branes and antibranes, and (b) the mass of the BPS configuration of $`N\overline{N}`$ branes <sup>13</sup><sup>13</sup>13for $`N>\overline{N}`$; for $`N<\overline{N}`$ these will be $`\overline{N}N`$ antibranes.. In the case $`N=\overline{N}`$ (zero RR charge) the point (b) represents flat space. The case $`N=\overline{N}`$ is extensively studied from the point of view of open strings living on the brane-antibrane system, and we recognised the solutions (a) and (b) as the supergravity background corresponding to the maximum and the minimum of the tachyon potential. This lead us to interpret the extra parameter in our solution as the supergravity manifestation of an expectation value of the tachyon. We matched the qualitative behaviour of the ADM mass as a function of this extra parameter with the behaviour of the tachyon potential $`V(T)`$.
We noticed the absence of a decoupling of the bulk closed strings from the brane-antibrane open strings. This means that the interactions between the open and closed strings remain relevant and suggets that there is also a limitation on the quantitative understanding of the tachyon condensation process by using only the open string description.
We briefly discussed a more general (four-parameter) space of solutions in which we assume only rotational invariance in the spatial directions on the world-volume. This space includes brane-antibrane pairs, BPS D-branes, the black $`p`$-branes of and Schwarzschild black holes. The detailed understanding of this four-parameter space in terms of brane variables is an outstanding problem.
Acknowledgement: We would like to thank A. Kumar and P. Townsend for discussions.
## Appendix
## A Real Sections of the Supergravity Solution
As remarked in the text, the three parameters $`(r_0,c_1,c_2)`$ characterising the supergravity solution (6),(7) appear as integration constants in the solution of differential equations and as such could be complex. However, this would generically make the metric, dilaton and gauge field also complex. We find that there are three distinct 3-dimensional domains of $`(r_0,c_1,c_2)`$, described below as Branches I, II and III, where the supergravity fields remain real.
Branch I:
$$\begin{array}{c}c_1(0,c_m),c_m=\sqrt{\frac{8p}{8(p+1)(7p)}}\hfill \\ c_2(\mathrm{},1)(1,\mathrm{})\hfill \\ \mu r_0^{7p}R\hfill \\ \eta =\pm 1\hfill \end{array}$$
(65)
We will assume in this section that we have already fixed the $`Z_2`$ symmetries (8) of the solution by implementing (9),(10). For Branch I, the remaining choices of signs are best discussed by thinking of four sub-branches, depending on whether the signs of $`(c_2,\mu r_0^{7p})`$ are $`++,+,+`$ and $``$ respectively. We denote these as $`I_{++},I_+,I_+,I_{}`$ respectively (each of these will also contain $`\eta =\pm `$). The formulae for the ADM mass and charge for Branch I is given by (19),(16). Explicitly
$$\begin{array}{c}M=N_pr_0^{7p}\left[\frac{3p}{2}c_1+2c_2\sqrt{\frac{2(8p)}{7p}\frac{(p+1)(7p)}{16}c_1^2}\right]\hfill \\ Q=2\eta N_pr_0^{7p}\sqrt{\frac{2(8p)}{7p}\frac{(p+1)(7p)}{16}c_1^2}\sqrt{c_2^21},\hfill \end{array}$$
(66)
The behaviour of these functions depends on the signs of $`c_2`$ and $`\mu `$. We find that it is the branch $`I_{++}`$ for $`p=3,4,5,6`$ which lends to a tachyon interpretation (Section 3). For $`p=0,1,2,3`$ it is $`I_{}`$.
Branch II
$$\begin{array}{c}c_1(c_m,\mathrm{})k=i\stackrel{~}{k},\stackrel{~}{k}=\sqrt{\frac{2(8p)}{7p}+\frac{(p+1)(7p)}{16}c_1^2}\hfill \\ c_2=i\stackrel{~}{c}_2,\stackrel{~}{c}_2R\hfill \\ \mu r_0^{7p}R\hfill \\ \eta =\pm 1\hfill \end{array}$$
(67)
The mass and charge for this branch read
$$\begin{array}{c}M=N_pr_0^{7p}\left[\frac{3p}{2}c_1+2c_2\sqrt{\frac{2(8p)}{7p}+\frac{(p+1)(7p)}{16}c_1^2}\right]\hfill \\ Q=2\eta N_pr_0^{7p}\sqrt{\frac{2(8p)}{7p}\frac{(p+1)(7p)}{16}c_1^2}\sqrt{(\stackrel{~}{c}_2)^2+1},\hfill \end{array}$$
(68)
Branch III
$$\begin{array}{c}c_1=i\stackrel{~}{c}_1,\stackrel{~}{c}_1R^+\hfill \\ c_2=i\stackrel{~}{c}_2,\stackrel{~}{c}_2R\hfill \\ \mu r_0^{7p}=i\stackrel{~}{\mu },\stackrel{~}{\mu }R\hfill \\ \eta =\pm 1\hfill \end{array}$$
(69)
The mass and charge for this branch read
$$\begin{array}{c}M=N_pr_0^{7p}\left[\frac{3p}{2}c_1+2c_2\sqrt{\frac{2(8p)}{7p}+\frac{(p+1)(7p)}{16}(\stackrel{~}{c}_1)^2}\right]\hfill \\ Q=2\eta N_pr_0^{7p}\sqrt{\frac{2(8p)}{7p}+\frac{(p+1)(7p)}{16}(\stackrel{~}{c}_1)^2}\sqrt{(\stackrel{~}{c}_2)^2+1},\hfill \end{array}$$
(70)
## B Details of the 4-parameter solution
The equations of motion that follow from (3) for the ansatz (64) are
$`A^{\prime \prime }+(p+1)(A^{})^2+(7p)A^{}B^{}+{\displaystyle \frac{8p}{r}}A^{}+{\displaystyle \frac{1}{2}}(\mathrm{ln}f)^{}A^{}={\displaystyle \frac{7p}{16}}S^2,`$
$`A^{\prime \prime }+(p+1)(A^{})^2+(7p)A^{}B^{}+{\displaystyle \frac{8p}{r}}A^{}+{\displaystyle \frac{1}{2}}(\mathrm{ln}f)^{\prime \prime }+`$
$`{\displaystyle \frac{1}{2}}(\mathrm{ln}f)^{}\left((d+1)A^{}+{\displaystyle \frac{1}{2}}(\mathrm{ln}f)^{}+(7p)B^{}+{\displaystyle \frac{8p}{r}}\right)={\displaystyle \frac{7p}{16}}S^2,`$
$`B^{\prime \prime }+(p+1)A^{}B^{}+{\displaystyle \frac{p+1}{r}}A^{}+(7p)(B^{})^2+{\displaystyle \frac{1}{2}}(\mathrm{ln}f)^{}\left(B^{}+{\displaystyle \frac{1}{r}}\right)+`$
$`{\displaystyle \frac{152p}{r}}B^{}={\displaystyle \frac{1}{2}}{\displaystyle \frac{p+1}{8}}S^2,`$
$`dA^{\prime \prime }+(8p)B^{\prime \prime }+(p+1)(A^{})^2+{\displaystyle \frac{8p}{r}}B^{}(p+1)A^{}B^{}+`$
$`{\displaystyle \frac{1}{2}}(\mathrm{ln}f)^{\prime \prime }+{\displaystyle \frac{1}{4}}\left((\mathrm{ln}f)^{}\right)^2+{\displaystyle \frac{1}{2}}(\varphi ^{})^2={\displaystyle \frac{1}{2}}{\displaystyle \frac{7p}{8}}S^2,`$
$`\varphi ^{\prime \prime }+\left((p+1)A^{}+(7p)B^{}+{\displaystyle \frac{8p}{r}}+{\displaystyle \frac{1}{2}}(\mathrm{ln}f)^{}\right)\varphi ^{}={\displaystyle \frac{a}{2}}S^2,`$
$`\left({\displaystyle \frac{\mathrm{\Lambda }^{}}{f^{\frac{1}{2}}}}e^{\mathrm{\Lambda }+a\varphi (p+1)A+(7p)B}r^{8p}\right)^{}=0,`$ (71)
where
$$S=\frac{\mathrm{\Lambda }^{}}{f^{\frac{1}{2}}}e^{\frac{1}{2}a\varphi +\mathrm{\Lambda }dA}.$$
(72)
The solutions depend on four parameters $`r_0,c_1,c_2,c_3`$ (we have interchanged the labels $`c_2,c_3`$ for convenience, compared to ), and are given by
$`f(r)`$ $`=`$ $`e^{c_3h(r)},`$
$`A(r)`$ $`=`$ $`{\displaystyle \frac{(7p)}{32}}\left({\displaystyle \frac{3p}{2}}c_1+(1+{\displaystyle \frac{(3p)^2}{8(7p)}})c_3\right)h(r)`$
$`{\displaystyle \frac{7p}{16}}\mathrm{ln}\left[\mathrm{cosh}(kh(r))c_2\mathrm{sinh}(kh(r))\right],`$
$`B(r)`$ $`=`$ $`{\displaystyle \frac{1}{7p}}\mathrm{ln}\left[f_{}(r)f_+(r)\right]+{\displaystyle \frac{(p3)}{64}}\left((p+1)c_1{\displaystyle \frac{3p}{4}}c_3\right)h(r)`$
$`+{\displaystyle \frac{p+1}{16}}\mathrm{ln}\left[\mathrm{cosh}(kh(r))c_2\mathrm{sinh}(kh(r))\right],`$
$`\varphi (r)`$ $`=`$ $`{\displaystyle \frac{(7p)}{16}}\left((p+1)c_1{\displaystyle \frac{3p}{4}}c_3\right)h(r)`$
$`+{\displaystyle \frac{3p}{4}}\mathrm{ln}\left[\mathrm{cosh}(kh(r))c_2\mathrm{sinh}(kh(r))\right],`$
$`e^{\mathrm{\Lambda }(r)}`$ $`=`$ $`\eta (c_2^21)^{1/2}{\displaystyle \frac{\mathrm{sinh}(kh(r))}{\mathrm{cosh}(kh(r))c_2\mathrm{sinh}(kh(r))}},`$ (73)
where
$`f_\pm (r)`$ $``$ $`1\pm \left({\displaystyle \frac{r_0}{r}}\right)^{7p},`$
$`h(r)`$ $`=`$ $`\mathrm{ln}\left[{\displaystyle \frac{f_{}(r)}{f_+(r)}}\right],`$
$`k^2`$ $`=`$ $`{\displaystyle \frac{2(8p)}{7p}}c_1^2+{\displaystyle \frac{1}{4}}\left({\displaystyle \frac{3p}{2}}c_1+{\displaystyle \frac{7p}{8}}c_3\right)^2{\displaystyle \frac{7}{16}}c_3^2,`$
$`\eta `$ $`=`$ $`\pm 1.`$ (74)
The parameter $`\eta `$ describes whether we are measuring the “brane” charge or the “antibrane” charge of the system.
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# Untitled Document
Witten solution of the Gelfand-Dikii hierarchy
S.M.Natanzon
Moscow State University, e-mail: natanzonmccme.ru
Abstract. We produce formulas, permiting to find the coefficients of Taylor–series expanded of some important solution of the Gelfand-Dikii hierarchy. By the Witten conjecture these coefficients are numbers of intersection of Mumford-Morita-Muller stable cohomological classes of moduli space of $`n`$-spin bundles on Riemann surfaces with punctures.
Mathematics Subject Classifications (1991). 14, 35.
Key words: Gelfand–Dikii hierarchy, $`KP`$–hierarchy, moduly space, Witten conjecture.
1. Introduction
The $`n`$-Gelfand-Dikii hierarchy ($`n`$-K$``$V hierarchy) usually is described in term of formal pseudo-differential operators. Consider functions $`u_j=u_j(x)(j=0,\mathrm{},n2)`$ of infinite number of variable $`x=(x_1,x_2,\mathrm{}.)`$. Let us $`_i=\frac{}{x_i},=_1`$ and $`L=^n+u_{n2}^{n2}+\mathrm{}+u_1+u_0`$. Denote by $`L^{\frac{i}{n}}`$ the pseudo-differential operator such that $`(L^{\frac{i}{n}})^n=L^i`$. Let $`[L^{\frac{i}{n}}]_+`$ be its differential part. Then the $`n`$-Gelfand-Dikii hierarchy \[GD\] is the system of differential equations on $`u_i`$, which follow from infinite system of equations
$$_iL=[[L^{\frac{i}{n}}]_+,L](i=1,2,\mathrm{}).$$
Any solution $`(u_0,\mathrm{},u_{n2})`$ is generated by a function $`v(x_1,x_2,\mathrm{})`$. Among these solutions there exists a remarkable solution $`W`$, which satisfies the string equation, generates a vacuum vector of $`W`$–algebra and has a representation in a form of matrix integral \[AM\]. We call it Witten’s solution because according to the Witten conjecture the function $`F(x_1,x_2,x_3,\mathrm{})=W(x_1,\frac{x_2}{2},\frac{x_3}{3},\mathrm{})`$ is the generating function for numbers of intersection of Mumford-Morita-Muller stable cohomological classes of moduli space of $`n`$-spin bundles on Riemann surfaces with punctures \[W\]. This conjecture was proved by Kontsevich for $`n=2`$ \[Ko\] and by Witten himself for surfaces of genus 0 \[W\].
In this paper we find recurrence relation between coefficients of Taylor series of W. This reduces the Witten conjecture to conjecture that $`n`$-spin Mamford-Morita-Muller numbers satisfy to the same relations. These relations give also an algorithm for calculation of $`n`$-spin Mamford-Morita-Muller numbers in assuming that the Witten conjecture is true. Moreover we prove that $`F(x_1,x_2,\mathrm{},x_{n1},0,0,\mathrm{})=W(x_1,\frac{x_2}{2},\mathrm{},\frac{x_n}{n1},0,0,\mathrm{})`$ and thus the Witten conjecture is true for numbers $`(c_D(\nu ),\overline{M}_{g,s}^{})`$ in notation \[W\].
We prove also that the solution $`W`$ has a representation $`W=\underset{g=0}{\overset{\mathrm{}}{}}W^g`$, where $`W^g`$ are quasihomohenious series of degrees $`(1g)(2+\frac{2}{n})`$. This is some indirect corroboration of the Witten conjecture because according to \[W\] the function $`F`$ has a representation $`F=\underset{g=0}{\overset{\mathrm{}}{}}F^g`$, where $`F^g`$ are quasihomohenious series of degrees $`(1g)(2+\frac{2}{n})`$, corresponding to surfaces of genius $`g`$.
We prove that $`W^0(x_1,\mathrm{},x_{n1},0,\mathrm{})=W(x_1,\mathrm{},x_{n1},0,\mathrm{})`$ and therefore $`F^0(x_1,\mathrm{},x_{n1},0,\mathrm{})=F(x_1,\mathrm{},x_{n1},0,\mathrm{})`$. Last function is polynomial solution of WDVV equations. Some simple formulas for calculation of this solution was found in \[N2\].
Organisation of the paper is as follow. In sect 2-4 we following by \[DN, N1\] represent KP hierarchy as a system of differential equation for $`v=\mathrm{ln}\tau `$. In section 5 we prove that the $`n`$-hierarchy of Gelfand-Dikii is equivalent of a system of differential equations in a form
$$_{i_1}\mathrm{}_{i_k}v=\underset{m=1}{\overset{\mathrm{}}{}}N_{i_1\mathrm{}i_k}^m\left(\begin{array}{ccc}s_1& \mathrm{}& s_m\\ t_1& \mathrm{}& t_m\end{array}\right)_{s_1}_1^{t_1}v\mathrm{}_{s_m}_1^{t_m}v,$$
$`1.1`$
where $`i_j,t_i1,s_i<n`$ and $`N_{i_1\mathrm{}i_k}^m\left(\begin{array}{ccc}s_1& \mathrm{}& s_m\\ t_1& \mathrm{}& t_m\end{array}\right)`$ are rational constants.
In section 6 we investigate the Witten solution W of the system (1.1).
The author thanks B.A.Dubrovin for fruitful discussions.
2. Combinatorial lemma
For natural $`s,i_1,\mathrm{},i_n`$ and integer not negative $`j_1,\mathrm{},j_n`$ defind $`P_s\left(\begin{array}{ccc}i_1& \mathrm{}& i_n\\ j_1& \mathrm{}& j_n\end{array}\right)`$ by recurence formulas:
$$1)P_s\left(\begin{array}{ccc}i_1& \mathrm{}& i_n\\ 0& \mathrm{}& 0\end{array}\right)=0;2)P_s\left(\begin{array}{c}i\\ j\end{array}\right)=C^j_s\text{for}j>0;$$
$$3)P_s\left(\begin{array}{ccc}i_1& \mathrm{}& i_n\\ j_1& \mathrm{}& j_n\end{array}\right)=\frac{1}{n!}C_s^{j_1+\mathrm{}+j_n}\frac{(j_1+\mathrm{}+j_n)!}{j_1!\mathrm{}j_n!}$$
$$\underset{q=1}{\overset{n1}{}}P_s\left(\begin{array}{ccc}i_1& \mathrm{}& i_q\\ j_1& \mathrm{}& j_q\end{array}\right)\frac{1}{(nq)!}C_{s(i_1+\mathrm{}+i_q+j_1+\mathrm{}+j_q)}^{j_{q+1}+\mathrm{}+j_n}\frac{(j_{q+1}+\mathrm{}+j_n)!}{j_{q+1}!\mathrm{}j_n!}$$
for $`(j_1,\mathrm{},j_n)(0,\mathrm{},0)`$.
Let $`\left[\begin{array}{ccc}i_1& \mathrm{}& i_n\\ j_1& \mathrm{}& j_n\end{array}\right]`$ be the set of all matrices, which appear from $`\left(\begin{array}{ccc}i_1& \mathrm{}& i_n\\ j_1& \mathrm{}& j_n\end{array}\right)`$ by permutation of columns. Let $`\begin{array}{ccc}i_1& \mathrm{}& i_n\\ j_1& \mathrm{}& j_n\end{array}`$ be the number of such matrices. Put us
$$P_s\left[\begin{array}{ccc}i_1& \mathrm{}& i_n\\ j_1& \mathrm{}& j_n\end{array}\right]=P_s\left(\begin{array}{ccc}a_1& \mathrm{}& a_n\\ b_1& \mathrm{}& b_n\end{array}\right),$$
where the sum is taken by all
$$\left(\begin{array}{ccc}a_1& \mathrm{}& a_n\\ b_1& \mathrm{}& b_n\end{array}\right)\left[\begin{array}{ccc}i_1& \mathrm{}& i_n\\ j_1& \mathrm{}& j_n\end{array}\right]$$
Lemma 2.1. Let $`m>0,k>0`$ and $`j_n1`$ for $`nm`$. Then
$$P_s\left[\begin{array}{cccccc}i_1& \mathrm{}& i_m& i_{m+1}& \mathrm{}& i_{m+k}\\ j_1& \mathrm{}& j_m& 0& \mathrm{}& 0\end{array}\right]=$$
$$=\{\begin{array}{cc}0,\text{if }si_1+\mathrm{}+i_m+j_1+\mathrm{}+j_m\text{,}\hfill & \\ \frac{1}{k!}\begin{array}{ccc}i_{m+1}& \mathrm{}& i_{m+k}\\ 0& \mathrm{}& 0\end{array}P_s\left[\begin{array}{ccc}i_1& \mathrm{}& i_m\\ j_1& \mathrm{}& j_m\end{array}\right],\text{if }s<i_1+\mathrm{}+i_m+j_1+\mathrm{}+j_m.\hfill & \end{array}$$
Proof: Prove at first the lemma for $`m=1`$ using induction by k. For $`m=k=1`$
$$P_s\left[\begin{array}{c}i_1i_2\\ j_10\end{array}\right]=P_s\left(\begin{array}{c}i_1i_2\\ j_10\end{array}\right)+P_s\left(\begin{array}{c}i_2i_1\\ 0j_1\end{array}\right)=$$
$$=\frac{1}{2}C_s^{j_1}P_s\left(\begin{array}{c}i_1\\ j_1\end{array}\right)C_{s(i_1+j_1)}^0+\frac{1}{2}C_s^{j_1}P_s\left(\begin{array}{c}i_2\\ 0\end{array}\right)C_{si_2}^0=C_s^{j_1}C_s^{j_1}C_{s(i_1+j_1)}^0=$$
$$=\{\begin{array}{cc}0,\hfill & \text{if }si_1+j_1,\hfill \\ C_s^{j_1}=P_s\left[\begin{array}{c}i_1\\ j_1\end{array}\right],\hfill & \text{if }s<i_1+j_1.\hfill \end{array}$$
Prove the lemma for $`m=1,k=N`$, considering that it is proved for $`m=1,k<N`$. If $`si_1+j_1`$, then
$$P_1\left[\begin{array}{cccc}i_1& i_2& \mathrm{}& i_{k+1}\\ j_1& 0& \mathrm{}& 0\end{array}\right]=$$
$$=\frac{1}{k!}C_s^{j_1}\begin{array}{ccc}i_2& \mathrm{}& i_{k+1}\\ 0& \mathrm{}& 0\end{array}P_s\left(\begin{array}{c}i_1\\ j_1\end{array}\right)C_{s(i_1+j_1)}^0\frac{1}{k!}\begin{array}{ccc}i_2& \mathrm{}& i_{k+1}\\ 0& \mathrm{}& 0\end{array}=0.$$
If $`s<i_1+j_1`$ than
$$P_1\left[\begin{array}{cccc}i_1& i_2& \mathrm{}& i_{k+1}\\ j_1& 0& \mathrm{}& 0\end{array}\right]=\frac{1}{k!}C_s^{j_1}\begin{array}{ccc}i_2& \mathrm{}& i_{k+1}\\ 0& \mathrm{}& 0\end{array}AC_{s(i_1+j_1)}^0=$$
$$=\frac{1}{k!}\begin{array}{ccc}i_2& \mathrm{}& ,i_{k+1}\\ 0& \mathrm{}& 0\end{array}P_s\left[\begin{array}{c}i_1\\ j_1\end{array}\right].$$
Thus the lemma is proved for $`m=1`$.
Prove the lemma for $`m=N`$ considering that it is proved for $`m<N`$. Then
$$P_s\left[\begin{array}{cccccc}i_1& \mathrm{}& i_m& i_{m+1}& \mathrm{}& i_{m+k}\\ j_1& \mathrm{}& j_m& 0& \mathrm{}& 0\end{array}\right]=$$
$$=\underset{\left(\begin{array}{ccc}\alpha _1& \mathrm{}& \alpha _m\\ \beta _1& \mathrm{}& \beta _m\end{array}\right)\left[\begin{array}{ccc}i_1& \mathrm{}& i_m\\ j_1& \mathrm{}& j_m\end{array}\right]}{}(\frac{1}{(m+k)!}C_s^{\beta _1+\mathrm{}+\beta _m}\frac{(\beta _1+\mathrm{}+\beta _m)!}{\beta _1!\mathrm{}\beta _m!}C_{m+k}^k$$
$$\begin{array}{ccc}i_{m+1}& \mathrm{}& i_{m+k}\\ 0& \mathrm{}& 0\end{array}$$
$$\underset{q=1}{\overset{m}{}}P_s\left(\begin{array}{ccc}\alpha _1& \mathrm{}& \alpha _q\\ \beta _1& \mathrm{}& \beta _q\end{array}\right)\frac{1}{(m+kq)!}C_{s(\alpha _1+\mathrm{}+\alpha _q+\beta _1+\mathrm{}+\beta _q)}^{\beta _{q+1}+\mathrm{}+\beta _m}\frac{(\beta _{q+1}+\mathrm{}+\beta _m)!}{\beta _{q+1}!\mathrm{}\beta _m!}C_{m+kq}^k$$
$$\begin{array}{ccc}i_{m+1}& \mathrm{}& i_{m+k}\\ 0& \mathrm{}& 0\end{array})=P_s\left[\begin{array}{ccc}i_1& \mathrm{}& i_m\\ j_1& \mathrm{}& j_m\end{array}\right]\frac{1}{k!}\begin{array}{ccc}i_{m+1}& \mathrm{}& i_{m+k}\\ 0& \mathrm{}& 0\end{array}$$
$$P_s\left[\begin{array}{ccc}i_1& \mathrm{}& i_m\\ j_1& \mathrm{}& j_m\end{array}\right]C_{s(i_1+\mathrm{}+i_m+j_1+\mathrm{}+j_m)}^0\frac{1}{k!}\begin{array}{ccc}i_{m+1}& \mathrm{}& i_{m+k}\\ 0& \mathrm{}& 0\end{array}=$$
$$=\{\begin{array}{cc}0,\text{if }si_1+\mathrm{}+i_m+j_1+\mathrm{}+j_m,\hfill & \\ P_s\left[\begin{array}{ccc}i_1& \mathrm{}& i_m\\ j_1& \mathrm{}& j_m\end{array}\right]\frac{1}{k!}\begin{array}{ccc}i_{m+1}& \mathrm{}& i_{m+k}\\ 0& \mathrm{}& 0\end{array},\hfill & \text{if }s<i_1+\mathrm{}+i_m+j_1+\mathrm{}+j_m\text{.}\mathrm{}\hfill \end{array}$$
3. Equations for the Bacher-Akhiezer function
Consider the KP hierarchy. This is a condition of compatibility of the infinite system of the differential equations
$$\frac{\psi }{x_n}=L_n\psi ,$$
$`3.1`$
where
$$L_n=\frac{^n}{x_1^n}+\underset{i=2}{\overset{n}{}}B_n^i(x)\frac{^{ni}}{^{ni}x_1},$$
and $`\psi `$ is a function of type
$$\psi (x,k)=\text{exp}(\underset{j=1}{\overset{\mathrm{}}{}}x_jk^j)(1+\underset{i=1}{\overset{\mathrm{}}{}}\xi _ik^i),$$
(here $`kC`$ belong to some neighbourhood of $`\mathrm{}`$ and $`x=(x_1,x_2,\mathrm{})`$ — is a finite sequence).
Put us
$$_i=\frac{}{x_i},=_1.$$
A direct calculation gives
Lemma 3.1. Conditions of compatibility of (1) are
$$B_s^t=\underset{i=1}{\overset{t1}{}}C_s^i^i\xi _{ti}\underset{j=2}{\overset{t1}{}}B_s^j\underset{i=0}{\overset{tj1}{}}C_{sj}^i^i\xi _{tij},$$
$`3.2`$
$$_n\xi _i=\underset{j=1}{\overset{n+i1}{}}C_n^j^j\xi _{i+nj}+\underset{k=2}{\overset{n}{}}B_n^k\underset{j=0}{\overset{nk}{}}C_{nk}^j^j\xi _{i+njk}.\mathrm{}$$
$`3.3`$
In this case $`\psi `$ is called a Bacher-Akhiezer function.
Consider now the function
$$\mathrm{ln}\psi (x,k)=\underset{j=1}{\overset{\mathrm{}}{}}x_jk^j+\underset{j=1}{\overset{\mathrm{}}{}}\eta _jk^j,$$
where
$$\xi _j=\underset{n=1}{\overset{\mathrm{}}{}}\frac{1}{n!}\underset{i_1+\mathrm{}+i_n=j}{}\eta _{i_1}\mathrm{}\eta _{i_n}.$$
Lemma 3.2. Let $`2ts`$. Then
$$B_s^t=\underset{n=1}{\overset{\mathrm{}}{}}P_s\left(\begin{array}{ccc}i_1& \mathrm{}& i_n\\ j_1& \mathrm{}& j_n\end{array}\right)^{j_1}\eta _{i_1}\mathrm{}^{i_n}\eta _{i_n},$$
where the second sum is taken by all matrices $`\left(\begin{array}{ccc}i_1& \mathrm{}& i_n\\ j_1& \mathrm{}& j_n\end{array}\right)`$ such that $`i_m1,j_m1`$ and $`i_1+\mathrm{}+i_n+j_1+\mathrm{}+j_n=t`$.
Proof: An induction by $`t`$. For $`t=2`$ according to (3.2), $`B_s^2=s\xi _1=P_s\left(\genfrac{}{}{0pt}{}{1}{1}\right)\eta _1`$. Prove the lemma for $`t=N`$ considering that it is proved for $`t<N`$. According to (3.2)
$$B_s^t=\underset{i=1}{\overset{t1}{}}C_s^i^i(^{\mathrm{}}_{n=1}\frac{1}{n!}_{i_1+\mathrm{}+i_n=ti}\eta _{i_1}\mathrm{}\eta _{i_n})+$$
$$+\underset{j=2}{\overset{t1}{}}\left(\underset{n=1}{\overset{\mathrm{}}{}}\underset{i_1+\mathrm{}+j_n=j}{}P_s\left(\begin{array}{ccc}i_1& \mathrm{}& i_n\\ j_1& \mathrm{}& j_n\end{array}\right)^{j_1}\eta _{i_1}\mathrm{}^{j_n}\eta _{i_n}\right)$$
$$\left(\underset{i=0}{\overset{tj1}{}}C_{sj}^i^j\left(\underset{n=1}{\overset{\mathrm{}}{}}\frac{1}{n!}\underset{i_1+\mathrm{}+i_n=tij}{}\eta _{i_1}\mathrm{}\eta _{i_n}\right)\right)=$$
$$=\underset{n=1}{\overset{\mathrm{}}{}}(\frac{1}{n!}C_s^{j_1+\mathrm{}+j_n}\frac{(j_1+\mathrm{}+j_n)!}{j_1!\mathrm{}j_n!}$$
$$\underset{q=1}{\overset{n1}{}}P_s\left(\begin{array}{ccc}i_1& \mathrm{}& i_q\\ j_1& \mathrm{}& j_q\end{array}\right)\frac{1}{(nq)!}C_{s(i_1+\mathrm{}i_q+j_1+\mathrm{}+j_q)}^{j_{q+1}\mathrm{}j_n}\frac{(j_{q+1}+\mathrm{}j_n)!}{j_{q+1}!\mathrm{}j_n!})^{j_1}\eta _{i_1}\mathrm{}^{j_n}\eta _{i_n}=$$
$$=\underset{n=1}{\overset{\mathrm{}}{}}P_s\left(\begin{array}{ccc}i_1& \mathrm{}& i_n\\ j_1& \mathrm{}& j_n\end{array}\right)^{j_1}\eta _{i_1}\mathrm{}^{j_n}\eta _{i_n},$$
where the second sums are taken by all matrices $`\left(\begin{array}{ccc}i_1& \mathrm{}& i_n\\ j_1& \mathrm{}& j_n\end{array}\right)`$ such that $`i_1+\mathrm{}+i_n+j_1+\mathrm{}+j_n=t`$, $`i_m1,j_m0`$. According to lemma 2.1 it is possible to consider that in the last sum $`j_m>0`$ for all $`m`$. $`\mathrm{}`$
Lemma 3.3.
$$_s\eta _r=\underset{n=1}{\overset{\mathrm{}}{}}P_s\left(\begin{array}{ccc}i_1& \mathrm{}& i_n\\ j_1& \mathrm{}& j_n\end{array}\right)^{j_1}\eta _{i_1}\mathrm{}^{j_n}\eta _{i_n},$$
where the second sum is taken by all matrices $`\left(\begin{array}{ccc}i_1& \mathrm{}& i_n\\ j_1& \mathrm{}& j_n\end{array}\right)`$ such that $`i_m1,j_m1`$ and $`i_1+\mathrm{}+i_n+j_1+\mathrm{}+j_n=r+s`$.
Proof: An induction by $`r`$. According to (3.3) and lemma 3.2 for $`r=1`$
$$_s\eta _1=_s\xi _1=\underset{j=1}{\overset{\mathrm{}}{}}C_s^j^j\xi _{s+1j}+\underset{k=2}{\overset{s}{}}B_s^k\underset{j=0}{\overset{sk}{}}C_{sk}^j^j\xi _{1+sjk}=$$
$$=\underset{j=1}{\overset{s}{}}C_s^j^j\xi _{s+1j}\underset{k=2}{\overset{\mathrm{}}{}}\left(\underset{n=1}{\overset{\mathrm{}}{}}\underset{i_1+\mathrm{}+j_n=k}{}P_s\left(\begin{array}{ccc}i_1& \mathrm{}& i_n\\ j_1& \mathrm{}& j_n\end{array}\right)^{j_1}\eta _{i_1}\mathrm{}^{j_n}\eta _{i_n}\right)$$
$$\left(\underset{j=0}{\overset{sk}{}}C_{sk}^j^j\xi _{1+sjk}\right)=\underset{n=1}{\overset{\mathrm{}}{}}P_s\left(\begin{array}{ccc}i_1& \mathrm{}& i_n\\ j_1& \mathrm{}& j_n\end{array}\right)^{j_1}\eta _{i_1}\mathrm{}^{j_n}\eta _{i_n},$$
where the second sum in the last formula is taken by all matrices $`\left(\begin{array}{ccc}i_1& \mathrm{}& i_n\\ j_1& \mathrm{}& j_n\end{array}\right)`$ such that $`i_1+\mathrm{}+i_n+j_1+\mathrm{}+j_n=s+1,i_m1,j_m0`$. According to lemma 1 in this sum it is sufficient consider only matrices, where $`j_m>0`$ for all $`m`$.
Prove now the lemma for $`r=N`$, considering that it is proved for $`r<N`$. According to (3.3)
$$_s\left(\underset{n=1}{\overset{\mathrm{}}{}}\frac{1}{n!}\underset{i_1+\mathrm{}+i_n=r}{}\eta _{i_1}\mathrm{}\eta _{i_n}\right)=\underset{j=1}{\overset{s+r1}{}}C_s^j^j\xi _{s+rj}+\underset{k=2}{\overset{\mathrm{}}{}}B_s^k\underset{j=0}{\overset{sk}{}}C_{sk}^j^j\xi _{r+sjk}.$$
Thus according to lemma 3.2, lemma 2.1 and inductive hypothesis,
$$_s\eta _r=\underset{j=1}{\overset{\mathrm{}}{}}C_s^j^j\left(\underset{n=1}{\overset{\mathrm{}}{}}\frac{1}{n!}\underset{i_1+\mathrm{}+i_n=s+rj}{}\eta _{i_1}\mathrm{}\eta _{i_n}\right)+$$
$$+\underset{k=2}{\overset{\mathrm{}}{}}\left(\underset{i_1+\mathrm{}+j_n=k}{}P_s\left(\begin{array}{ccc}i_1& \mathrm{}& i_n\\ j_1& \mathrm{}& j_n\end{array}\right)^{j_1}\eta _{i_1}\mathrm{}^{j_n}\eta _{i_n}\right)$$
$$\underset{j=0}{\overset{sk}{}}C_{sk}^j^j\left(\underset{n=1}{\overset{\mathrm{}}{}}\frac{1}{n!}\underset{i_1+\mathrm{}+i_n=r+sjk}{}\eta _{i_1}\mathrm{}\eta _{i_n}\right)_s\left(\underset{n=2}{\overset{\mathrm{}}{}}\frac{1}{n!}\underset{i_1+\mathrm{}+i_n=r}{}\eta _{i_1}\mathrm{}\eta _{i_n}\right)=$$
$$=\underset{n=1}{\overset{\mathrm{}}{}}P_s\left(\begin{array}{ccc}i_1& \mathrm{}& i_n\\ j_1& \mathrm{}& j_n\end{array}\right)^{j_1}\eta _{i_1}\mathrm{}^{j_n}\eta _{i_n},$$
where the second sum in the last formula is taken by all matrices $`\left(\begin{array}{ccc}i_1& \mathrm{}& i_n\\ j_1& \mathrm{}& j_n\end{array}\right)`$ such that $`i_1+\mathrm{}+i_n+j_1+\mathrm{}+j_n=s+r`$, $`i_m1,j_m1`$. $`\mathrm{}`$
4. KP hierarchy
According to \[DKJM\] the Bacher-Akhiezer function $`\psi `$ is
$$\psi (x,k)=\text{exp}(x_jk^j)\frac{\tau (x_1k^1,x_2\frac{1}{2}k^2,x_3\frac{1}{3}k^3,\mathrm{})}{\tau (x_1,x_2,x_3,\mathrm{})}$$
for some function $`\tau (x_1,x_2,\mathrm{})`$. By analogy of \[N1\] this gives a possibility to describe the KP hierarchy as an infinite system of differential equations on $`v(x,k)=\mathrm{ln}\tau (x,k)`$. Really
$$\underset{j=1}{\overset{\mathrm{}}{}}\eta _jk^j=\mathrm{ln}\psi (x,k)\underset{j=1}{\overset{\mathrm{}}{}}x_jk^j=v(x_1k^1,x_2\frac{1}{2}k^2,\mathrm{})+v(x)=$$
$$=\underset{n=1}{\overset{\mathrm{}}{}}\underset{i_1+\mathrm{}+i_n=j}{}\frac{(1)^{n+1}}{n!i_1\mathrm{}i_n}_{i_1}\mathrm{}_{i_n}v(x)k^j.$$
Therefore
$$\eta _r=\underset{n=1}{\overset{\mathrm{}}{}}\underset{i_1+\mathrm{}+i_n=r}{}\frac{(1)^{n+1}}{n!i_1\mathrm{}i_n}_{i_1}\mathrm{}_{i_n}v.$$
$`4.1`$
Theorem 4.1. There exist universal rational coefficients
$$R_r\left(\begin{array}{ccc}s_1& \mathrm{}& s_n\\ t_1& \mathrm{}& t_n\end{array}\right),R_{ij}\left(\begin{array}{ccc}s_1& \mathrm{}& s_n\\ t_1& \mathrm{}& t_n\end{array}\right)$$
such that
$$\eta _r=\frac{1}{r}_rv+\underset{n=1}{\overset{\mathrm{}}{}}R_r\left(\begin{array}{ccc}s_1& \mathrm{}& s_n\\ t_1& \mathrm{}& t_n\end{array}\right)_{s_1}^{t_1}v\mathrm{}_{s_n}^{t_n}v,$$
$`4.2`$
$$_i_jv=\underset{n=1}{\overset{\mathrm{}}{}}R_{ij}\left(\begin{array}{ccc}s_1& \mathrm{}& s_n\\ t_1& \mathrm{}& t_n\end{array}\right)_{s_1}^{t_1}v\mathrm{}_{s_n}^{t_n}v,$$
$`4.3`$
where the second sums are taken by all matrices $`\left(\begin{array}{ccc}s_1& \mathrm{}& s_n\\ t_1& \mathrm{}& t_n\end{array}\right)`$ such that $`s_m,t_m1`$, and the sum $`s_1+\mathrm{}+s_n+t_1+\mathrm{}+t_n`$ is equal $`r`$ for (4.2) and $`i+j`$ for (4.3).
Proof: An induction by $`k`$ and $`i+j`$. For $`i+j=2`$ the theorem is obviously. For $`r=1`$ it follows from (4). Prove the theorem for $`i+j=N`$ and $`r=N1`$, considering that it is proved for $`i+j<N`$ and $`r<N1`$. Later we consider that $`s_m,t_m1`$ and $`\sigma _n=s_1+\mathrm{}+s_n+t_1+\mathrm{}+t_n`$. Then according to (4.1) and (4.3)
$$\eta _r=\frac{1}{r}_rv+\underset{n=2}{\overset{\mathrm{}}{}}\underset{s_1+\mathrm{}+s_n=r}{}\frac{(1)^{n+1}}{n!s_1\mathrm{}s_n}_{s_1}\mathrm{}_{s_n}v(x)=$$
$$=\frac{1}{r}_rv+\underset{n=1}{\overset{\mathrm{}}{}}\underset{\sigma _n=r}{}R_r\left(\begin{array}{ccc}s_1& \mathrm{}& s_n\\ t_1& \mathrm{}& t_n\end{array}\right)_{s_1}^{t_1}v\mathrm{}_{s_n}^{t_n}v.$$
Thus according to (4.2), (4.3) and lemma 3.3
$$\frac{1}{j}_i_jv=_i\eta _j_i\left(\underset{n=1}{\overset{\mathrm{}}{}}\underset{\sigma _n=j}{}R_j\left(\begin{array}{ccc}s_1& \mathrm{}& s_n\\ t_1& \mathrm{}& t_n\end{array}\right)_{s_1}^{t_1}v\mathrm{}_{s_n}^{t_n}v\right)=$$
$$=\underset{n=1}{\overset{\mathrm{}}{}}\underset{\sigma _n=i+j}{}P_i\left(\begin{array}{ccc}s_1& \mathrm{}& s_n\\ t_1& \mathrm{}& t_n\end{array}\right)^{t_1}\eta _{s_1}\mathrm{}^{t_n}\eta _{s_n}$$
$$_i\left(\underset{n=1}{\overset{\mathrm{}}{}}\underset{\sigma _n=j}{}R_j\left(\begin{array}{ccc}s_1& \mathrm{}& s_n\\ t_1& \mathrm{}& t_n\end{array}\right)^{t_1}_{s_1}v\mathrm{}^{t_n}_{s_n}v\right)=$$
$$=\underset{n=1}{\overset{\mathrm{}}{}}\underset{s_1+\mathrm{}+s_n+n=i+j}{}P_i\left(\begin{array}{ccc}s_1& \mathrm{}& s_n\\ 1& \mathrm{}& 1\end{array}\right)(\frac{1}{s_1}_{s_1}v)\mathrm{}(\frac{1}{s_n}_{s_n}v)+$$
$$+\underset{n=1}{\overset{\mathrm{}}{}}\underset{\sigma _n=i+j,t_1+\mathrm{}+t_n>n}{}R_{ij}\left(\begin{array}{ccc}s_1& \mathrm{}& s_n\\ t_1& \mathrm{}& t_n\end{array}\right)_{s_1}^{t_1}v\mathrm{}_{s_n}^{t_n}v.\mathrm{}$$
Remark 4.1. The system (4.3) was at first deduced in \[DN\]. The set of its solution bijectively correspond to the set of solutions of (3.1). Up to a constant formal solution of (4.3) is defined by an infinite set of functions of one variable $`f_i(x_1)=_iv|_{x_2=x_3=\mathrm{}=0}(i=1,2,\mathrm{})`$.
Remark 4.2. The algorithm described in the proof ot theorem 4.1 gives an algorithm for calculation of all rational constants $`R_{ij}\left(\begin{array}{ccc}s_1& \mathrm{}& s_m\\ t_1& \mathrm{}& t_m\end{array}\right)`$. The first equations of hierarchy (4.3) are:
$$_2^2v=\frac{4}{3}_3v\frac{1}{3}^4v+2(^2v)^2,$$
$`4.4`$
$$_3_2v=\frac{3}{2}_4v\frac{3}{2}_2^3v+3_2v^2v,$$
$$_3^2v=\frac{9}{5}_5v_3^3v+\frac{1}{5}^6v+3_3v^2v+\frac{9}{4}(_2v)^23^4v^2v\frac{9}{4}(^3v)^2+3(^2v)^3.$$
The equation (4.4) is KP equation twice integrated over $`x_1`$.
Theorem 4.2. If $`_{i=1}^m(t_i+1)1(\text{mod 2})`$, then $`R_{ij}\left(\begin{array}{ccc}s_1& \mathrm{}& s_m\\ t_1& \mathrm{}& t_m\end{array}\right)=0`$.
Proof: The equations of KP–hierarchy are equivalents of equations on function $`\tau (x)`$ \[DKJM\]. All these equations can be written simply by means of the ”bilinear Hirota operators”. We recall the definition of them. If $`f(x)`$ is a function of one variable, then for any polynomial (or power series) $`Q`$ the action of the Hirota operator $`Q(D_x)f(x)f(x)`$ is defined by
$$Q(D_x)f(x)f(x)=Q(_y)[f(x+y)f(xy)]_{y=0}.$$
For functions of several variables the definition is similar. The generating function for the equations of the KP hierarchy has the form
$$\underset{j=0}{\overset{\mathrm{}}{}}p_j(2y)p_{j+1}(\stackrel{~}{D})\mathrm{exp}\left(\underset{i=1}{\overset{\mathrm{}}{}}y_iD_i\right)\tau \tau =0,$$
$`4.5`$
where $`y=(y_1,y_2,\mathrm{},)`$ are auxiliary independent variables, $`\stackrel{~}{D}=(D_1,2^1D_2,3^1D_3,\mathrm{})`$, $`D_j`$ is the Hirota operator in the variable $`x_j`$ and $`p_j`$ are the Schur polynomials defined from the following expansion:
$$\mathrm{exp}\left(\underset{j=1}{\overset{\mathrm{}}{}}x_jk^j\right)=\underset{j=0}{\overset{\mathrm{}}{}}k^jp_j(x_1,\mathrm{},x_j).$$
All monomials of odd degree give trivial Hirota operators. Therefore if $`\tau (x)`$ is a solution of (4.5), then $`\stackrel{~}{\tau }(x)=\tau (x)`$ is also solution of (4.5). Moreover, according to \[DN\] a function $`\tau `$ is a solution of the system (4.5) if and only if $`v=\mathrm{ln}(\tau )`$ is a solution of the system (4.3). Thus, $`v(x)`$ is a formal solution of the system (4.3), if and only if $`\stackrel{~}{v}(x)=v(x)`$ is a formal solution of the system (4.3). This is equivalent of the affirmation of theorem 4.2. $`\mathrm{}`$.
5. Gelfand–Dikii hierarchy
According to \[S\] the set of solution of $`n`$–Gelfand–Dikii hierarchy bijectively correspond to the set of nondepending from $`x_n`$ solutions of KP hierarchy. In this case according to theorem 4.1
$$0=_m_nv=\frac{mn}{m+n1}_{n+m1}v+\underset{m=1}{\overset{\mathrm{}}{}}R_{mn}\left(\begin{array}{ccc}s_1& \mathrm{}& s_m\\ t_1& \mathrm{}& t_m\end{array}\right)_{s_1}^{t_1}v\mathrm{}_{s_m}^{t_m}v,$$
where $`1s_jn+m2,t_j1`$. This gives recurrence formulas expressing $`_kv`$ for $`k>n`$ via $`_rv`$ for $`r<n`$. Thus we have relations
$$_{n+r}v=\underset{m=1}{\overset{\mathrm{}}{}}N_{1(n+1)}^m\left(\begin{array}{ccc}s_1& \mathrm{}& s_m\\ t_1& \mathrm{}& t_m\end{array}\right)_{s_1}^{t_1}v\mathrm{}_{s_m}^{t_m}v,$$
$`5.1`$
where $`t_j1,s_j<n,_{j=1}^m(s_j+t_j)=n+r+1`$.
Example. For $`n=2`$ the system (5.1) passes to $`KV`$ hierarchy.
Compering the systems (5.1) and (4.3) we find the system
$$_i_jv=\underset{m=1}{\overset{\mathrm{}}{}}N_{i_j}^m\left(\begin{array}{ccc}s_1& \mathrm{}& s_m\\ t_1& \mathrm{}& t_m\end{array}\right)_{s_1}^{t_1}v\mathrm{}_{s_m}^{t_m}v,$$
$`5.2`$
where $`i,j1,1s_\alpha n1,t_\alpha 1`$, $`_{i=1}^m(s_\alpha +t_\alpha )=i+j`$ and $`N_{ij}^m\left(\begin{array}{ccc}s_1& \mathrm{}& s_m\\ t_1& \mathrm{}& t_m\end{array}\right)`$ are some universal rational coefficients.
Examples.
1. For $`n=3`$ the first equation from the system (5.2) is the Boussinesq equation
$$_2^2v=\frac{1}{3}^4v+2(^2v)^2.$$
2. For $`n=4`$ the first equations of the system (5.2) are
$$_2^2v=\frac{4}{3}_3v\frac{1}{3}^4v+2(^2v)^2,$$
$$_3_2v=\frac{3}{2}_2^3v+3_2v^2v,$$
$$_3^2v=\frac{1}{4}_3^3v+\frac{1}{8}^6v+\frac{9}{8}(_2v)^2\frac{9}{8}(^3v)^2\frac{9}{4}^4v^2v+3(^2v)^3.\mathrm{}$$
Theorem 5.1. The Gelfand–Dikii hierarchy is equivalent to a system of differential equations in a form
$$_{i_1}\mathrm{}_{i_k}v=\underset{m=1}{}\underset{i_1\mathrm{}i_k}{}N_{i_1\mathrm{}i_k}^m\left(\begin{array}{ccc}s_1& \mathrm{}& s_m\\ t_1& \mathrm{}& t_m\end{array}\right)_{s_1}^{t_1}v\mathrm{}_{s_m}^{t_m}v,$$
$`5.3`$
where $`k2,t_j1,s_j<n,_{j=1}^ki_j=_{j=1}^m(s_j+t_j),_{j=1}^mt_jm+k2`$ $`k+m+_{j=1}^mt_j0(\text{mod}2)`$.
Proof: For $`k=2`$ the equations (5.3) coincide with the equations (5.2). For $`k>2`$ the equations (5.3) are received from equations for $`k1`$ by differentiation by $`_{i_k}`$ and replacing $`_i_{i_k}v`$ by (5.2). This gives monomials $`_{s_1}^{t_1}\mathrm{}_{s_m}^{t_m}v`$, where the $`_{j=1}^ki_j=_{j=1}^m(s_j+t_j)`$ $`_{j=1}^mt_jm+k2`$. The condition $`k+m+_{j=1}^mt_j0(\text{mod}2)`$ follows from theorem 4.2. $`\mathrm{}`$
Remark 5.1. Structure of the system (5.3) such that its formal solutions are defined up to constant by arbitrary set of $`n1`$ series from one variable $`f_i(x_1)=_iv|_{x_2=x_3=\mathrm{}=0}(i=1,\mathrm{},n1)`$.
Remark 5.2. Coefficients $`N_{i_1\mathrm{}i_k}^m\left(\begin{array}{ccc}s_1& \mathrm{}& s_m\\ t_1& \mathrm{}& t_m\end{array}\right)`$ are rational constants. The constructions, described in the proofs of theorems 4.1 5.1, give recurrent formulas for its calculation.
6. Witten solution of the Gelfand–Dikii hierarchy
Follow by Witten \[W\] let us consider the space $`M_{g,s}`$ of Riemann surfaces of genus $`g`$ with $`s`$ punctures. Correspond to any puncture a pair $`(k_i,m_i)`$, where $`1k_i<n,m_i0`$. Witten \[W\] connects with the set $`\{(k_i,m_i)|i=1,\mathrm{},s\}`$ a number (correlator) $`<\underset{k,m}{}\tau _{k,m}^{d_{k,m}}>_g`$, where $`d_{k,m}`$ is the number of pairs $`(k_i,m_i)`$, that equal to $`(k,m)`$. The number $`<\underset{k,m}{}t_{k,m}^{d_{k,m}}>_g`$ is equals to the value of some class of cohomology on a compactification of a space of $`n`$–spin bundles over $`PM_{g,s}`$ \[W\]. Put us
$$F^g(t_{1,0},t_{1,1},\mathrm{})=\underset{d_{k,m}}{}<\underset{k,m}{}\tau _{k,m}^{d_{k,m}}>_g\underset{k,m}{}\frac{t_{k,m}^{d_{k,m}}}{d_{k,m}!}.$$
According to the Witten conjecture the serias $`F=\underset{g=0}{\overset{\mathrm{}}{}}F^g`$ after the change $`t_{k,m}(mn+k)x_{mn+k}`$ pass to a formal solution $`v`$ of the system (5.3) satisfying the equation
$$v=\frac{1}{2}\underset{i+j=n}{}ijx_ix_j+\underset{i=1}{\overset{\mathrm{}}{}}(i+n)x_{i+n}_iv,$$
$`6.1`$
$$0=v(0)=_iv(0)(i=1,2,\mathrm{}).$$
The single such solution $`W`$ we call the Witten solution of the Gelfand–Dikii hierarchy.
Theorem 6.1. The Witten solution of the Gelfand–Dikii hierarchy is $`W=\underset{g=0}{\overset{\mathrm{}}{}}W^g`$, where
$$W^g(x_1,x_2,\mathrm{})=$$
$$=\underset{k=2}{\overset{\mathrm{}}{}}\underset{i_1+\mathrm{}+i_k=(n+1)(2g2+k)}{}\frac{(n1)^{2g2+k}}{k!}N_{i_1\mathrm{}i_k}^{2g2+k}\left(\begin{array}{ccc}n1& \mathrm{}& n1\\ 2& \mathrm{}& 2\end{array}\right)x_{i_1}\mathrm{}x_{i_k}.$$
Moreover the function $`W^g`$ is a quasihomogeous series of degree $`(1g)(2+\frac{2}{n})`$ by $`x_i`$ of degrees $`1+\frac{1}{n}\frac{i}{n}`$.
Proof: Compatibility of the equations (5.3) (6.1) follows from \[AM\], where this solution is represented in a form of matrix integral. According to (6.1) $`_iW|_{x_2=x_3=\mathrm{}=0}=\delta _{n1,i}(n1)x_1`$. These conditions and the equations (5.3) uniquely determine all functions $`f_{i_1\mathrm{}i_k}(x_1)=_{i_1}\mathrm{}_{i_k}W|_{x_2=x_3=\mathrm{}=0}`$. According to theorem 5.1 if $`f_{i_1\mathrm{}i_k}(0)0`$, that $`(i_1+\mathrm{}+i_k)0(\text{mod}(n+1))`$ and
$$f_{i_1\mathrm{}i_k}(0)=(n1)^mN_{i_1\mathrm{}i_k}^m\left(\begin{array}{ccc}n1& \mathrm{}& n1\\ 2& \mathrm{}& 2\end{array}\right),$$
where $`m=\frac{i_1+\mathrm{}+i_k}{n+1}`$. From this theorem follow also that $`mk2`$ $`k+m0(\text{mod}2)`$. Therefore $`m=2g+k2`$, where $`g0`$ is a natural number. Thus, $`W=W^g`$, where
$$W^g(x_1,x_2,\mathrm{})=\underset{k=2}{\overset{\mathrm{}}{}}\underset{i_1+\mathrm{}+i_k=(n+1)(2g+k2)}{}\frac{1}{k!}f_{i_1\mathrm{}i_k}(0)x_{i_1}\mathrm{}x_{i_k}=$$
$$=\underset{k=2}{\overset{\mathrm{}}{}}\underset{i_1+\mathrm{}+i_k=(n+1)(2g+k2)}{}\frac{1}{k!}(n1)^{2g2+k}N_{i_1\mathrm{}i_k}^{2g2+k}\left(\begin{array}{ccc}n1& \mathrm{}& n1\\ 2& \mathrm{}& 2\end{array}\right)x_{i_1}\mathrm{}x_{i_k}.$$
The quasihomogeneity of the series $`W^g`$ follows from $`i_1+\mathrm{}+i_k=(n+1)(2g2+k)`$. $`\mathrm{}`$
Corollary 6.1. The Witten solution of the Gelfand–Dikii hierarchy has a representation by the sum of quasihomogeneous series $`W^g`$ of the same degrees that $`F^g`$.
Proof: According to \[W\] $`F^g`$ is a quasihomogeneous series of degree $`(1g)(2+\frac{2}{n})`$ by $`t_{k,m}`$ of degree $`1+\frac{1}{n}k\frac{m}{n}`$. $`\mathrm{}`$
Theorem 6.2. The functions $`W`$ and $`W^0`$ coincide on the set $`L_0=(x_1,x_2,\mathrm{},x_{n1},0,0,\mathrm{})`$.
Proof: According to (6.1) $`^r_{\mathrm{}}W=0`$ on the set $`L_0`$ if $`\mathrm{}<n1,r>1`$, or $`\mathrm{}=n1,r>2`$. Besides according to (5.1)
$$_{n+\mathrm{}}W=\underset{m=1}{\overset{\mathrm{}}{}}N_{1(n+\mathrm{})}^m\left(\begin{array}{ccc}s_1& \mathrm{}& s_m\\ t_1& \mathrm{}& t_m\end{array}\right)_{s_1}^{t_1}W\mathrm{}_{s_m}^{t_m}W,$$
where $`\underset{i=1}{\overset{m}{}}(s_i+t_i)=n+\mathrm{}+1`$. Thus, if $`\mathrm{}<n`$ and $`_{n+\mathrm{}}W0`$, that all numbers $`t_j`$ are less than 3 and among them is not two number more than 1. But if among the numbers $`1t_1,\mathrm{},t_m2`$ there is exactly one $`t_j=2`$ then $`N_{1(n+\mathrm{})}^m\left(\begin{array}{ccc}s_1& \mathrm{}& s_m\\ t_1& \mathrm{}& t_m\end{array}\right)=0`$ by theorem 5.1. Thus,
$$_{n+\mathrm{}}W=\underset{m=1}{\overset{\mathrm{}}{}}N_{1(n+\mathrm{})}^m\left(\begin{array}{ccc}s_1& \mathrm{}& s_m\\ 1& \mathrm{}& 1\end{array}\right)_{s_1}W\mathrm{}_{s_m}W.$$
Moreover, according (6.1) $`_{\mathrm{}}W=_{n+\mathrm{}}W`$ and $`_sW=s(ns)x_{ns}`$ on the set $`L_0`$ . Thus from the equality $`\underset{i=1}{\overset{m}{}}(s_i+t_i)=n+\mathrm{}+1`$ follows that $`_{\mathrm{}}W|_{L_0}`$ is a quasihomogeneous polynomial of degree $`2+\frac{2}{n}(1+\frac{1}{n}\frac{\mathrm{}}{n})=1+\frac{1}{n}+\frac{\mathrm{}}{n}`$. Thus, $`W|_{L_0}`$ is a quasihomogeneous polynomial of degree $`2+\frac{2}{n}`$ $`W|_{L_0}=W^0|_{L_0}.\mathrm{}`$
Corollary 6.2. $`F(x_1,x_2,\mathrm{},x_{n1},0,0,\mathrm{})=W(x_1,\frac{x_2}{2},\frac{x_3}{3},\mathrm{},\frac{x_{n1}}{n1},0,0,\mathrm{})`$.
Proof: According to \[W\] $`F^0(x_1,x_2,\mathrm{})=W^0(x_1,\frac{x_2}{2},\frac{x_3}{3},\mathrm{})`$ and $`F(x_1,x_2,\mathrm{},x_{n1},0,0,\mathrm{})=F^0(x_1,x_2,\mathrm{},x_{n1},0,0,\mathrm{})`$. Thus the theorem 6.2 imply corollary 6.2. $`\mathrm{}`$
Remark 6.1. According to \[DVV, Kr\] the function $`W^0|_{L_o}`$ is the potential of Frobenius structure on the space of versal deformations of the singularity $`A_n`$. By theorem 6.2 we have $`W^0|_{L_0}=W|_{L_0}`$. A simple algorithm of calculation of this function describes in \[N2\].
References
\[AM\] M.Adler, P.van Moerbeke, A matrix integral solution to two–dimensional $`W_p`$–gravity. Commun. Math. Phys. 147 (1992), 25-56.
\[DKJM\] E.Date, M.Kashiwara, M.Jimbo, T.Miwa, Transformation groups for soliton equation, Proceedings of RIMS Symposium on Non-Linear Integrable System. Singapore: World Science Publ. Co., 1983, 39-119.
\[DN\] B.A.Dubrovin, S.M.Natanzon, Real theta-function solutions of the Kadomtsev–Petviashvili equation. Math. USSR Irvestiya, 32:2 (1989), 269-288.
\[DVV\] R.Dijkgraaf, E.Verlinde, H.Verlinde, Topological strings in $`d<1`$. Nucl. Phys., B 352 (1991), 59.
\[GD\] I.M.Gelfand, L.A.Dikii, The resolvent and Hamiltonian systems, Funct. Anal. Appl. 11 (2) (1977), 93-105.
\[Ko\] M.Kontsevich, Intersection theory on the moduli space of curves and the matrix airy function. Commun. Math. Phys. 147 (1992), 1-23.
\[Kr\] I.Krichever, The dispersionless Lax equations and topological minimal models. Commun. Math. Phys. 143 (1992), 415-429.
\[N1\] S.M.Natanzon, Real nonsingular finite zone solutions of solution equations. Amer. Math. Soc. Transl. (2). V.170 (1995), 153-183
\[N2\] S.M.Natanzon, Formulas for $`A_n`$ and $`B_n`$–solutions of WDVV equations, hep-th/9904103.
\[S\] M.Sato, Soliton equations and universal Grassmann manifold. Math. Lect. Notes Ser., Vol.18, Sophia University, Tokyo (1984).
\[W\] E.Witten, Algebraic geometry associated with matrix models of two dimensional gravity, Topological models in modern mathematics (Stony Brook, NY, 1991), Publish or Perish, Houston, TX 1993, 235-269.
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# CANONICAL TRANSFORMATION OF THE HUBBARD MODEL AND W=0 PAIRING: COMPARISON WITH EXACT DIAGONALIZATION RESULTS
## 1 Introduction: the Hubbard model
Let us consider the Hubbard model on a square lattice of $`N\times N`$ sites, with hamiltonian
$$H=H_0+W,$$
(1)
where
$$H_0=t\underset{\sigma }{}\underset{(r,r^{})}{}c_{r\sigma }^{}c_{r^{}\sigma }$$
(2)
with $`r,r^{}`$ nearest neighbours and
$$W=U\underset{r}{}n_rn_r.$$
(3)
We use periodic boundary conditions and the one-body wave vectors are $`k=(k_x,k_y)=\frac{2\pi }{N}(p,q)`$ with $`p`$ and $`q`$ integers.
## 2 W=0 pairs
The strong on-site repulsion between two opposite spin fermions normally prevents the formation of singlet bound states. However, the planar square symmetry of the Hubbard and related models allows singlet two body eigenstates of $`H_0`$ belonging to the zero eigenvalue of $`W`$. We will refer to these states as W=0 pairs. This means that the two fermions of a W=0 pair do not interact directly, but only by means of virtual electron-hole excitations. The background particles play a crucial role in determining the structure of the effective interaction which can be, in priciple, attractive or repulsive.
In previous works<sup>1,2</sup> we have shown how to build W=0 pairs with zero total momentum. There the idea was to project the determinantal state
$$|d(k)=c_k^{}c_k^{}|vac$$
(4)
on the Irreducible Representations (irreps) of the point symmetry group $`C_{4v}`$. Remarkably, by projecting on the irreps $`A_2`$, $`B_1`$ and $`B_2`$ one obtains exclusively W=0 pairs. Here we want to point out a new and more general criterion to get W=0 pairs. Let $`𝒢`$ a symmetry group of the non interacting Hubbard Hamiltonian $`H_0`$, big enough to justify the degeneracy of the single particle energy levels. Let us consider a two body state of opposite spins trasforming as the i-th component of the irrep $`\mathrm{\Gamma }`$ of $`𝒢`$:
$$|\mathrm{\Psi }_i^{(\mathrm{\Gamma })}=\underset{r_1r_2}{}\mathrm{\Psi }_i^{(\mathrm{\Gamma })}(r_1,r_2)c_{r_1}^{}c_{r_2}^{}|vac$$
(5)
Then we have
$$n_rn_r|\mathrm{\Psi }_i^{(\mathrm{\Gamma })}=\mathrm{\Psi }_i^{(\mathrm{\Gamma })}(r,r)c_r^{}c_r^{}|vac\mathrm{\Psi }_i^{(\mathrm{\Gamma })}(r,r)|r,r.$$
(6)
Let $`P_i^{(\mathrm{\Gamma })}`$ be the projection operator on the i-th component of the irrep $`\mathrm{\Gamma }`$. Since
$$P_i^{(\mathrm{\Gamma })}\underset{r}{}\mathrm{\Psi }_i^{(\mathrm{\Gamma })}(r,r)|r,r=\underset{r}{}\mathrm{\Psi }_i^{(\mathrm{\Gamma })}(r,r)|r,r$$
(7)
if
$$P_i^{(\mathrm{\Gamma })}|r,r=0r$$
(8)
then
$$\mathrm{\Psi }_i^{(\mathrm{\Gamma })}(r,r)=0r.$$
(9)
Clearly eq.(8) is true if and only if
$$P_i^{(\mathrm{\Gamma })}|r\sigma =0r$$
where $`|r\sigma =c_{r\sigma }^{}|vac`$.
It is always possible to write $`|r\sigma `$ as
$$|r\sigma =\underset{\mathrm{\Gamma }}{}\underset{i}{}c_i^{(\mathrm{\Gamma })}(r)|\phi _{i,\sigma }^{(\mathrm{\Gamma })}$$
(10)
where $``$ is the set of the irreps of the one-body spectrum of $`H_0`$ and $`|\phi _{i,\sigma }^{(\mathrm{\Gamma })}`$ the corresponding eigenstate with spin $`\sigma `$. From (10) it follows directly that if $`\mathrm{\Gamma }^{}`$ does not belong to $``$
$$P^{(\mathrm{\Gamma }^{})}|r\sigma =0$$
and so
$$P^{(\mathrm{\Gamma }^{})}|r,r=0.$$
We have proven the following
THEOREM: Let $`|\mathrm{\Psi }`$ be a two-body eigenstate of the kinetic energy $`H_0`$ with spin $`S_z=0`$. Projecting $`|\mathrm{\Psi }`$ on an irrep not contained in $``$, we get either zero or an eigenstate of $`H_0`$ with no double occupancy.
The singlet component of this state is a W=0 pair, and its momentum does not generally vanish.
## 3 Binding energy: analytic approach
From now on we call hole the fermion created by $`c^{}`$ and electron its antifermion.
The Schrödinger equation for the ground state of our $`N\times N`$ square lattice with $`n_h`$ holes is
$$H|\mathrm{\Psi }_0=E_h(n_h)|\mathrm{\Psi }_0.$$
(11)
If the non interacting ground state with $`n_h2`$ holes can be written in terms of a single determinantal state, the exact $`|\mathrm{\Psi }_0`$ can always be expanded in terms of excitations over it:
$$|\mathrm{\Psi }_0=\underset{m}{}a_m|m+\underset{\alpha }{}b_\alpha |\alpha +\underset{\beta }{}c_\beta |\beta +\mathrm{}.$$
(12)
here $`m`$ runs over pair states, $`\alpha `$ over 4-body states ($`2`$ holes and $`1`$ electron-hole pair), $`\beta `$ over 6-body ones ($`2`$ holes and $`2`$ electron-hole pairs) and so on. Eq.(12) is an expansion in the number of virtual excitations.
In the following we set up a procedure which 1) carries out the “Configuration Interaction” calculation in a compact and efficient way 2) separates neatly the effective interaction from the self-energy contributions to the ground state energy. To understand how this mechanism works, for the moment we truncate the expansion to the $`\beta `$ states. Then equation (11) yields
$$\left(E_mE_h(n_h)\right)a_m+\underset{m^{}}{}a_m^{}W_{m,m^{}}+\underset{\alpha }{}b_\alpha W_{m,\alpha }+\underset{\beta }{}c_\beta W_{m,\beta }=0$$
(13)
$$\left(E_\alpha E_h(n_h)\right)b_\alpha +\underset{m^{}}{}a_m^{}W_{\alpha ,m^{}}+\underset{\alpha ^{}}{}b_\alpha ^{}W_{\alpha ,\alpha ^{}}+\underset{\beta }{}c_\beta W_{\alpha ,\beta }=0$$
(14)
$$\left(E_\beta E_h(n_h)\right)c_\beta +\underset{m^{}}{}a_m^{}W_{\beta ,m^{}}+\underset{\alpha ^{}}{}b_\alpha ^{}W_{\beta ,\alpha ^{}}+\underset{\beta ^{}}{}c_\beta ^{}W_{\beta ,\beta ^{}}=0.$$
(15)
Choosing the $`\beta `$ states in such a way that
$$(H_0+W)_{\beta \beta ^{}}=E_\beta ^{}\delta _{\beta \beta ^{}}$$
we exactly decouple the 6-body states getting
$$\left(E_mE_h(n_h)\right)a_m+\underset{m^{}}{}a_m^{}W_{m,m^{}}^{}+\underset{\alpha }{}b_\alpha W_{m,\alpha }^{}=0$$
$$\left(E_\alpha E_h(n_h)\right)b_\alpha +\underset{m^{}}{}a_m^{}W_{\alpha ,m^{}}^{}+\underset{\alpha ^{}}{}b_\alpha ^{}W_{\alpha ,\alpha ^{}}^{}=0$$
where $`W^{}`$’s are the renormalized interaction coefficients. It is clear that if we had truncated the expansion to an arbitrary number $`n`$ of electron-hole virtual exitations, we should have obtained the same results but with further renormalizations. This is a recursion method to perform the full canonical transformation; it applies to all the higher order interactions, and we can recast our problem as if only $`2`$ and $`4`$-body states existed. Now we choose the $`\alpha `$ states in such a way that
$$(H_0+W^{})_{\alpha ,\alpha ^{}}=E_\alpha ^{}\delta _{\alpha ,\alpha ^{}}$$
getting the following eigenvalue problem
$$\left(E_h(n_h)E_m\right)a_m=\underset{m^{}}{}a_m^{}m|F[E_h(n_h)]+W_{eff}[E_h(n_h)]|m^{},$$
(16)
where
$$m|F\left[E_h(n_h)\right]+W_{eff}[E_h(n_h)]|m^{}=W_{m,m^{}}^{}+\underset{\alpha }{}\frac{m|W^{}|\alpha \alpha |W^{}|m^{}}{E_h(n_h)E_\alpha ^{}}.$$
(17)
Equation (16) determines the amplitudes $`a_m`$ of the $`m`$ states in the $`n_h`$ hole ground state and the corresponding eigenvalue $`E_h(n_h)`$ relative to the hole vacuum. Here $`F`$ is the forward scattering operator and $`W_{eff}`$ the effective interaction. After the shift
$$H_0H_0E_h(n_h2),$$
(18)
it is perfectly consistent to interpret $`a_m`$ as the wave function of the dressed pair. It is worth to underline that this canonical transformation enables us to identify the effective interaction between two holes in a $`n_h2`$ holes background self-consistently. In principle we can work it out analytically and the expansion is neither in $`U`$ nor in $`t`$ but in the number of vitual exitations. It permits us to identify the binding energy $`|\mathrm{\Delta }(n_h)|`$ too. After the shift (18) we write the ground state energy $`E_h(n_h)`$ of the Hubbard hamiltonian as
$$E_h(n_h)=2E_F+\mathrm{\Delta }(n_h)$$
(19)
where $`E_F`$ is the renormalized Fermi energy and $`\mathrm{\Delta }(n_h)`$ the two holes energy gain of the interacting system with respect to the non interacting one. If the effective interaction $`W_{eff}`$ is attractive and $`\mathrm{\Delta }(n_h)<0`$ we can speak of hole pairing.
## 4 Pairing in the 4$`\times `$4 lattice Hubbard model
In this section we want to compare the results for the binding energy $`|\mathrm{\Delta }(n_h)|`$ obtained 1) from eq.(19) by truncating the canonical transformation to the $`\alpha `$ states and 2) from the definition
$$\mathrm{\Delta }(n_h)=E_h(n_h)+E_h(n_h2)2E_h(n_h1)$$
(20)
by computing exactly the three ground state energies involved in eq.(20). The equivalence of the two definitions of the binding energy was shown in Refs,.
We consider a 4$`\times `$4 lattice Hubbard model with $`t=1`$ and periodic boundary conditions; we investigate the possibility of pairing when two holes are added to a two hole background. The one particle energy spectrum of $`H_0`$ has 5 equally spaced levels having degeneracy 1,4,6,4,1 respectively.
Our first task is to determine of the symmetry group $`𝒢`$. The Space Group containing the translations and the 8 $`C_{4v}`$ operations is not enough to explain the degeneracy 6. Indeed the largest dimension of the irreps of the Space Group for this 4$`\times `$4 lattice is 4. The Space Group has 128 elements and 20 classes. As observed by previous authors<sup>3,4</sup> there must be additional space symmetries. We have found how to obtain them. Label the sites from 1 to 16 in the natural way; then rotate the plaquettes 1,2,5,6 and 11,12,15,16 clockwise and the other two counterclockwise by 90 degrees. This is one of the missing space symmetries, as one can see that each site preserves its first neighbours. It is not an isometry. The other missing symmetries are generated by adding this one to the Space Group. In this way we obtain the Symmetry Group $`𝒢`$ of 384 elements, which is large enough to justify the degeneracy of each one particle energy level. It still has 20 classes; the dimensions of the irreps are 1,1,1,1,2,2,3,3,3,3,4,4,4,4,6,6,6,6,8,8.
The exact interacting ground state of the system with 4 holes is threefold degenerate and belongs to an irrep of $`𝒢`$ that contains the irreps $`A_1`$ and $`B_1`$ of $`C_{4v}`$ once and two times respectively. This irrep is not contained in the one particle spectrum of $`H_0`$, which does not admit degeneracy 3. By the above theorem, pairs belonging to it must be W=0 pairs. These W=0 pairs arise from the single-particle states of $`H_0`$ with eigenvalue -2, and their irrep is contained in the square of the irrep of the single-particle level.
The $`m`$ states of the expansion (12) are all the W=0 pairs belonging to the ground state irrep. We have computed the binding energy $`|\mathrm{\Delta }(4)|`$ analitically from eq.(19) by truncating the expansion (12) to the $`\alpha `$ states and also by means of exact diagonalization from eq.(20). The results are listed below for different values of the on-site interaction $`U`$ (energies are in eV):
| | $`U=`$0.1 | $`U=`$0.5 | $`U=`$0.7 | $`U=`$1 |
| --- | --- | --- | --- | --- |
| $`\mathrm{\Delta }(4)_{exact}`$ | -0.053 | -0.32 | -0.67 | -1.18 |
| $`\mathrm{\Delta }(4)_{analytic}`$ | -0.078 | -1.95 | -3.83 | -7.81 |
As expected, with increasing $`U`$ the difference $`|\mathrm{\Delta }(4)_{exact}\mathrm{\Delta }(4)_{analytic}|`$ increases because the renormalization induced by virtual electron-hole exitations becomes important and it is no longer a good approximation to consider the $`\alpha `$ states only. Nevertheless the canonical transformation still predicts the right sign of $`\mathrm{\Delta }`$.
The above canonical transformation applies when two holes are added to a determinantal vacuum. To study the system at and close to half filling, we have extended the above canonical transformation to the case when the vacuum is degenerate. We are able to demonstrate analytically that holes pair and to explain the ground state symmetries that were found numerically in Ref.,. Thus, W=0 pairs are responsible for pairing and for the symmetry of the interacting ground state. This is also confirmed by the superconducting quantization of magnetic flux; in previous works<sup>2,6</sup> we show that the $`C_{4v}`$ symmetry is restored exactly at half fluxon $`\varphi _0/2=hc/2e`$ allowing the existence of W=0 pairs. The symmetry of the interacting ground state is still the same of the possible W=0 pairs and the corresponding energy versus flux has there a second minimum. A full account of the theory will be submitted elsewhere<sup>7</sup>. It is clear by now, however, that the existence of bound pairs of nonvanihing momentum opens up the possibility of a Jahn-Teller distortion and the present approach is likely to predict charge inomogeneities.
## Acknowledgements
This work has been supported by the Istituto Nazionale di Fisica della Materia.
## References
1. Michele Cini, Gianluca Stefanucci and Adalberto Balzarotti, Solid State Communication 109, 229 (1999).
2. Michele Cini, Gianluca Stefanucci and Adalberto Balzarotti, European Physical Journal B 10, 293 (1999).
3. G. Fano, F. Ortolani and A. Parola, Phys. Rev. B 46, 1048 (1992).
4. J. Bonca, P. Prelovsek and I. Sega, Phys. Rev. B 39, 7074 (1989).
5. A. Parola, S. Sorella, M. Parrinello and E. Tosatti, Phys. Rev. B 43, 6190 (1991).
6. Michele Cini and Adalberto Balzarotti, Phys. Rev. B 56, 14711 (1997); Michele Cini Adalberto Balzarotti and Gianluca Stefanucci, European Physical Journal B 14, 269 (2000).
7. Michele Cini and Gianluca Stefanucci, to be published.
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# 1 Introduction
## 1 Introduction
After the sphaleron solution in the Weinberg–Salam model had been found, the temperature dependence of baryon number violating processes (BNVP) was studied extensively. To understand the overall features of BNVP over the entire range of temperature, the computation of periodic instantons and their corresponding classical actions is required. However, the calculation of these in the Weinberg–Salam model is a highly non–trivial problem, even if numerical techniques are employed. Hence in many cases simple toy models were used to explore the temperature dependence of BNVP.
An immediate candidate as a simple toy model is the $`d=2`$ Mottola–Wipf(MW) model, which shares many common features with $`d=4`$ electroweak theory. The scale invariance of the nonlinear $`O(3)`$ model is broken in the MW model by adding an explicit mass term. This has a close analogy to the fact that the conformal invariance of the electroweak theory is broken in the Higgs sector. Also, neither model supports a vacuum instanton which gives a dominant contribution to the winding number transition at low temperature. The transition between thermally assisted quantum tunneling dominated by periodic instantons and the classical crossover dominated by the sphaleron in the MW model has been analyzed in Refs. using the method of , and it has been shown that the instanton–sphaleron transition is of the sharp first–order type in the full range of parameter space.
Recently, however, a numerical study of the $`d=4`$ SU(2)–Higgs model – which is a bosonic sector of the electroweak theory – has shown that a smooth second–order transition occurs when $`6.665<M_H/M_W<12.03`$ although the first–order transition occurs when $`M_H/M_W<6.665`$. This implies that the MW model does not exhibit a proper transition of BNVP when heavy Higgs’s are involved.
Another candidate as a toy model is the $`d=2`$ Abelian-Higgs model which supports vortex solutions, in particular the vacuum instanton and the sphaleron simultaneously. The simultaneous existence of instanton and sphaleron causes the model to yield phase diagrams for the instanton–sphaleron transition which are completely different from those of electroweak theory, as shown in Ref.. Furthermore, numerical and analytical approaches have shown that the instanton–sphaleron transition in this model is always of the second–order type, regardless of the ratio $`M_H/M_W`$. Hence, contrary to the MW model, the ordinary Abelian–Higgs model does not describe the instanton–sphaleron transition of the electroweak theory properly when the Higgs mass is small.
In the following we study the instanton–sphaleron transition in the $`d=2`$ Abelian-Higgs model when the spatial coordinate is compactified to $`S^1`$. Quite apart from the question of the physical relevance of the investigation below, it is a natural theoretical curiosity to inquire what the order of thermal transitions would be in this case, and we present the answer here. Physically, of course, the transitions we investigate are not those with respect to an order parameter as in the Weinberg–Salam theory, but with respect to temperature or inverse period of the periodic instantons in the potential barrier. These transitions have physically the meaning of transitions between classical and quantum behavior . Nonetheless, as stated we consider the model as an analogy, which enables us to investigate corresponding behavior. Since, to our knowledge, the effect of the compactification of the spatial coordinate of this model has not yet been investigated, this is also of interest on its own. Furthermore, we show that this model exhibits both first–order and second–order transitions depending on the size of the circumference of the spatial coordinate domain, i.e. the first order transition disappears in the limit of the circumference becoming infinitely large, in fact, even beyond a finite critical value. This means that the Abelian–Higgs model defined on a circle can be a better toy model than the MW model or the uncompactified Abelian-Higgs model for an analysis which can be compared with that of BNVP. One may wonder how, if at all, this situation compares with finite size scaling effects, i.e. of lattices with periodic boundary conditions, in lattice gauge theory contexts. In the latter, see e.g., lattice sizes of $`4^4`$ to $`16^{16}`$ are used and the possible dependence of the order of thermal transitions on these is investigated. Nonetheless the lattice sizes used are presumably still too small to permit definite conclusions about the scaling regime. We do not think that our case is really comparable to that. Rather we view the present model as a testing ground for various aspects related to phase transitions, since the study of the latter, as can be seen from the computations needed in the following, are highly nontrivial (and we therefore have to present some technical details), so that any model that can be handled to a large extent analytically, is worth studying. Thus before higher dimensional cases can be attacked convincingly, it is essential to have a thorough understanding of a lower dimensional one like the one we study here. This is therefore another main objective of the following.
## 2 The Sphaleron Configuration
We begin with the Euclidean action
$$S_E^{(0)}=𝑑\tau 𝑑x\left[\frac{1}{4}F_{\mu \nu }F_{\mu \nu }+(D_\mu \varphi )^{}D_\mu \varphi +\lambda [\varphi ^2\frac{v^2}{2}]^2\right]$$
(1)
and its field equations
$`_\mu F_{\mu \nu }`$ $`=`$ $`ig\left[\varphi ^{}(D_\nu \varphi )(D_\nu \varphi )^{}\varphi \right],`$ (2)
$`D_\mu D_\mu \varphi `$ $`=`$ $`2\lambda \varphi (\varphi ^2{\displaystyle \frac{v^2}{2}}),`$
where $`D_\mu =_\mu igA_\mu `$. We define as mass–dimensional parameters
$`M_H`$ $``$ $`\sqrt{2\lambda }v,`$ (3)
$`M_W`$ $``$ $`gv,`$
which correspond to Higgs mass and gauge particle mass in electroweak theory respectively. It is easy to show that the static sphaleron solution in the $`A_0=0`$ gauge is given by
$`A_1`$ $`=`$ $`A=const,`$ (4)
$`\varphi _{sph}`$ $`=`$ $`{\displaystyle \frac{kb(k)}{\sqrt{\lambda }}}e^{igAx}sn[b(k)x],`$
where $`sn[z]`$ is a Jacobian elliptic function, $`k`$ is the modulus of the elliptic function, and
$$b(k)=\sqrt{\frac{\lambda }{2}}v\left(\frac{2}{1+k^2}\right)^{\frac{1}{2}}.$$
(5)
Since $`sn[z]`$ has period $`4K(k)`$, where $`K(k)`$ is the complete elliptic integral of the first kind, the circumference $`L`$ of $`S^1`$ is defined by
$$L_n=\frac{4nK(k)}{b(k)},n=1,2,3\mathrm{}.$$
(6)
Since the transition rate is negligible for large $`n`$, we restrict ourselves to the $`L=L_1`$ case here. In view of $`K(1)=\mathrm{}`$ we see that $`k=1`$ gives the uncompactified limit we investigated previously. Thus, since this case does not lead to a first order transition, we can expect one, if at all, only in the domain of small values of the elliptic modulus $`k`$, and in fact, we shall see that this is the case.
In order to examine the type of instanton–sphaleron transition we have to introduce the fluctuation fields around the sphaleron and expand the field equations (2) up to the third order in these fields. If, however, one expands Eq.(2) naively, one will realize that the fluctuation operators are not diagonalized and, hence, the computation of the spectra of these operators becomes a very non–trivial problem. To avoid this difficulty, we choose the $`R_\xi `$ gauge by adding to the original action (1) the gauge fixing term
$$S_{gf}=\frac{1}{2\xi }𝑑\tau 𝑑x\left[_\mu A_\mu +\frac{ig}{2}\xi (\varphi ^2\varphi ^2)\right]^2.$$
(7)
Then the field equations for the total Euclidean action $`S_E=S_E^{(0)}+S_{gf}`$ become
$`_\mu F_{\mu \nu }+{\displaystyle \frac{1}{\xi }}\left[_\mu _\nu A_\mu +ig\xi (\varphi _\nu \varphi \varphi ^{}_\nu \varphi ^{})\right]`$ $`=`$ $`ig\left[\varphi ^{}(D_\nu \varphi )(D_\nu \varphi )^{}\varphi \right],`$ (8)
$`D_\mu D_\mu \varphi +ig\varphi ^{}\left[_\mu A_\mu +{\displaystyle \frac{ig\xi }{2}}(\varphi ^2\varphi ^2)\right]`$ $`=`$ $`2\lambda \varphi (\varphi ^2{\displaystyle \frac{v^2}{2}}).`$
One can show that the sphaleron in this gauge is the same as that of Eq.(4) if $`A=0`$:
$`A_1`$ $`=`$ $`0,`$ (9)
$`\varphi _{sph}`$ $`=`$ $`{\displaystyle \frac{kb(k)}{\sqrt{\lambda }}}sn[b(k)x].`$
We have therefore determined the sphaleron configuration in the most optimal way to permit continuation with the following difficult computations.
## 3 Fluctuations about the Sphaleron
We now introduce the fluctuation fields around the sphaleron configuration by setting
$`A_0(\tau ,x)`$ $`=`$ $`a_0(\tau ,x),`$ (10)
$`A_1(\tau ,x)`$ $`=`$ $`a_1(\tau ,x),`$
$`\varphi (\tau ,x)`$ $`=`$ $`\varphi _{sph}(x)+{\displaystyle \frac{1}{\sqrt{2}}}\left(\eta _1(\tau ,x)+i\eta _2(\tau ,x)\right),`$
where $`a_0`$, $`a_1`$, $`\eta _1`$, and $`\eta _2`$ are real fields. Inserting (10) into Eq.(1) and Eq.(7) one can express $`S_E`$ for $`\xi =1`$ as
$$S_E=\frac{E_{sph}}{T}+S_2+S_3+S_4$$
(11)
where $`1/T`$ is the period of the sphaleron and
$`E_{sph}`$ $`=`$ $`\sqrt{2\lambda }v^3[[\left({\displaystyle \frac{2}{1+k^2}}\right)^{\frac{1}{2}}+{\displaystyle \frac{1+2k^2}{3}}\left({\displaystyle \frac{2}{1+k^2}}\right)^{\frac{3}{2}}2\left({\displaystyle \frac{2}{1+k^2}}\right)^{\frac{1}{2}}]K(k)`$
$`+[2\left({\displaystyle \frac{2}{1+k^2}}\right)^{\frac{1}{2}}{\displaystyle \frac{1+k^2}{3}}\left({\displaystyle \frac{2}{1+k^2}}\right)^{\frac{3}{2}}]E(k)],`$
$`S_2`$ $`=`$ $`{\displaystyle }d\tau dx[{\displaystyle \frac{1}{2}}a_0[_\mu _\mu +2g^2\varphi _{sph}^2]a_0+{\displaystyle \frac{1}{2}}a_1[_\mu _\mu +2g^2\varphi _{sph}^2]a_1`$
$`+{\displaystyle \frac{1}{2}}\eta _1\left[_\mu _\mu +2\lambda (3\varphi _{sph}^2{\displaystyle \frac{v^2}{2}})\right]\eta _1+{\displaystyle \frac{1}{2}}\eta _2\left[_\mu _\mu +2(\lambda +g^2)\varphi _{sph}^2\lambda v^2\right]\eta _2`$
$`+2\sqrt{2}g\varphi _{sph}^{}a_1\eta _2],`$
$`S_3`$ $`=`$ $`{\displaystyle }d\tau dx[2g(a_0\dot{\eta _1}\eta _2+a_1\eta _1^{}\eta _2)+\sqrt{2}g^2\varphi _{sph}(a_0^2+a_1^2)\eta _1`$
$`+\sqrt{2}\lambda \varphi _{sph}\eta _1^3+\sqrt{2}(\lambda +g^2)\varphi _{sph}\eta _1\eta _2^2],`$
$`S_4`$ $`=`$ $`{\displaystyle 𝑑\tau 𝑑x\left[\frac{g^2}{2}(a_0^2+a_1^2)(\eta _1^2+\eta _2^2)+\frac{\lambda }{4}(\eta _1^2+\eta _2^2)^2+\frac{g^2}{2}\eta _1^2\eta _2^2\right]}.`$
In these equations $`E(k)`$ is the complete elliptic integral of the second kind, and the dot and the prime denote differentiation with respect to $`\tau `$ and $`x`$ respectively. Owing to the final term in $`S_2`$ the fluctuation operators for $`a_1`$ and $`\eta _2`$ are not diagonalized although the $`R_{\xi =1}`$ gauge has been chosen. To obtain the diagonalization we introduce the fluctuation fields $`\rho _\pm `$ defined as
$$\rho _+=v_1a_1+v_2\eta _2,\rho _{}=v_2a_1+v_1\eta _2,$$
(13)
where
$$v_1=\sqrt{\frac{1(\varphi _{sph}^2\frac{v^2}{2})f_1^{\frac{1}{2}}}{2}},v_2=\sqrt{\frac{1+(\varphi _{sph}^2\frac{v^2}{2})f_1^{\frac{1}{2}}}{2}},$$
(14)
and
$$f_1=(\varphi _{sph}^2\frac{v^2}{2})^2\mathrm{cosh}^2\alpha \frac{v^4}{4}\left(\frac{1k^2}{1+k^2}\right)\mathrm{sinh}^2\alpha .$$
(15)
Here $`\alpha =\mathrm{sinh}^12\theta `$ and $`\theta `$ is the dimensionless parameter
$$\theta \frac{2M_W}{M_H}=\sqrt{\frac{2g^2}{\lambda }}.$$
(16)
Using the field redefinition (13) and the first–order differential equation for $`\varphi _{sph}`$,
$$\varphi _{sph}^{}+\sqrt{\lambda }\left[\frac{v^4}{4}\left(\frac{2k}{1+k^2}\right)^2v^2\varphi _{sph}^2+\varphi _{sph}^4\right]^{\frac{1}{2}}=0,$$
(17)
it is straightforward to show that $`S_2`$ becomes
$$S_2=\frac{1}{2}𝑑\tau 𝑑x[a_0D_0a_0+\eta _1D_1\eta _1+\rho _+D_+\rho _++\rho _{}D_{}\rho _{}],$$
(18)
where
$`D_0`$ $`=`$ $`_\mu _\mu +2g^2\varphi _{sph}^2,`$ (19)
$`D_1`$ $`=`$ $`_\mu _\mu +2\lambda (3\varphi _{sph}^2{\displaystyle \frac{v^2}{2}}),`$
$`D_\pm `$ $`=`$ $`_\mu _\mu +2g^2\varphi _{sph}^2+\lambda (\varphi _{sph}^2{\displaystyle \frac{v^2}{2}})\lambda \sqrt{f_1}.`$
After inserting the field redefinition (13) into $`S_3`$ and $`S_4`$, one can derive the field equations for the fluctuation fields by varying the total action $`S_E`$, i.e. (the method of Ref. to determine the order of thermal transitions requires all the terms written out explicitly here)
$$\widehat{l}\left(\begin{array}{c}a_0\\ \rho _+\\ \rho _{}\\ \eta _1\end{array}\right)=\widehat{h}\left(\begin{array}{c}a_0\\ \rho _+\\ \rho _{}\\ \eta _1\end{array}\right)+\left(\begin{array}{c}G_2^{a_0}\\ G_2^{\rho _+}\\ G_2^\rho _{}\\ G_2^{\eta _1}\end{array}\right)+\left(\begin{array}{c}G_3^{a_0}\\ G_3^{\rho _+}\\ G_3^\rho _{}\\ G_3^{\eta _1}\end{array}\right)+\mathrm{}$$
(20)
where
$$\begin{array}{cc}\widehat{l}=\left(\begin{array}{cccc}\frac{^2}{z_0^2}& 0\hfill & 0& \hfill 0\\ 0& \frac{^2}{z_0^2}\hfill & 0& \hfill 0\\ 0& 0\hfill & \frac{^2}{z_0^2}& \hfill 0\\ 0& 0\hfill & 0& \hfill \frac{^2}{z_0^2}\end{array}\right),\hfill & \widehat{h}=\left(\begin{array}{cccc}\widehat{h}_{a_0}& 0\hfill & 0& \hfill 0\\ 0& \widehat{h}_{\rho _+}\hfill & 0& \hfill 0\\ 0& 0\hfill & \widehat{h}_\rho _{}& \hfill 0\\ 0& 0\hfill & 0& \hfill \widehat{h}_{\eta _1}\end{array}\right),\end{array}$$
and
$`G_2^{a_0}`$ $`=`$ $`{\displaystyle \frac{2g}{b(k)}}(v_2\rho _++v_1\rho _{})\dot{\eta _1}+{\displaystyle \frac{2\sqrt{2}g^2}{b^2(k)}}\varphi _{sph}a_0\eta _1,`$ (21)
$`G_3^{a_0}`$ $`=`$ $`{\displaystyle \frac{g^2}{b^2(k)}}a_0\left[\eta _1^2+(v_2\rho _++v_1\rho _{})^2\right],`$
$`G_2^{\rho _+}`$ $`=`$ $`{\displaystyle \frac{2g}{b(k)}}[v_2a_0\dot{\eta _1}+(v_1^2v_2^2)\rho _{}\eta _1^{}+2v_1v_2\rho _+\eta _1^{}],`$
$`+{\displaystyle \frac{2\sqrt{2}\lambda }{b^2(k)}}\varphi _{sph}[v_2^2\rho _+\eta _1+v_1v_2\rho _{}\eta _1]+{\displaystyle \frac{2\sqrt{2}g^2}{b^2(k)}}\varphi _{sph}\rho _+\eta _1,`$
$`G_3^{\rho _+}`$ $`=`$ $`{\displaystyle \frac{g^2}{b^2(k)}}[\rho _+\eta _1^2+v_2^2a_0^2\rho _++v_1v_2a_0^2\rho _{}+2v_1^2v_2^2\rho _+^3+3v_1v_2(v_1^2v_2^2)\rho _+^2\rho _{}`$
$`+(v_1^44v_1^2v_2^2+v_2^4)\rho _+\rho _{}^2v_1v_2(v_1^2v_2^2)\rho _{}^3]`$
$`+`$ $`{\displaystyle \frac{\lambda }{b^2(k)}}\left[v_2^2\rho _+\eta _1^2+v_1v_2\rho _{}\eta _1^2+v_2^4\rho _+^3+3v_1v_2^3\rho _+^2\rho _{}+3v_1^2v_2^2\rho _+\rho _{}^2+v_1^3v_2\rho _{}^3\right],`$
$`G_2^\rho _{}`$ $`=`$ $`{\displaystyle \frac{2g}{b(k)}}[v_1a_0\dot{\eta _1}+(v_1^2v_2^2)\rho _+\eta _1^{}2v_1v_2\rho _{}\eta _1^{}]`$
$`+{\displaystyle \frac{2\sqrt{2}\lambda }{b^2(k)}}\varphi _{sph}[v_1^2\rho _{}\eta _1+v_1v_2\rho _+\eta _1]+{\displaystyle \frac{2\sqrt{2}g^2}{b^2(k)}}\varphi _{sph}\rho _{}\eta _1,`$
$`G_3^\rho _{}`$ $`=`$ $`{\displaystyle \frac{g^2}{b^2(k)}}[\rho _{}\eta _1^2+v_1^2a_0^2\rho _{}+v_1v_2a_0^2\rho _++2v_1^2v_2^2\rho _{}^3+v_1v_2(v_1^2v_2^2)\rho _+^3`$
$`+(v_1^44v_1^2v_2^2+v_2^4)\rho _+^2\rho _{}3v_1v_2(v_1^2v_2^2)\rho _+\rho _{}^2]`$
$`+{\displaystyle \frac{\lambda }{b^2(k)}}\left[v_1^2\rho _{}\eta _1^2+v_1v_2\rho _+\eta _1^2+v_1^4\rho _{}^3+v_1v_2^3\rho _+^3+3v_1^2v_2^2\rho _+^2\rho _{}+3v_1^3v_2\rho _+\rho _{}^2\right],`$
$`G_2^{\eta _1}`$ $`=`$ $`{\displaystyle \frac{2g}{b(k)}}[v_2(\dot{a_0}\rho _++a_0\dot{\rho _+})+v_1(\dot{a_0}\rho _{}+a_0\dot{\rho _{}})+2(v_1v_1^{}v_2v_2^{})\rho _+\rho _{}`$
$`+`$ $`(v_1^2v_2^2)(\rho _+^{}\rho _{}+\rho _+\rho _{}^{})+v_1^{}v_2(\rho _+^2\rho _{}^2)+v_1v_2^{}(\rho _+^2\rho _{}^2)+2v_1v_2(\rho _+\rho _+^{}\rho _{}\rho _{}^{}]`$
$`+{\displaystyle \frac{\sqrt{2}\lambda }{b^2(k)}}\varphi _{sph}(3\eta _1^2+v_2^2\rho _+^2+v_1^2\rho _{}^2+2v_1v_2\rho _+\rho _{})+{\displaystyle \frac{\sqrt{2}g^2}{b^2(k)}}\varphi _{sph}(a_0^2+\rho _+^2+\rho _{}^2),`$
$`G_3^{\eta _1}`$ $`=`$ $`{\displaystyle \frac{g^2}{b^2(k)}}(a_0^2+\rho _+^2+\rho _{}^2)\eta _1+{\displaystyle \frac{\lambda }{b^2(k)}}[\eta _1^3+(v_2\rho _++v_1\rho _{})^2\eta _1].`$
Here $`z_0b(k)\tau `$, $`z_1b(k)x`$, and the dot and the prime denote differentiation with respect to $`z_0`$ and $`z_1`$ respectively. Also, the fluctuation operators $`\widehat{h}_{a_0}`$, $`\widehat{h}_{\rho _+}`$, $`\widehat{h}_\rho _{}`$, and $`\widehat{h}_{\eta _1}`$ are
$`\widehat{h}_{a_0}`$ $`=`$ $`{\displaystyle \frac{^2}{z_1^2}}+{\displaystyle \frac{2g^2}{b^2(k)}}\varphi _{sph}^2,`$ (22)
$`\widehat{h}_{\rho _+}`$ $`=`$ $`{\displaystyle \frac{^2}{z_1^2}}+{\displaystyle \frac{1}{b^2(k)}}\left[2g^2\varphi _{sph}^2+\lambda (\varphi _{sph}^2{\displaystyle \frac{v^2}{2}})+\lambda \sqrt{f_1}\right],`$
$`\widehat{h}_\rho _{}`$ $`=`$ $`{\displaystyle \frac{^2}{z_1^2}}+{\displaystyle \frac{1}{b^2(k)}}\left[2g^2\varphi _{sph}^2+\lambda (\varphi _{sph}^2{\displaystyle \frac{v^2}{2}})\lambda \sqrt{f_1}\right],`$
$`\widehat{h}_{\eta _1}`$ $`=`$ $`{\displaystyle \frac{^2}{z_1^2}}+{\displaystyle \frac{2\lambda }{b^2(k)}}\left[3\varphi _{sph}^2{\displaystyle \frac{v^2}{2}}\right].`$
In the following section the resulting fluctuations and the characteristics of the thermal transitions are analyzed in detail.
## 4 Fluctuation Analysis and Quantum–Classical Transitions
Having derived the fluctuation equations, our next aim is to derive the eigenvalues and then with knowledge of the negative mode (as required in the method of Ref. ) to investigate quantum–classical transitions and their order.
The lowest few eigenvalues of $`\widehat{h}_{a_0}`$ and $`\widehat{h}_{\eta _1}`$ can be obtained exactly by using Lam$`\stackrel{´}{e}`$’s equation . It is easy to show that the spectrum of $`\widehat{h}_{a_0}`$ consists of only positive modes whose explicit forms are not needed here for further study. Also, of the lowest eigenstates of $`\widehat{h}_{\eta _1}`$, we need only the $`2K`$–antiperiodic eigenfunctions to recover the proper uncompactified limit as shown in Ref.. The lowest two $`2K`$–antiperiodic eigenstates of $`\widehat{h}_{\eta _1}`$ are summarized in Table I. It may be impossible to obtain the higher states analytically at present. Using $`_K^K\psi _i^{(\eta _1)}\psi _j^{(\eta _1)}𝑑z_1=\delta _{ij}`$, one can show (using formulae of Ref. ) that the normalization constant $`N_1`$ in Table I is given by
$$N_1=\sqrt{\frac{3k^2}{2[(1k^2)K(12k^2)E]}}.$$
(23)
We now consider the eigenstates of $`\widehat{h}_{\rho _+}`$ and $`\widehat{h}_\rho _{}`$. In Appendix A we explain how the eigenstates of $`\widehat{h}_{\rho _+}`$ and $`\widehat{h}_\rho _{}`$ are computed numerically. Following the method of Appendix A, one can show that the eigenstates of $`\widehat{h}_{\rho _+}`$ also consist of only positive modes which we do not need. What we need (as pointed out earlier), is only the negative mode of $`\widehat{h}_\rho _{}`$. If one performs the numerical calculation, one finds that $`\widehat{h}_\rho _{}`$ has two negative modes, one of which is $`2K`$–periodic and the other $`2K`$–antiperiodic. Fig. 1 shows the $`k`$–dependence of the negative eigenvalues for $`\theta =1`$. Since the $`2K`$–antiperiodic boundary condition is required for the proper continuum limit, we have to use the solid line in Fig. 1 as a negative eigenvalue. One should note that this negative eigenvalue approaches zero in the small $`k`$ region. We show in the following that this effect guarantees that the instanton–sphaleron transition in the small $`k`$–region is different from that in the large $`k`$–region. Fig. 2 shows normalized $`2K`$–antiperiodic eigenfunctions for the negative mode of $`\widehat{h}_\rho _{}`$ at ($`\theta =1`$, $`k=0.6`$) and ($`\theta =1`$, $`k=0.99`$). Their Gaussian shape is indicative of their ground–state nature (below the zero–eigenvalue of the translational mode).
We let $`\psi _1^{(\rho _{})}`$ and $`ϵ_1^{(\rho _{})}`$ be respectively the $`2K`$–antiperiodic eigenfunction and corresponding eigenvalue for the negative mode. To obtain the criterion for the sharp first–order instanton–sphaleron transition we have to compute the nonlinear correction to the frequency of the periodic instanton around the sphaleron. This can be carried out by solving Eq.(20) perturbatively. The perturbation procedure is briefly summarized in Appendix B. The criterion for the first-order transition is expressed as an inequality
$$\mathrm{\Omega }\mathrm{\Omega }_{sph}>0,$$
(24)
where $`\mathrm{\Omega }`$ is the frequency involving the nonlinear correction and $`\mathrm{\Omega }_{sph}\sqrt{ϵ_1^{(\rho _{})}}`$.
In Appendix B it is shown that the inequality (24) can be expressed as
$$<\psi _1^{(\rho _{})}D_1(z_1)><0$$
(25)
where
$$D_1(z_1)=D_1^{(1)}(z_1)+D_1^{(2)}(z_1)+D_1^{(3)}(z_1).$$
(26)
Here
$`D_1^{(1)}(z_1)`$ $`=`$ $`{\displaystyle \frac{2\sqrt{2(1+k^2)}}{v}}\psi _1^{(\rho _{})}(z_1)[k(v_1^2+{\displaystyle \frac{s_1(s_1+1)}{2}})sn[z_1]g_{\eta _1,1}(z_1)`$
$`\sqrt{s_1(s_1+1)}v_1v_2g_{\eta _1,1}^{}(z_1)],`$
$`D_1^{(2)}(z_1)`$ $`=`$ $`{\displaystyle \frac{\sqrt{2(1+k^2)}}{v}}\psi _1^{(\rho _{})}(z_1)[k(v_1^2+{\displaystyle \frac{s_1(s_1+1)}{2}})sn[z_1]g_{\eta _1,2}(z_1)`$
$`\sqrt{s_1(s_1+1)}v_1v_2g_{\eta _1,2}^{}(z_1)],`$
$`D_1^{(3)}(z_1)`$ $`=`$ $`{\displaystyle \frac{3(1+k^2)}{4v^2}}[v_1^4+s_1(s_1+1)v_1^2v_2^2]\psi _1^{(\rho _{})3}(z_1),`$ (27)
where $`s_1\sqrt{\theta ^2+\frac{1}{4}}\frac{1}{2}`$ and
$`g_{\eta _1,1}(z_1)`$ $`=`$ $`\widehat{h}_{\eta _1}^1q(z_1)>,`$ (28)
$`g_{\eta _1,2}(z_1)`$ $`=`$ $`(\widehat{h}_{\eta _1}+4\mathrm{\Omega }_{sph}^2)^1q(z_1)>,`$
$`q(z_1)>`$ $`=`$ $`{\displaystyle \frac{1}{v}}\sqrt{{\displaystyle \frac{1+k^2}{2}}}[\theta ((v_1v_2)^{}\psi _1^{(\rho _{})2}+2v_1v_2\psi _1^{(\rho _{})}\psi _1^{(\rho _{})})`$
$`+k(v_1^2+{\displaystyle \frac{\theta ^2}{2}})sn[z_1]\psi _1^{(\rho _{})2}].`$
It is now necessary to evaluate $`g_{\eta _1,1}`$ and $`g_{\eta _1,2}`$ explicitly. Although one can calculate $`g_{\eta _1,1}`$ exactly by following the procedure given in the Appendix of Ref., this is not necessary here. We already know the type of instanton–sphaleron transition at $`k=1`$ so that our interest concerns only the domain of small $`k`$. We can therefore adopt the following approximate procedure which has been shown to be valid in the small $`k`$ region. Using the completeness relation one can express $`g_{\eta _1,1}`$ as
$$g_{\eta _1,1}=\underset{n=0}{\overset{\mathrm{}}{}}\frac{<\psi _n^{(\eta _1)}q>}{ϵ_n^{(\eta _1)}}\psi _n^{(\eta _1)}>.$$
(29)
Since $`q>`$ is an odd function, the zero mode of $`\widehat{h}_{\eta _1}`$ does not contribute to the r.h.s. of Eq.(29). Hence the first approximation of $`g_{\eta _1,1}`$ is
$$g_{\eta _1,1}\frac{<\psi _1^{(\eta _1)}q>}{ϵ_1^{(\eta _1)}}\psi _1^{(\eta _1)}>$$
(30)
which can be evaluated numerically. In fact, this approximation is valid when $`\psi _1^{(\eta _1)}>`$ is an isolated discrete mode and the density of higher states is very dilute. Ref. shows these conditions are fulfilled in the small $`k`$–region if $`\widehat{h}_{\eta _1}`$ is a Lam$`\stackrel{´}{e}`$ operator as is the case here. In the same way $`g_{\eta _1,2}`$ is approximately
$$g_{\eta _1,2}\frac{<\psi _1^{(\eta _1)}q>}{3k^2+4\mathrm{\Omega }_{sph}^2}\psi _1^{(\eta _1)}>.$$
(31)
The plots of Fig. 3 show the $`k`$-dependence of
$$J_i<\psi _1^{(\rho _{})}D_1^{(i)}>$$
and of the sum $`J_1+J_2+J_3`$ for $`\theta =1`$. One can see that this sum becomes negative at approximately $`k=0.2`$ and therefore satisfies the inequality (25) and so (24) for the existence of a first order transition. Thus Fig. 3 demonstrates that the sharp first–order instanton–sphaleron transition occurs at $`k<k_c0.2`$ for $`\theta =1`$. Although the result is not included in this paper, we have checked also the $`\theta =3`$ case and have found a similar behavior: a sharp transition occurs in the small $`k`$ region.
## 5 Conclusions
The study of phase transitions is of considerable significance in many areas of physics, but – as is also evident from the above – this requires highly nontrivial efforts, both analytically and numerically. In the above we studied a model which permits a considerable fraction of analytical investigation, but finally requires also highly nontrivial computational work. The results we presented above answer the naturally asked question as to what behavior the Abelian–Higss model would reveal if the spacial coordinate is compactified on a circle. We have found that indeed a change occurs as compared to the uncompactified case, i.e. in the region of small elliptic modulus $`k`$ of the periodic instantons that we used, which corresponds to circle–circumferences below a critical size (a specific critical value was given for appropriate values of other parameters). Hence, depending on $`k`$, this model allows both smooth second–order transitions in the large $`k`$ region and sharp first–order transitions in the small $`k`$ region. These findings are similar to those of $`d=4`$ SU(2)–Higgs theory in which the type of transition depends on the ratio of $`M_H`$ and $`M_W`$. Thus our findings can be seen as an analogy. Of course, our model lacks direct physical application, but this was also not anticipated. Rather we explored the model also for the other reasons stated, i.e. as a matter of curiosity as to what type of thermal behavior will be found once the spatial coordinate is compactified, and as a further testing ground for methods of investigation of phase transitions, here in the sense of transitions from quantum to classical behavior. Such investigations are usually very complicated and there are few models that permit also transparent analytical investigation, at least in part.
Acknowledgement: This work has been supported in part by the Korean Research Foundation (Contract Number: KRF-2000-D00073).
Table
Appendix A
Here we explain how the spectrum of $`\widehat{h}_\rho _{}`$ is obtained. The spectrum of $`\widehat{h}_{\rho _+}`$ can be obtained similarly. The eigenvalue equation of $`\widehat{h}_\rho _{}`$ is
$$\left[\frac{^2}{z_1^2}+f(k,\theta ,z_1)\right]\psi _n^{(\rho _{})}=\zeta \psi _n^{(\rho _{})}$$
(32)
where
$`f(k,\theta ,z_1)`$ $`=`$ $`(1+\theta ^2)k^2sn^2[z_1]\sqrt{(1+4\theta ^2)\left(k^2sn^2[z_1]{\displaystyle \frac{1+k^2}{2}}\right)^2\theta ^2(1k^2)^2},`$
$`\zeta `$ $`=`$ $`ϵ^{(\rho _{})}+{\displaystyle \frac{1+k^2}{2}}.`$ (33)
We first choose the $`4K`$–periodic boundary condition. In this case we can use the Fourier expansions
$$f(k,\theta ,z_1)=\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}a_ne^{i\frac{n\pi }{l}z_1},\psi _n^{(\rho _{})}=\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}b_ne^{i\frac{n\pi }{l}z_1},$$
(34)
where $`l=2K`$ and the coefficients $`a_n`$ are given by
$$a_n=\frac{1}{2l}_l^lf(k,\theta ,z_1)e^{i\frac{n\pi }{l}z_1}.$$
(35)
Inserting (34) into (32) and using the property of linear independence of the exponential functions one obtains
$$\underset{m}{}\left[\left(\frac{n\pi }{l}\right)^2\delta _{mn}+a_{nm}\right]b_m=\zeta b_n.$$
(36)
Solving this matrix equation numerically, one can evaluate the eigenvalue $`ϵ_n^{(\rho _{})}`$ and eigenfunction $`\psi _n^{(\rho _{})}`$. After that we choose only $`2K`$–antiperiodic eigenfunctions and determine the corresponding eigenvalues for the proper $`k=1`$ limit.
Appendix B
In this appendix we show briefly how the inequality (25) is derived for the criterion of the sharp first–order transition by solving Eq.(20) perturbatively. First we choose an ansatz
$$\left(\begin{array}{c}a_0\\ \rho _+\\ \rho _{}\\ \eta _1\end{array}\right)=\mathrm{\Delta }\left(\begin{array}{c}a_{0,0}(z_1)\\ \rho _{+,0}(z_1)\\ \rho _{,0}(z_1)\\ \eta _{1,0}(z_1)\end{array}\right)\mathrm{cos}\mathrm{\Omega }_{sph}z_0$$
(37)
where $`\mathrm{\Delta }`$ is a small oscillation amplitude around the sphaleron. After inserting (37) into Eq.(20) and neglecting higher order terms, one obtains
$`\mathrm{\Omega }_{sph}`$ $`=`$ $`\sqrt{ϵ_1^{(\rho _{})}},`$ (38)
$`a_{0,0}`$ $`=`$ $`0,\rho _{+,0}=0,`$
$`\rho _{,0}`$ $`=`$ $`\psi _1^{(\rho _{})},\eta _{1,0}=0.`$
For the next order perturbation we set
$$\left(\begin{array}{c}a_0\\ \rho _+\\ \rho _{}\\ \eta _1\end{array}\right)=\left(\begin{array}{c}\mathrm{\Delta }^2a_{0,1}(z_0,z_1)\\ \mathrm{\Delta }^2\rho _{+,1}(z_0,z_1)\\ \mathrm{\Delta }\rho _{,0}(z_1)\mathrm{cos}\mathrm{\Omega }z_0+\mathrm{\Delta }^2\rho _{,1}(z_0,z_1)\\ \mathrm{\Delta }^2\eta _{1,1}(z_0,z_1)\end{array}\right).$$
(39)
Inserting Eq.(39) into Eq.(20) and considering only terms up to quadratic order, one can show there is no frequency shift to this order. It is also straightforward to show that $`a_{0,1}=0`$, $`\rho _{+,1}=0`$, $`\rho _{,1}=0`$, and
$$\eta _{1,1}=g_{\eta _1,1}(z_1)+g_{\eta _1,2}(z_1)\mathrm{cos}2\mathrm{\Omega }_{sph}z_0$$
(40)
where $`g_{\eta _1,1}`$ and $`g_{\eta _1,2}`$ are given by Eq.(28). For the next order perturbation we set
$$\left(\begin{array}{c}a_0\\ \rho _+\\ \rho _{}\\ \eta _1\end{array}\right)=\left(\begin{array}{c}\mathrm{\Delta }^3a_{0,2}(z_0,z_1)\\ \mathrm{\Delta }^3\rho _{+,2}(z_0,z_1)\\ \mathrm{\Delta }\rho _{,0}(z_1)\mathrm{cos}\mathrm{\Omega }z_0+\mathrm{\Delta }^3\rho _{,2}(z_0,z_1)\\ \mathrm{\Delta }^2\eta _{1,1}(z_0,z_1)+\mathrm{\Delta }^3\eta _{1,2}(z_0,z_1)\end{array}\right).$$
(41)
Inserting this into Eq.(20) and considering contributions up to cubic order, one can show that there is a frequency change in this order given by
$$\mathrm{\Omega }_{sph}^2\mathrm{\Omega }^2=\mathrm{\Delta }^2<\rho _{,0}D_1>$$
(42)
which proves Eq.(25).
Figure Captions
Fig.1
$`k`$–dependence of the negative eigenvalues $`ϵ^{(\rho _{})}`$ for $`\widehat{h}_\rho _{}`$ for $`\theta =1`$. The dotted line and the solid line represent the negative eigenvalues for the $`2K`$–periodic and $`2K`$–antiperiodic eigenfunctions respectively. For the correct $`k=1`$ limit we have to choose the solid line as the negative eigenvalue.
Fig.2
The normalized $`2K`$–antiperiodic eigenfunctions for the negative mode of $`\widehat{h}_\rho _{}`$ for (a) $`\theta =1`$, $`k=0.6`$, and (b) $`\theta =1`$, $`k=0.99`$.
Fig.3
$`k`$–dependence of $`J_1`$, $`J_2`$, $`J_3`$, and $`J_1+J_2+J_3`$ for $`\theta =1`$. This shows that the sharp first–order instanton–sphaleron transition occurs for $`k<k_c0.2`$.
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# About the Dependence of the Currency Exchange Rate at Time and National Dividend, Investments Size, Difference Between Total Demand and Supply
\[
## Abstract
The time dependence of the currency exchange rate $`K`$ treated as a function of national dividend, investments and difference between total demand for a goods and supply is considered. To do this a proposed earlier general algorithm of economic processes describing on the basis of the equations for $`K`$ like the equations of statistical physics of open systems is used. A number of differential equations (including nonlinear ones too) determining the time dependence of the exchange rate (including oscillations) is obtained.
\]
1. L.Ya.Kobelev et al. in order to describe economic phenomena offered a system of nonlinear differential equations based on the mathematical methods of statistical physics of open systems ().(see also -). The present paper is devoted to consideration of a model example of calculation of time dependence of national currency exchange rate treated as a function of national dividend $`F`$, investments C and total difference between demand and supply (with all the values being measured in national currency units). The modality of the example considered is in neglecting of dependencies of national currency upon the others, not mentioned economic variables, so the main purpose of this paper is not so to elucidate real economic laws and regularities (although this also takes place) as to illustrate the possibilities of the offered in the named paper method if applied to a concrete problem.
2. Thus, to show the advantages of the method used consider the dependence of the currency exchange rate upon investments $`C`$, difference between total demand and supply P, and national dividend $`F`$. Let K be a variation of the currency exchange rate per time unit. In this case the equations describing the exchange rate as a function of $`C`$, $`P`$ and $`F`$ takes the form of the simplified equations (3), (4), (7) of the paper
$$\frac{dK}{dt}=\frac{K}{t}+\frac{K}{C}\frac{dC}{dt}+\frac{K}{P}\frac{dP}{dt}+\frac{K}{F}\frac{dF}{dt}=I_C$$
(1)
$`I_C=\phi (K)+{\displaystyle \frac{}{C}}(D_C{\displaystyle \frac{}{C}}K)+`$ (2)
$`+{\displaystyle \frac{}{p}}(D_P{\displaystyle \frac{}{P}}K)+{\displaystyle \frac{}{F}}(D_F{\displaystyle \frac{}{F}}K)`$ (3)
The equations describing time dependencies of $`C`$, $`P`$, $`F`$ and the dependence of the currency exchange rate upon them write down as follows:
$$\frac{dK}{dx_i}=F_i(x_i,K)+\frac{}{x_j}[D_{j\beta }\frac{}{x_\beta }(x_iA_jx_i)]$$
(4)
$$\frac{dx_i}{dt}=\phi (x_j,\dot{x_j},K,\dot{K},\frac{K}{x_j})$$
(5)
where $`i=1,2,3`$ ($`x_1=C`$, $`x_2=P`$, $`x_3=F`$), $`F_i`$ and $`\varphi `$ are nonlinear functions of their own arguments. In eqs. (2)-(4) the terms containing $`D_C`$, $`D_P`$ and $`D_F`$ coefficients describe ”diffusion” in the space of the exchange rate alteration variables $`C`$, $`P`$, $`F`$, i.e., the distribution of $`K`$ over these variables, and the $`\varphi (K)`$ terms describe the dependence of $`\frac{dK}{dt}`$ upon the processes regulating the velocity of the exchange rate alteration (in particular, a currency exchange rate relaxation, bistable states existence etc.).
3. In the simplest case chose $`{}_{C}{}^{}=0`$, that is, consider the equation
$$\frac{dK}{dt}=\frac{K}{t}+\frac{K}{C}\frac{dC}{dt}+\frac{K}{P}\frac{dP}{dt}+\frac{K}{F}\frac{dF}{dt}=0$$
(6)
The equality $`I_C`$=0 corresponds to neglecting by the currency exchange rate alteration smoothing processes influence on $`\frac{dK}{t}`$. The $`\frac{K}{C}`$, $`\frac{K}{P}`$ and $`\frac{K}{F}`$ derivatives in (6) are determined, in difference with kinetic theory in statistical physics, by eqs.(4) and the derivatives $`\dot{x}_i`$ are governed by (5).
4. Consider now, for the sake of simplicity, the case when the influence of $`P`$ and $`F`$ on the currency exchange rate alteration is very small and may not be taken into account. Then
$$\frac{K}{t}+\frac{K}{C}\frac{dC}{dt}=0$$
(7)
To define $`\frac{K}{C}\frac{dC}{dt}`$ we use a reduced equations (1)-(5)
$$\frac{dC}{dt}=a_0^tK(t)𝑑t$$
(8)
$$\frac{K}{C}=\alpha $$
(9)
where $`\alpha `$ in (8) is a velocity of changing in investments attributed to the sum of the currency exchange rate alteration over the time interval of $`t`$ and $`\alpha `$ in (9) is the changing in the currency exchange rate alteration per investments unit.
Assume then that these values do not vary significantly, that is $`\alpha =const`$ and $`\frac{dC}{dt}=b=const`$. In this case we will have $`\frac{K}{t}=K_1=ab`$ and $`K=K_1t+K_2`$ ($`K_2=const.`$). So $`K`$ will increase or decrease, depending upon the sign of $`K_1`$. In case of (8), differentiating (7) in time and using (13) from paper and (9) yield
$`{\displaystyle \frac{^2K}{t^2}}+\omega _0^2K=0`$ (10)
where $`\omega _0=\sqrt{a\alpha }`$. The solution of (10) has the form
$`K=K_0\mathrm{sin}\omega _0t`$ (11)
where $`K_0`$ is constant. So with the assumption made, the currency exchange rate alteration $`K`$ varies periodically in time with the frequency of $`\sqrt{a\alpha }`$. If the dependence of $`K`$ upon the investments made is weak, the frequency of oscillations will be small and K will take the form ($`\omega _0t1`$)
$$K=K_0\sqrt{a\alpha }t+\mathrm{}\sqrt{a\alpha }t1$$
At times large enough, nevertheless, the periodicity of the currency exchange rate varying with time will be observed.
5. Introduce now into the right-hand side of eq.(6) a parameter characterizing the deviation of $`K`$ from an equilibrium: $`I_c=\frac{KK_0}{\tau (K)}`$, where $`\tau (K)`$ is the relaxation time for the currency exchange rate alteration if diverged from a stationary state $`K_0`$ (in general case $`\tau `$ is a function of $`K`$, $`\dot{K}`$, $`C`$, $`F`$ etc.)
Then instead of (7) we have
$$\frac{K}{t}+\frac{K}{C}\frac{dC}{dt}=\frac{KK_0}{\tau (K)}$$
(12)
Assuming $`K_0=const`$ and $`\tau =const`$ differentiating in time, if eqs. (8) and (9) are valid, yields:
$`\ddot{K}+\omega _0^2K{\displaystyle \frac{1}{\tau }}\dot{K}=0`$ (13)
whose solution is
$`K=k_0e^{\frac{t}{\tau }}\mathrm{cos}\omega _0t`$ (14)
In this case the oscillations of the currency exchange rate changing will damp tending to zero when $`t\tau `$.
6. Take into account in equation (12) the difference $`P`$ between total demand and supply influence on the exchange rate alteration, determining $`\frac{dP}{dt}`$ and $`\frac{K}{P}`$ from the equations (a special case of (1) and (5))
$`a){\displaystyle \frac{dP}{dt}}=\gamma \dot{K},{\displaystyle \frac{K}{P}}=\beta K^2`$ (15)
$`b){\displaystyle \frac{dP}{dt}}=\beta K^2,{\displaystyle \frac{K}{P}}=\gamma \dot{K}`$ (16)
with linear (a) or nonlinear (b) dependence of $`\frac{dP}{dt}`$ upon $`\dot{K}`$ or $`K`$ correspondingly. Then instead of (12) we get
$$\frac{K}{t}+a\alpha _0^tK(t^{})𝑑t^{}+\beta \gamma K^2\dot{K}=\frac{KK_0}{\tau (K)}$$
(17)
or, after differentiating with respect to time
$$(1+\beta \gamma K^2)\ddot{K}(\frac{1}{\tau }2\beta \gamma K\dot{K})\dot{K}+\omega _0^2K=0$$
(18)
If $`\beta \gamma K^21`$ and $`2\beta \gamma K\dot{K}\frac{1}{\tau }`$, then (17) reduces to
$$\ddot{K}+2\delta (K\dot{K})\dot{K}+\omega _0^2K=0$$
(19)
where
$`\delta (K\dot{K})=\beta \gamma K\dot{K}`$ (20)
Eq. (18) is an equation of oscillations with a nonlinear friction $`2\delta (K\dot{K})`$ and its solutions contains all the peculiarities of nonlinear systems. Take into consideration in (17) the term $`\frac{}{P}(D_P\frac{}{P})K`$, assuming $`D_P=const`$. This yields a change in coefficient at in (19) in the case (a)
$`(1+\beta \gamma K^2)\ddot{K}`$ $``$ $`({\displaystyle \frac{1}{\tau }}2\beta \gamma K\dot{K}+D_P\beta ^2K^2)\dot{K}+`$ (21)
$`+`$ $`\omega _0^2K=0`$ (22)
and appearance of a term containing $`\frac{d^3K}{dt^3}`$ in the case (b)
$`\gamma ^2D_P^2{\displaystyle \frac{d^3K}{dt^3}}(1+\beta \gamma K^2)\ddot{K}`$ (23)
$`2\beta (K\dot{K})\dot{K}\omega _0^2K=0`$ (24)
7. One may take into account in (17) the terms containing $`\frac{K}{F}\frac{dF}{dt}`$ assuming, for example, $`D_F=const`$, and (a special case of (1)-(5))
$`a){\displaystyle \frac{K}{F}}=b{\displaystyle }_0^tK(t^{})dt^{},{\displaystyle \frac{dF}{dt}}=d=const`$ (25)
$`b){\displaystyle \frac{K}{F}}=d=const,{\displaystyle \frac{dF}{dt}}={\displaystyle }_0^tK(t^{})dt^{}`$ (26)
Eqs. (21) then takes the from ((25a) case)
$`(1+\beta \gamma K^2)\ddot{K}({\displaystyle \frac{1}{\tau }}2\beta \gamma K\dot{K}+D_P\beta ^2K^2)\dot{K}+`$ (27)
$`[(\omega _0^2+bd)KD_F\beta ^2{\displaystyle _0^t}K(t^{})𝑑t^{}]=0`$ (28)
and after differentiating in $`t`$
$`(1+\beta \gamma K^2){\displaystyle \frac{d^3K}{dt^3}}({\displaystyle \frac{1}{\tau }}2\beta \gamma K\dot{K}+D_P\beta ^2K^2)^{}\dot{K}+`$ (29)
$`+({\displaystyle \frac{1}{\tau }}2\beta \gamma K\dot{K}+D_P\beta ^2K^2)\ddot{K}+`$ (30)
$`+[(\omega _0^2+bd)\dot{K}D_F\beta ^2K(t)]+2\beta \gamma K\dot{K}\ddot{K}=0`$ (31)
Using (25b) one obtains
$`(1+\beta \gamma K^2)\ddot{K}({\displaystyle \frac{1}{\tau }}2\beta \gamma K\dot{K}+D_P\beta ^2K^2)\dot{K}+`$ (32)
$`+(\omega _0^2+bd)K=0`$ (33)
Consider finally a case of nonlinear dependence of $`I_C`$ upon $`K`$. Let, remaining the assumptions of the previous paragraphs concerning the dependencies of $`K`$ upon $`D`$, $`C`$ and $`F`$, $`I_C`$ has the form
$`I_C={\displaystyle \underset{0}{\overset{t}{}}}(l_0`$ $``$ $`lK^2)Kdt+{\displaystyle \frac{KK_0}{\tau }}+D_P{\displaystyle \frac{^2K}{P^2}}+`$ (35)
$`+D_C{\displaystyle \frac{^2K}{C^2}}+D_F{\displaystyle \frac{^2K}{F^2}}`$
The choice of (35) corresponds to nonlinear type of $`K(t)`$ relaxation. Substituting $`I_C`$ from (35) (e.g., for (25b) case) into (2) yields (using the proper equations for derivatives of $`K`$, $`C`$, $`P`$, $`F`$):
$`(1+\beta \gamma K^2)\ddot{K}`$ $`+(2\beta \gamma K\dot{K}D_P\beta ^2K^2{\displaystyle \frac{1}{\tau }})\dot{K}+`$ (37)
$`+[(\omega _0^2+bd)(l_0lK^2)]K=0`$
Eq. (37) has a bifurcation point at
$`K=\pm {\displaystyle \frac{1}{\sqrt{l}}}\sqrt{l_0\omega _0^2bd}`$ (38)
(which corresponds to the appearance of a bistable state for the currency exchange rate alteration) and a number of other interesting peculiarities as well (in particular, at $`K=\pm \sqrt{\frac{1}{|\gamma \beta |}}`$, $`(\gamma \beta <0)`$,
$`atK={\displaystyle \frac{1}{D_P\beta ^2}}(\beta \gamma \dot{K}\pm \sqrt{(\beta \gamma \dot{K})^2+{\displaystyle \frac{2D_P\beta ^2}{\tau }}},`$ (39)
$`at\dot{K}={\displaystyle \frac{D_P\beta ^2K\pm \sqrt{D_P\beta ^2K^28\beta \gamma (\omega _0^2+bdl_0lK^2)}}{4\gamma \beta }},`$ (40)
etc.)
8. Note, that wide opportunities to chose the equations for $`\frac{K}{C}`$, $`\frac{K}{P}`$, $`\frac{K}{F}`$ and $`\frac{dC}{dt}`$, $`\frac{dP}{dt}`$, $`\frac{dF}{dt}`$ and $`I_C`$ are not exhausted by the selection used in the example considered and the concrete forms of the equations are determined by the state of the economic processes and correlations between them. For example, one may treat as $`x_i`$ variables in the case considered realized $`Y_e`$ and produced $`Y_l`$ gross national products, the total sum of money in use, the number of workable population $`N`$, absolute level of unemployment $`\delta N`$, the part of gross national product used by the state (these variables were used by Bystrai () while defining the economic entropy.
CONCLUSION
A general algorithm of describing economic processes basing on the equations of statistical physics of open systems developed by Kobelev L.Ya. et al. was used to describe the time dependence of the national currency exchange rate as a function of national dividend, investments size and difference between total demand for the goods and its supply. A number of nonlinear differential equations describing the time dependence of exchange rate (in particular, oscillations in different cases) were obtained.
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# COSMIC RAY SIGNATURES OF MASSIVE RELIC PARTICLES
## 1 Introduction
The only massive particles in the Standard Model to have survived from the Big Bang are nucleons — protons and (bound) neutrons — along with a commensurate number of electrons to yield the observed charge neutrality of the universe.<sup>a</sup><sup>a</sup>aWe know now that massive relic neutrinos contribute at least as much as the luminous component of nucleons to the present energy density. However they are unlikely to be the dominant component of the dark matter, based on arguments concerning structure formation. Considerations of primordial nucleosynthesis restrict the nucleonic contribution to the density parameter to $`\mathrm{\Omega }_\mathrm{N}<0.1`$ and it is widely accepted that the dark matter in galaxies and clusters which contributes $`\mathrm{\Omega }_{\mathrm{DM}}>0.3`$ is non-nucleonic and probably composed of a new stable relic particle. There are many candidates for the identity of this particle but the most popular notion is that it is associated with the new physics beyond the Standard Model necessary to stabilize the hierarchy between the Fermi scale, $`G_F^{1/2}300`$ GeV, and the Planck scale, $`G_N^{1/2}10^{19}`$ GeV. In particular theories of (softly broken) low energy supersymmetry (SUSY) typically imply that the lightest SUSY partner is a neutralino with mass of order the Fermi scale, which is absolutely stable if the discrete symmetry termed $`R`$-parity is exactly conserved. Interestingly enough the relic abundance of such a weakly interacting particle which was in thermal equilibrium in the early universe can account for the dark matter.
In supergravity theories, there is a new energy scale of $`𝒪(10^{11})`$ GeV — the geometric mean of the Fermi and Planck scales. This is the scale of the ‘hidden sector’ in which SUSY is broken through gaugino condensation induced by a new strong interaction, and communicated to the visible sector through gravitational interactions. Following the emergence of superstrings (for which $`N=1`$ supergravity is the effective field theory) it was realised $`^\mathrm{?}`$ that the hidden sector can also serve to confine fractionally charged states which are a generic prediction $`^\mathrm{?}`$ of string theory. This avoids a serious conflict with the unsuccessful experimental searches for fractional charges but necessarily implies the existence of (integrally charged) bound states with mass of $`𝒪(10^{11})`$ GeV. In a specific construction with $`SU(5)U(1)`$ unification, it was noted $`^\mathrm{?}`$ that most such states would be short-lived but that the lightest such state would only decay through non-renormalizable operators of dimension $`8`$ and thus have a lifetime exceeding the age of the universe. This introduces a new candidate for the constituent of the dark matter — named “cryptons” — interestingly similar to nucleons which too are bound states of fractional charges and can only decay through non-renormalizable operators.
However, just as with nucleons, their cosmological origin is a puzzle. If such particles were ever in thermal equilibrium their relic abundance would have been excessive since their self-annihilations are rather inefficient. For nucleons the problem is just the opposite and their very existence today requires an out-of-equilibrium origin. If the same were true of cryptons, their relic abundance may well have a cosmologically interesting value.<sup>b</sup><sup>b</sup>bIt has recently been noted $`^\mathrm{?}`$ that particles with mass of $`𝒪(H_{\mathrm{inf}})10^{13}`$ GeV — also dubbed “wimpzillas” — can be created with a cosmologically interesting abundance through quantum vacuum fluctuations during inflation or during the subsequent (re)heating process. It is then interesting to ask what the observational signatures of such particles might be.
Reviving an old suggestion $`^\mathrm{?}`$, we recognised $`^\mathrm{?}`$ that the most sensitive probe would be in extremely high energy cosmic rays (EHECR), specifically in the flux of high energy neutrinos which would necessarily be created by crypton decays. The best constraint we obtained followed from the upper limit on deeply penetrating air showers set by the Fly’s Eye atmospheric fluorescence experiment; this implied that such particles must have a lifetime exceeding $`10^{18}`$ yr if they are an important constituent of the dark matter. As this was close to the theoretically expected lifetime in the “flipped” $`SU(5)`$ model, I was optimistic enough to suggest in a conference talk $`^\mathrm{?}`$ that “ …some improvement of these experimental sensitivities can rule out (or detect!) such particles”.
Just a few months later the Fly’s Eye array detected $`^\mathrm{?}`$ an event, consistent with a proton primary, but with an energy of $`(3.0\pm 0.9)\times 10^{11}`$ GeV. This was well above the Greisen-Zatsepin-Kuzmin (GZK) cutoff $`^\mathrm{?}`$ energy of $`5\times 10^{10}`$ GeV, beyond which resonant photopion production losses on the cosmic microwave background should limit the propagation distance of any such strongly interacting particle to less than about a hundred Mpc. Over a dozen such events have been detected subsequently by the Akeno airshower array (AGASA) as well as HiRes, the successor to Fly’s Eye, so the absence of the GZK cutoff $`^\mathrm{?}`$ is now well established. However contrary to the expectation that such high energy particles, being essentially undeflected by the weak intergalactic magnetic fields, should point back to their sources, the observed distribution on the sky $`^\mathrm{?}`$ is consistent with isotropy. This is quite baffling given that that only a few astrophysical sites (active galactic nuclei or the extended lobes of radio galaxies) are capable of accelerating such particles, even in principle, and there are none $`^\mathrm{?}`$ along the arrival directions within the propagation range. Hence it is generally acknowledged $`^\mathrm{?}`$ that there is no “conventional” astrophysical explanation for the observed EHECR.
## 2 EHECR from decaying dark matter
Faced with the above conundrum, some authors have resorted to desperate measures, e.g. postulating that the intergalactic magnetic field may be a thousand times stronger than usually believed, so capable of isotropising particles from a nearby active galaxy. However, following from our previous discussion, there is a natural explanation $`^\mathrm{?}`$ for both the observed isotropy and absence of the GZK cutoff if the EHECR originate from the decays of metastable cryptons which are part of the dark matter.<sup>c</sup><sup>c</sup>cThis was independently proposed by Berezinsky, Kachelrieß and Vilenkin $`^\mathrm{?}`$, without, however, a specific particle candidate in mind. Kuzmin and Rubakov $`^\mathrm{?}`$ also made a qualitative suggestion that EHECR may originate from relic particle decays, however they did not make the crucial observation that such particles would be highly concentrated in our Galactic halo. This is because such particles will behave as cold dark matter (CDM) and hence cluster in the halo of the Milky Way with a concentration $`10^4`$ times higher than the cosmic average. The local flux of EHECR will thus be dominated by decays of cryptons in the halo, implying two distinct observational tests of the hypothesis. First, the energy spectrum and cosmposition (nucleons, gammas, neutrinos) beyond the GZK cutoff will be determined $`^\mathrm{?}`$ essentially by the physics of crypton decays. Second, there will be a small anisotropy $`^\mathrm{?}`$ in the arrival directions of EHECR since we are located $`8`$ kpc away from the centre of the Galaxy and should therefore observe more particles arriving from that direction than from the anticentre. There may also be measurable correlations between arrival times of high energy nucleons, gammas and neutrinos.
### 2.1 Particle candidates
As noted above, the possibility of metastable relic particles with mass of $`𝒪(10^{11})`$ GeV had been proposed $`^{\mathrm{?},\mathrm{?}}`$ before the observations of EHECR beyond the GZK cutoff. An updated discussion $`^\mathrm{?}`$ of such particles in string/M-theory confirms that cryptons are indeed favoured over other possibilities such as the Kaluza-Klein states associated with new compact dimensions (which are too short-lived). The most likely candidate is still a neutral pion-like ‘tetron’ composed of four constituents, with a minimum lifetime of
$$\tau _X\frac{1}{m_X}\left(\frac{M}{m_X}\right)^{10},$$
(1)
where $`m_X10^{1213}`$ GeV, and the scale $`M`$ of suppression of non-renormalizable terms is of $`𝒪(10^{18})`$ GeV.<sup>d</sup><sup>d</sup>dOther authors $`^\mathrm{?}`$ have also considered string candidates for superheavy dark matter, and discussed $`^\mathrm{?}`$ the confinement of fractionally charged particles into baryon-like states in the hidden sector and the discrete symmetries required to ensure their longevity. Thus both the mass and lifetime of the candidate particle are motivated by topical physical considerations. This is in contrast to other proposals $`^{\mathrm{?},\mathrm{?}}`$ where the mass scale is not given any physical motivation and the decays are presumed to be mediated by unspecified instanton or quantum gravity effects so as to yield a suitably long lifetime.
### 2.2 Calculation of decay spectrum
Nevertheless all such proposals have a common phenomenology in that regardless of the decay mechanism, the spectra of the decay products is essentially determined by the physics of QCD fragmentation $`^\mathrm{?}`$ and has no major astrophysical uncertainties. In particular given that the propagation distance in the halo is $`<100`$ kpc, much shorter than the GZK range of $`100`$ Mpc, the EHECR spectrum at Earth will be the same as the decay spectrum (apart from the decay photons which will be degraded through scattering on background photon fields). Of course the decay mode (e.g. 2-body vs many-body) may well play an important role. However in our picture the decaying particle is a singlet under Standard Model interactions and has a mass which significantly exceeds the Fermi scale so the inclusive spectra of final state nucleons, photons and neutrinos should be relatively insensitive to the precise decay channel.
We can thus imagine that we have say a $`e^+e^{}`$ collider at our disposal with a centre-of-mass energy $`\sqrt{s}`$ sufficient to create a supermassive particle such as a crypton, rather than just a $`Z^0`$ as at LEP (Figure 1). This then decays into quarks and gluons which initiate multi-parton cascades through gluon bremsstrahlung. These finally hadronize to yield high multiplicity jets when the momentum scale of the process drops below $`\mathrm{\Lambda }_{\mathrm{QCD}}`$. In the present context we are only interested in the final yields of nucleons, photons and neutrinos into which all the produced hadrons will decay. The production of different hadron species is quantified by their respective ‘total fragmentation functions’ $`F^h(x,s)=\sigma _{\mathrm{tot}}^1\mathrm{d}\sigma /\mathrm{d}x`$, viz. the probability distributions for their inclusive production as a function of the scaled hadron energy $`x2E_h/\sqrt{s}`$. These can be factorized as the sum of contributions from different primary partons $`i=u,d,\mathrm{},g`$:
$$F^h(x,s)=\underset{i}{}\frac{\mathrm{d}z}{z}C_i(s;z,\alpha _\mathrm{s}(s))D_i^h(x/z,s),$$
(2)
where $`C_i`$ are the ‘coefficient functions’ dependent on the production process, and $`D_i^h`$ is the ‘universal fragmentation function’ for parton $`i`$ hadron $`j`$. The essential physics in the hard parton cascade is the logarithmic evolution of the strong coupling $`\alpha _\mathrm{s}`$ with energy and sophisticated techniques have been developed $`^\mathrm{?}`$ to handle divergences associated with collinear and soft gluon emission. The formation of the final hadrons is however an inherently non-perturbative process and can only be described at present by empirical models encoded in Monte Carlo event generators, e.g. JETSET $`^\mathrm{?}`$ based on the ‘string fragmentation’ model, or HERWIG $`^\mathrm{?}`$ based on the ‘cluster hadronization’ model. These also account for the subsequent decay of the hadrons into the observed particles, taking into account all experimentally measured branching ratios, resonances etc, so can make detailed predictions of measurable quantities.
Although the fragmentation functions are not perturbatively calculable, their evolution as a function of the momentum scale is governed by the Dokshitzer-Gribov-Lipatov-Altarelli-Parisi (DGLAP) equation $`^\mathrm{?}`$
$$s\frac{}{s}D_i^h(x,s)=\underset{j}{}_x^1\frac{dz}{z}P_{ji}(z,\alpha _\mathrm{s}(s))D_i^h(x/z,s),$$
(3)
where $`P_{ji}`$ are the ‘splitting functions’ for the process parton $`ij`$. Thus by measuring the fragmentation functions at one momentum scale, one can evaluate them at another scale. As seen in Figure 2 the DGLAP equations predict violations of ‘Feynmann scaling’ — a softening of the spectrum with increasing energy — in good agreement with data, in this case measured at PETRA ($`\sqrt{s}=22`$ GeV) and LEP ($`\sqrt{s}=91.2`$ GeV).
At small values of $`x`$, multiple soft gluon emission gives rise to higher-order corrections, which turn out to be resummable by altering the scale in the DGLAP equation (3) from $`sz^2s`$; this yields a simple Gaussian function in the variable $`\xi \mathrm{ln}(1/x)`$:
$$xF^h(x,s)\mathrm{exp}\left[\frac{1}{2\sigma ^2}(\xi \xi _\mathrm{p})^2\right],$$
(4)
which has a characteristic peak at $`\xi _\mathrm{p}\mathrm{ln}s/4`$, with width $`\sigma (\mathrm{ln}s)^{3/4}`$. Including ‘next-to-leading’ corrections to these predictions yields the ‘modified leading log approximation’ (MLLA) $`^{\mathrm{?},\mathrm{?}}`$ which accounts very well for the shape of the observed fragmentation functions at small $`x`$. In comparing with data one has to further assume ‘local parton hadron duality’ (LPHD) $`^\mathrm{?}`$, viz. that the hadron distribution is simply proportional to the parton distribution. Thus the prediction cannot distinguish between the individual hadronic species. Moreover although this kinematic region dominates the total multiplicity, it accounts for only a small fraction of the energy in the cascade, hence the MLLA spectrum cannot be correctly normalized.
With this background, we can review what has been done so far to explain the EHECR data in terms of decaying halo particles. Berezinsky et al $`^\mathrm{?}`$ adopted the gaussian approximation (4) to MLLA to infer the spectrum of nucleons from the decay of a particle of mass $`m_X`$. Although this approximation is only valid for small $`x`$ ($`0.1`$), these authors nevertheless normalized it by requiring that $`_0^1dxxF^h(x)=f_\mathrm{N}`$ where $`f_\mathrm{N}0.05`$ is the assumed fraction of the decaying particle mass transferred to nucleons (on the basis of $`Z^0`$ decay data). The rest is assumed to go into pions which decay to yield photons and neutrinos, with neutral pions taking a third of the total energy. Figure 3 shows their fit to the EHECR data (multiplied by $`E^3`$ for clarity) with decaying particles of mass $`m_X=10^{13}`$ GeV which contribute a fraction $`\xi _X`$ of the CDM density in our halo (taken to have a radius of 100 kpc). Their adopted normalization then implies a particle lifetime $`\tau _x/t_0=2\times 10^{10}\xi _X`$, where $`t_012\times 10^9`$ yr is the age of the universe. Note that the constant suppression with energy of the nucleon flux with respect to photons (and neutrinos) follows from the assumed proportionality of the nucleon and pion fragmentation functions.
Because of the above problems with the MLLA spectrum, we had already considered and rejected this convenient approximation and chosen instead to embark on a time-consuming calculation of the fragmentation functions in the kinematic region of relevance to the data, using the HERWIG $`^\mathrm{?}`$ event generator. In doing so we were initially motivated to test an argument due to Hill $`^\mathrm{?}`$ that the fragmentation spectrum at large $`x`$ should be $`(1x)^2`$. By normalizing to the total multiplicity and demanding energy conservation Hill was then able to obtain an empirical fragmentation function which could be fitted to extant data. Assuming the multiplicity to be $`s^{1/4}`$ as in the naive statistical model of jet fragmentation, this was
$$F_h=\frac{15}{16}x^{3/2}(1x)^2,$$
(5)
while by adopting the leading-log QCD prediction for the multiplicity ($`\mathrm{exp}\sqrt{\mathrm{ln}(s/\mathrm{\Lambda }^2)}`$), he obtained
$$F^h=N(b)\frac{\mathrm{exp}[b\sqrt{\mathrm{ln}(1/x)}](1x)^2}{x\sqrt{\mathrm{ln}(1/x)}}.$$
(6)
By fitting this form to PETRA data, Hill found $`N(b)=0.08`$ and $`b=2.6`$. On the basis of the same data he assumed that 3% of the hadronic jets form nucleons and the other 97% are pions which decay into photons and neutrinos.
Many authors $`^\mathrm{?}`$ who have investigated the annihilation of GUT-scale relic topological defects (TD) as the source of EHECR have used these expressions to estimate the fluxes. However this is clearly inaccurate since there would be large scaling violations (see Figure 2) in going from the PETRA energy scale of 22 GeV up to the very much higher GUT energy scale. This is just what our calculations $`^\mathrm{?}`$ using HERWIG demonstrate. The functional form (6) continues to provide a good fit to the fragmentation function of nucleons but as shown in Figure 4 the spectrum becomes significantly softer with increasing particle mass, e.g. for $`m_X=10^{13}`$ GeV, the normalization $`N(b)`$ drops to $`0.0078`$ (with $`b=2.8`$). Thus TD models $`^\mathrm{?}`$ of EHECR which use the fragmentation functions (5,6) overestimate nucleon production by a factor of $`10`$ at high energies.
The EHECR spectrum in the energy range $`10^{911}`$ GeV is well fitted $`^\mathrm{?}`$ as the sum of two power-laws — the extrapolation of the $`E^{3.3}`$ spectrum from lower energies and a new flatter component $`E^{2.7}`$ which dominates above $`10^{10}`$ GeV. (The AGASA data $`^\mathrm{?}`$ gives the slope of the new component as $`2.78_{0.33}^{+0.25}`$.) There is some indication $`^\mathrm{?}`$ that the composition also changes from iron-group nuclei to protons at this energy. In Figure 5 we see that the HERWIG generated spectrum $`^\mathrm{?}`$ for a decaying particle mass of $`10^{12}`$ GeV is indeed in reasonable agreement with this new component of cosmic rays. Our normalization requires a lifetime $`\tau _X/t_0=3\times 10^9\xi _X`$ in the notation of Berezinsky et al $`^\mathrm{?}`$, i.e. a factor of $`6`$ smaller. (This is because they seem to have normalized to photons rather than nucleons at $`10^{10}`$ GeV; see Figure 3).
However our own calculations $`^\mathrm{?}`$ suffer from two problems. The first, which we were initially unaware of, is that HERWIG has a known tendency $`^\mathrm{?}`$ to overproduce baryons at large $`x`$ (essentially due to the fragmentation of the leading quark from the initial hard process, viz. particle decay in the present case). Although the overall multiplicity is correctly predicted at LEP energies (e.g. 0.953 protons per event vs $`0.98\pm 0.09`$ observed in $`Z^0`$ decay), HERWIG overproduces nucleons at $`x>0.3`$ by a factor of $`23`$. Secondly, in studying the evolution of the fragmentation function to very high energies we should take into account that the running of the strong coupling $`\alpha _\mathrm{s}`$ would be altered above $`10^3`$ GeV when SUSY particles begin to be excited from the vacuum.
Both these issues have been addressed recently in unpublished work $`^\mathrm{?}`$ by Rubin. As shown in Figure 6 he finds that agreement of HERWIG with LEP data $`^\mathrm{?}`$ on baryon production is significantly improved by isotropizing the decay of the hadron cluster formed from the hard process (rather than have the hadron leave the cluster along the direction of the initial quark). An additional suppression of the probability for cluster decay by $`20\%`$ improves the fit further. To take SUSY particles into account, he evolves the DGLAP equation (3) from LEP energies upwards (with the initial sparton fragmentation functions calculated with the PYTHIA $`^\mathrm{?}`$ event generator). The SUSY $`\beta `$ function for $`\alpha _\mathrm{s}`$ is used, with the flavour thresholds corresponding to the sparticle spectrum of a typical minimal supergravity model (with parameters $`M_0=800`$ GeV, $`m_{1/2}=200`$ GeV, $`A_0=0`$, $`\mathrm{tan}\beta =10`$, sgn$`(\mu )=+`$).
Figure 7 shows his results for both the non-SUSY and SUSY cases. The high energy “bump” in our proton spectrum (see Figure 4) has been erased but the non-SUSY spectrum continues to reproduce the shape of the data for a decaying halo particle mass of $`𝒪(10^{12})`$ GeV. The effect of including the effects of SUSY on the evolution of the parton cascade is to flatten the spectrum further so that a $`10`$ times larger mass is still acceptable. We note that the spectral shape differs considerably from the “SUSY-QCD” spectrum calculated by Berezinsky and Kachelrieß $`^\mathrm{?}`$ using MLLA. This is not unexpected since as emphasized earlier, this approximation is unjustified at large $`x`$ so cannot be normalized (as these authors do) to the energy released in the decay. Moreover their assumption of an energy-independent ratio between nucleons and pions is invalid; as is evident from Figure 7 this ratio increases with energy. <sup>e</sup><sup>e</sup>eHowever it never exceeds unity as in our previous work $`^\mathrm{?}`$ using HERWIG which suffered from overproduction of hard baryons and gave an incorrect prediction of the $`p/\nu `$ ratio.
## 3 Conclusions
Although some progress has been made in sharpening the spectral predictions of the decaying halo particle model for EHECR, much work still needs to be done. The calculations so far have assumed the simplest decay channel — into two partons. However non-renormalizable operators are in fact likely to induce many-body decays. The effects of supersymmetry also need to be investigated more carefully, e.g. the effects of varying the SUSY parameters and inclusion of sparticle decay channels. Nevertheless it is already clear that the general trend in the EHECR data can be accounted for by this hypothesis, if the particle mass is $`m_X10^{1213}`$ GeV and its lifetime is $`\tau _X10^{16}\mathrm{yr}(\xi _X/3\times 10^4)`$, so that even with a very long lifetime such particles need constitute only a tiny fraction $`\xi _X`$ of the halo CDM. It is also clear that TD models $`^\mathrm{?}`$, in which $`m_X`$ corresponds to the GUT-scale, are already ruled out by the spectral data.
The next generation of large area cosmic ray, gamma-ray and neutrino observatories (Auger, Amanda, Antares, …) is now under construction so it is important to refine these calculations in order to make specific predictions for the expected fluxes. We emphasize that previous estimates of high energy gamma-ray and neutrino fluxes from TD $`^\mathrm{?}`$ are based on the Hill fragmentation functions (5,6), while other work $`^\mathrm{?}`$ use the (M)LLA spectrum (4) or its SUSY variant. Blasi $`^\mathrm{?}`$ has calculated in detail the flux of $`\gamma `$-rays in the decaying halo particle model but he too uses the Hill and the MLLA spectra. All these approximations are inapplicable at the high energies of interest as explained earlier, and moreover the spectra of pions are not simply proportional to that of nucleons as assumed. Hence it is clear that all these estimates are unreliable. It is essential that further work use the physically more realistic approach to calculating fragmentation spectra outlined above in order to devise definitive experimental tests $`^\mathrm{?}`$ of the decaying particle hypothesis.
## Acknowledgments
I wish to thank the organisers of this enjoyable conference, and Neil Rubin for discussions and permission to quote his unpublished work.
## References
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# The Yang-Mills Measure in the Kauffman Bracket Skein Module
## 1. Introduction
Since the introduction of quantum invariants of $`3`$-manifolds the fact that they are only defined at roots of unity has been an obstruction to analyzing their properties. One approach has been to study the perturbative theory of quantum invariants . However, there is ample evidence quantum invariants of three manifolds exist as holomorphic functions on the unit disk, that diverge everywhere on the unit circle but at roots of unity . This paper takes a step towards seeing that this holds in general. The Yang-Mills measure is the path integral on a topological quantization of the $`SU(2)`$-characters of the fundamental group of a closed surface. The measure displays the same convergence properties as are expected of quantum invariants of $`3`$-manifolds.
The Yang-Mills measure in the Kauffman bracket skein algebra of a cylinder over a closed surface $`\mathrm{\Sigma }_g`$ is a local, diffeomorphism invariant trace. It quantizes the symplectic measure on the space $`(\mathrm{\Sigma }_g)`$ of conjugacy classes of representations of the fundamental group of $`\mathrm{\Sigma }_g`$ into $`SU(2)`$. The definition of the symplectic structure and formulas for its computation are in . The volume of $`(\mathrm{\Sigma }_g)`$ was computed by Witten in in two ways: via the equivalence of two computations in quantum field theory, and by noting that the symplectic measure is equal to the measure coming from Reidemeister torsion. In Witten’s setting the Yang-Mills measure is a path integral in a lattice model of field theory that depends on area. Forman gave a direct proof that Witten’s measure converges to the symplectic measure as the area goes to zero.
Alekseev, Grosse and Schomerus conceived of a method of constructing lattice gauge field theory based on a quantum group. This idea was further developed by Buffenoir and Roche who gave a construction of the algebra, its Wilson loops and a trace called the Yang-Mills measure that were completely analogous to Witten’s construction. Their theory is topological when the area is set to zero.
The method of constructing the algebras in is combinatorial and based on generators and relations. We gave a new construction of lattice gauge field theory in that is “coordinate free”. The connections form a co-algebra and the product on the gauge fields is a convolution with respect to the co-multiplication of connections. This allows the structure of the observables to be elucidated. We found working over formal power series, basing the theory on quantum $`sl_2`$, that the observables are the Kauffman bracket skein algebra of a cylinder over a regular neighborhood of the $`1`$-skeleton. In we recover the same result working over the complex numbers.
These considerations lead one to expect that the Yang-Mills measure exists as a trace on the Kauffman bracket skein algebra of a closed surface. In this paper we affirm this fact, with the only reservation that if the deformation parameter $`t`$ is a generic point on the unit circle, then the measure does not converge. However, at roots of unity the trace exists and is well known. Furthermore, at $`t=1`$ the Yang-Mills measure is the symplectic measure on $`(\mathrm{\Sigma }_g)`$.
This paper is organized as follows. Section 2 recalls definitions, associated formulas and the algebra structure of the Kauffman bracket skein module of a cylinder over a surface. In section 3 the Yang-Mills measure is defined for compact surfaces with boundary, and is proved to be a trace. In section 4, working with the parameter $`t`$ such that $`|t|1`$, we obtain estimates for the absolute value of the tetrahedral coefficients and use these to show that the Yang-Mills measure can be defined for closed surfaces. In section 5 we define and investigate the measure when $`t`$ is a root of unity.
## 2. Preliminaries
Let $`M`$ be an orientable $`3`$-manifold. A framed link in $`M`$ is an embedding of a disjoint union of annuli into $`M`$. Framed links are depicted by showing the core of an annulus lying parallel to the plane of the paper (i.e. with blackboard framing). Two framed links in $`M`$ are equivalent if there is an isotopy of $`M`$ taking one to the other. Let $``$ denote the set of equivalence classes of framed links in $`M`$, including the empty link. Fix a complex number $`t0`$. Consider the vector space $``$ with basis $``$. Define $`S(M)`$ to be the smallest subspace of $``$ containing all expressions of the form $`\text{ }\text{}\text{ }t\text{ }\text{}\text{ }t^1\text{ }\text{}\text{ }`$ and $`+t^2+t^2`$, where the framed links in each expression are identical outside balls pictured in the diagrams. The Kauffman bracket skein module $`K_t(M)`$ is the quotient
$$/S(M).$$
Let $`F`$ be a compact orientable surface and let $`I=[0,1]`$. There is an algebra structure on $`K_t(F\times I)`$ that comes from laying one link over the other. Suppose that $`\alpha ,\beta K_t(F\times I)`$ are skeins represented by links $`L_\alpha `$ and $`L_\beta `$. After isotopic deformations, to “raise” the first link and “lower” the second, $`L_\alpha F\times (\frac{1}{2},1]`$ and $`L_\beta F\times [0,\frac{1}{2})`$. The skein $`\alpha \beta `$ is represented by $`L_\alpha L_\beta `$. This product extends to a product on $`K_t(F\times I)`$. We denote the resulting algebra by $`K_t(F)`$ to emphasize that it comes from viewing the underlying three manifold as a cylinder over $`F`$.
The notation and the formulas in this paper are taken from . However, the variable $`t`$ replaces $`A`$, and we use quantum integers
$$[n]=\frac{t^{2n}t^{2n}}{t^2t^2}.$$
When $`t=\pm 1`$, $`[n]=n`$. Note that $`\mathrm{\Delta }_n`$ from is equal to $`(1)^n[n+1]`$.
There is a standard convention for modeling a skein in $`K_t(M)`$ on a framed trivalent graph $`\mathrm{\Gamma }M`$. When $`\mathrm{\Gamma }`$ is represented by a diagram we assume blackboard framing. An admissible coloring of $`\mathrm{\Gamma }`$ is an assignment of a nonnegative integer to each edge so that the colors at trivalent vertices form admissible triples (defined below). The corresponding skein in $`K_t(M)`$ is obtained by replacing each edge labeled with the letter $`m`$ by the $`m`$-th Jones–Wenzl idempotent (see , or , p.136), and replacing trivalent vertices with Kauffman triads (see \[15, Fig. 14.7\]).
Recall the fusion identity:
$$\text{}a\text{ }b=\underset{c}{}(1)^c\frac{[c+1]}{\theta (a,b,c)}\text{}ca\text{ }ba\text{ }b$$
where the sum is over all $`c`$ so that the triples $`(a,b,c)`$ are admissible, i.e. $`a+b+c`$ is even, $`ab+c`$, $`ba+c`$, and $`ca+b`$. Value of $`\theta (a,b,c)`$ is given by equation (4) below. The fusion relation is satisfied in $`K_t(M)`$ unless $`t`$ is a root of unity other than $`\pm 1`$.
## 3. The Yang-Mills Measure in a Handlebody
Throughout this section we assume that $`t`$ is not a root of unity. The first result is well known and comes from Przytycki’s construction of examples of torsion in skein modules.
###### Lemma 1 (The Sphere Lemma).
Let $`s_c`$ be a skein represented by coloring a trivalent framed graph in the manifold $`M`$. Suppose further that there is a sphere embedded in $`M`$ which intersects the underlying graph transversely in a single point in the interior of an edge, and the color of that edge is not zero. Then $`s_c=0`$.
###### Proof.
Using the “light bulb trick” isotope the framed graph $`s_c`$ so that it is the same graph, but the framing on the edge intersecting the sphere has been changed by adding two kinks. Using the formula for eliminating a kink, notice that $`s_c`$ is a nontrivial complex multiple of itself. Ergo, $`s_c`$ represents zero in $`K_t(M)`$. ∎
Consider now $`K_t(\mathrm{\#}_gS^1\times S^2)`$, the Kauffman bracket skein module of the connected sum of $`g`$ copies of $`S^1\times S^2`$.
###### Proposition 1.
The skein module $`K_t(\mathrm{\#}_gS^1\times S^2)`$ is canonically isomorphic to $``$. The isomorphism is given by writing each skein as a complex multiple of the empty skein.
###### Proof.
This follows easily from theorems of Hoste and Przytycki . In the Kauffman bracket skein module of $`S^1\times S^2`$ is computed over $`[t,t^1]`$. This along with the results in on the Kauffman bracket skein module of a connected sum over rational functions in $`t`$, combined with the universal coefficient theorem stated in , proves the desired result.
We outline the actual isomorphism with the complex numbers. Choose a system of spheres in $`\mathrm{\#}_gS^1\times S^2`$ that cut it down to a punctured ball. Given a skein in $`\mathrm{\#}_gS^1\times S^2`$, represent it as a linear combination of colored, framed, trivalent graphs intersecting the spheres transversely in interior of edges, and so that each graph intersects any sphere at most once. This is done by fusing multiple edges passing through the same sphere. By the sphere lemma, we can assume the graphs miss the spheres. Now take the Kauffman bracket of the skein in the punctured ball to write it as a complex multiple of the empty skein. ∎
Given a handlebody $`H`$ of genus $`g`$ its double is $`\mathrm{\#}_gS^1\times S^2`$. There is a linear functional $`𝒴:K_t(H)`$ computed by taking the inclusion of $`H`$ into $`\mathrm{\#}_gS^1\times S^2`$ followed by taking the “ Kauffman bracket” as above. Let $`F`$ be a compact, oriented surface with boundary. Since $`F\times I`$ is a handlebody the linear functional
$$𝒴:K_t(F),$$
is defined. We call this the Yang-Mills measure.
Choose a trivalent spine of $`F`$. The admissible colorings of that spine form a basis for $`K_t(F)`$. The skein modules of the disk and annulus are exceptions; the first is spanned by the empty skein and the latter is described in Section 4. In terms of this basis the Yang-Mills measure is just the coefficient of the skein coming from labeling all the edges of the spine with $`0`$.
###### Proposition 2.
The Yang-Mills measure is a trace, that is
$$𝒴(\alpha \beta )=𝒴(\beta \alpha ).$$
Furthermore, the trace is invariant under the action of the diffeomorphisms of $`F\times I`$ on $`K_t(F)`$.
###### Proof.
Let $`L`$ be the link $`F\times \{1/2\}`$. The result of removing $`L`$ from the double of $`F\times I`$ is homeomorphic to the Cartesian product of the interior of $`F`$ with a circle. Given any skein in $`F\times I`$ we can represent it by a linear combination of framed links that miss $`L`$. Hence, the Yang-Mills measure factors through the skein module of $`F\times S^1`$. In $`F\times S^1`$ the skeins $`\alpha \beta `$ and $`\beta \alpha `$ are the same.
The group of diffeomorphisms of the handlebody $`F\times I`$ acts on $`K_t(F)`$ in the obvious way. If $`f:F\times IF\times I`$ is a diffeomorphism then it can be extended to $`Df:\mathrm{\#}_gS^1\times S^2\mathrm{\#}_gS^1\times S^2`$. Since the image of the empty skein under a diffeomorphism is the empty skein, the action of $`Df`$ on $`K_t(\mathrm{\#}_gS^1\times S^2)`$ is trivial. Therefore, $`𝒴(f(\alpha ))=𝒴(\alpha )`$. ∎
The final commonly used property of the Yang-Mills measure is that it is local. Suppose that $`k`$ is a proper arc in $`F`$. Cut $`F`$ along $`k`$ to get a surface $`F^{}`$. It is evident that if we write a skein $`\alpha `$ as a linear combination of admissibly colored graphs, each one intersecting $`k`$ transversely in at most a single point, then we can throw out any graph such that the edge intersecting $`k`$ carries a nonzero label. This yields a skein in $`F^{}`$, denoted by $`\alpha _k`$. Then $`𝒴(\alpha )=𝒴(\alpha _k)`$.
## 4. The Yang-Mills measure on a closed surface
Throughout this section assume that $`|t|1`$. In fact, we only work with $`0<t<1`$. However, it is evident that the same proofs are valid when $`1<t`$ since the formulas are symmetric in $`t`$ and $`t^1`$. Finally, the arguments extend to the case where $`t`$ is not real by replacing the estimates for $`t`$ by estimates of the absolute value of $`t`$.
Recall the Kauffman bracket skein algebra of a cylinder over an annulus $`A`$. The central core of the annulus can be seen as a link by giving it the blackboard framing. Let $`s_i`$ be the skein in the annulus which is the result of plugging the $`i`$-th Jones-Wenzl idempotent into the core. The skein module $`K_t(A)`$ is the vector space with basis $`\{s_i\}`$, where $`i`$ runs from zero to infinity. The product with respect to this basis is given by
(1)
$$s_is_j=\underset{q|ij|,\text{by 2’s}}{\overset{i+j}{}}s_q.$$
Use the Yang-Mills measure on $`K_t(A)`$ to define a pairing:
(2)
$$\alpha ,\beta =𝒴(\alpha \beta ).$$
The $`s_i`$ form an orthonormal basis with respect to (2). This pairing identifies the linear dual of $`K_t(A)`$ with series of the form $`_i\alpha _is_i`$, where the $`\alpha _i`$ are complex numbers. Note that:
$$\underset{i=0}{\overset{\mathrm{}}{}}\alpha _is_i,\underset{j=0}{\overset{n}{}}\beta _js_j=\underset{i=0}{\overset{n}{}}\alpha _i\beta _i.$$
Let $`\mathrm{\Sigma }_{g,1}`$ denote the compact orientable surface of genus $`g`$ with one boundary component. There is a pairing,
$$K_t(A)K_t(\mathrm{\Sigma }_{g,1})K_t(\mathrm{\Sigma }_{g,1})$$
given by representing the skein in $`K_t(\mathrm{\Sigma }_{g,1})`$ by a linear combination of links disjoint from some collar of the boundary, and plugging the skein in $`K_t(A)`$ into the collar. The Yang-Mills measure can then be applied to give a pairing,
(3)
$$K_t(A)K_t(\mathrm{\Sigma }_{g,1}).$$
This means there is a well defined map,
$$Y:K_t(\mathrm{\Sigma }_{g,1})K_t(A)^{}.$$
Topologize $`K_t(A)`$ by giving it the weak topology from $`Y`$. That is a sequence $`\sigma _nK_t(A)`$ is Cauchy if for every skein $`\alpha K_t(\mathrm{\Sigma }_{g,1})`$, the sequence of complex numbers $`Y(\alpha )(\sigma _n)`$ is Cauchy. A linear functional on $`K_t(\mathrm{\Sigma }_{g,1})`$ that comes from an element of this completion via the pairing (3) is called a distribution. It is interesting to note that the weak topology from $`Y`$ on $`K_t(A)`$ depends on the genus of the surface.
If $`g>1`$ there is a distribution on $`K_t(\mathrm{\Sigma }_{g,1})`$ which annihilates all “handle-slides” (Skeins that are represented by the difference of two links such that one can be obtained from the other by a slide across an imagined disk filling the boundary of $`\mathrm{\Sigma }_{g,1}`$). This linear functional descends to the skein module of the closed surface. Yang-Mills measure on a closed surface is the result of evaluating this distribution followed by a normalization.
Let’s think about what a skein in $`K_t(A)`$ would be like if it annihilated all handle-slides. Begin by writing it as $`_i\alpha _is_i`$ and solve for the $`\alpha _i`$. A simple computation shows that if $`\alpha _0`$ is zero then all $`\alpha _i`$ are zero. Normalize so that $`\alpha _0=1`$. Notice that if our skein annihilates handle-slides then the skein $`s_1+[2]s_0`$ must be annihilated. Using the rules for multiplication (1) we see that the coefficient $`\alpha _1`$ is equal to $`[2]`$. Continuing on this way we see that this skein has to be
$$\underset{i}{}(1)^i[i+1]s_i,$$
which is of course not in $`K_t(A)`$.
The first goal is to show that for $`g>1`$ the sequence of partials sums $`_{i=0}^n(1)^i[i+1]s_i`$ is Cauchy in the weak topology from $`Y`$, and so defines a distribution.
The notation $`\text{Tet}\left(\begin{array}{ccc}a& b& e\\ c& d& f\end{array}\right)`$ stands for the Kauffman bracket of the skein pictured in Figure 1 on the left. The explicit formula is given in . We also need the quantity $`\theta (a,b,c)`$ which is the Kauffman bracket of the colored graph on the right in Figure 1. In terms of quantum integers
(4)
$$\theta (a,b,c)=(1)^{\frac{a+b+c}{2}}\frac{[\frac{a+b+c}{2}+1]![\frac{a+bc}{2}]![\frac{b+ca}{2}]![\frac{c+ab}{2}]!}{[a]![b]![c]!}.$$
Another quantity, called a $`6j`$ symbol, is derived from the tetrahedral evaluation. Specifically,
(5)
$$\left\{\begin{array}{ccc}a& b& e\\ c& d& f\end{array}\right\}=\frac{\text{Tet}\left(\begin{array}{ccc}a& b& e\\ c& d& f\end{array}\right)(1)^e[e+1]}{\theta (a,d,e)\theta (c,b,e)}.$$
The $`6j`$ symbols can be woven together to give a change of basis matrix for the Whitehead move on graphs. As a consequence they satisfy an orthogonality equation:
(6)
$$\underset{e}{}\left\{\begin{array}{ccc}a& b& e\\ c& d& f\end{array}\right\}\left\{\begin{array}{ccc}d& a& g\\ b& c& e\end{array}\right\}=\delta _f^g,$$
where $`\delta _f^g`$ is the Kronecker delta.
The following proposition seems quite weak, but turns out to be a powerful tool for gauging the convergence of series of Kauffman brackets.
###### Proposition 3.
$$\left|\text{Tet}\left(\begin{array}{ccc}a& b& e\\ c& d& f\end{array}\right)\right|\sqrt{\frac{\theta (b,c,e)\theta (a,d,e)\theta (a,b,f)\theta (c,d,f)}{(1)^{e+f}[e+1][f+1]}}$$
###### Proof.
In order for all the triples at the vertices of a tetrahedron to be admissible , the parity of the sum of the entries in any two columns of
$$\text{Tet}\left(\begin{array}{ccc}a& b& e\\ c& d& f\end{array}\right)$$
has to be the same. Use (5) to expand the formulas for the $`6j`$ symbols in the orthogonality relation (6), with $`g=f`$. The tetrahedral evaluations are equal and the signs of the $`\theta `$’s and the $`(1)^{e+f}`$ cancel so that each term in the sum is positive. Hence every term in the sum is less than $`1`$. Fixing $`e`$ and putting everything except for the tetrahedral evaluations on the right hand side, and taking square roots yields the desired result.∎
###### Corollary 1.
There is a real valued function $`C(k_1,k_2,k_3)`$ so that
(7)
$$\frac{|\text{Tet}\left(\begin{array}{ccc}i& i& i\\ k_1& k_2& k_3\end{array}\right)|}{\sqrt{|\theta (i,i,k_1)\theta (i,i,k_2)\theta (i,i,k_3)}|}$$
is less than $`t^iC(k_1,k_2,k_3)`$ whenever the graphs corresponding to the functions in the formula are admissibly labeled.
###### Proof.
Substitute into the inequality from Proposition 3 to get,
(8)
$$\left|\text{Tet}\left(\begin{array}{ccc}i& i& i\\ k_1& k_2& k_3\end{array}\right)\right|\sqrt{\frac{\theta (k_1,k_2,k_3)\theta (i,i,k_1)\theta (i,i,k_2)\theta (i,i,k_3)}{(1)^{i+k_3}[k_3+1][i+1]}}.$$
Shift $`\sqrt{\theta (i,i,k_1)\theta (i,i,k_2)\theta (i,i,k_3)}`$ to the left hand side. Use the fact that $`\frac{1}{[i+1]}t^{2i}`$ to make the right hand side bigger. Finally, note that the remaining factor on the right hand side is a function of $`k_1`$, $`k_2`$ and $`k_3`$. ∎
###### Theorem 1.
The sequence $`_{i=0}^n(1)^i[i+1]s_i`$ defines a distribution for $`g>1`$. That is, the limit
$$𝒴_D(\alpha )=\underset{n\mathrm{}}{lim}𝒴(\alpha \underset{i=0}{\overset{n}{}}(1)^i[i+1]s_i)$$
exists and gives a well defined trace on $`K_t(\mathrm{\Sigma }_{g,1})`$ when $`g>1`$.
###### Proof.
Choose a trivalent spine for $`\mathrm{\Sigma }_{g,1}`$ with $`4g2`$ vertices and $`6g3`$ edges. Basis elements $`s_c`$ for $`K_t(\mathrm{\Sigma }_{g,1})`$ correspond to labeling the edges admissibly with integers $`k_j`$, where $`j`$ runs from $`1`$ to $`6g3`$. Let $`s_i`$ denote the core of an annulus that runs parallel to the boundary, labeled with the $`i`$th Jones-Wenzl idempotent. In order to compute $`𝒴(s_cs_i)`$ place both skeins in the same diagram. Choose a system of arcs, each intersecting this configuration transversely in three points, that isolate the vertices from one another. The transverse points of intersection are labeled $`i`$, $`k_j`$, $`i`$ as you traverse each arc. Fuse along these arcs, until the resulting graphs intersect each arc in at most one point. Discard any term where the label on an edge intersecting an arc is not zero. Given a vertex $`v`$, let $`(k_{v1},k_{v2},k_{v3})`$ be the triple of colors appearing there. The resulting answer is:
(9)
$$𝒴(s_cs_i)=\underset{j=1}{\overset{6g3}{}}\frac{1}{\theta (i,i,k_j)}\underset{v}{}\text{Tet}\left(\begin{array}{ccc}i& i& i\\ k_{v1}& k_{v2}& k_{v3}\end{array}\right).$$
Each edge appears at exactly two vertices, so (9) can be written as a product of $`4g2`$ factors like (7). By Corollary 1 the absolute value of $`𝒴(s_cs_i)`$ is less than $`C(k_j)t^{i(4g2)}`$, where $`C(k_j)`$ is a number depending only on the $`k_j`$. The $`n`$th partial sum for $`𝒴_D(s_c)`$ is
$$\underset{i=0}{\overset{n}{}}(1)^i[i+1]\underset{j=1}{\overset{6g3}{}}\frac{1}{\theta (i,i,k_j)}\underset{v}{}\text{Tet}\left(\begin{array}{ccc}i& i& i\\ k_{v1}& k_{v2}& k_{v3}\end{array}\right).$$
Note that $`[i+1]`$ is less than $`(i+1)t^{2i}`$. Hence the $`i`$-th summand is less than $`(i+1)(1)^iC(k_j)t^{i(4g4)}`$. The ratio test implies that the sequence of partial sums is absolutely convergent for $`0<t<1`$.
Finally, $`𝒴_𝒟`$ is a trace since the partial sums $`_{i=0}^n(1)^i[i+1]s_i`$ can be seen as lying in the center of $`K_t(\mathrm{\Sigma }_{g,1})`$. ∎
###### Theorem 2.
$`𝒴_D`$ descends to give a well defined trace
$$𝒴:K_t(\mathrm{\Sigma }_g).$$
###### Proof.
There is a homomorphism $`K_t(\mathrm{\Sigma }_{g,1})K_t(\mathrm{\Sigma }_g)`$ induced by inclusion. The surface $`K_t(\mathrm{\Sigma }_g)`$ is the result of adding a disk to the boundary of surface $`K_t(\mathrm{\Sigma }_{g,1})`$. The kernel of this homomorphism consists of all skeins that can be written as a linear combination of handle-slides. The next step is to show that the linear functional $`𝒴_D`$ annihilates all handle-slides. To this end we analyze the difference of the two skeins in the annulus (relative to a pair of points in the boundary).
(10)
$$\underset{i=0}{\overset{n}{}}(1)^i[i+1]\left(\text{}\text{}\right)$$
The analysis of the diagram (10) diagram is due to Lickorish, . It is equal to:
(11)
$$(1)^n[n+1]\left(\text{}\text{}\right).$$
This diagram needs to be set in place. Using standard arguments as in yields that we only need to check handle-slides of the following form. Take a skein corresponding to a colored spine, and separate one strand along an edge.
Now slide the strand over the added disk, locally the diagram looks like:
Multiplying each of the diagrams above by $`_{i=0}^n(1)^i[i+1]s_i`$, taking their difference, and using the identity (10)=(11), we get a difference of two terms like the one below. In the first one the label $`u=n`$ and the label $`v=n+1`$, and in the second one $`u=n+1`$ and $`v=n`$.
Fusing to isolate the vertices of this diagram requires two more cross cuts than the diagrams we have been working with up till now. We get the product of
(12)
$$(1)^n[n+1]\frac{1}{\theta (u,k,u)\theta (u,k1,v)}\text{Tet}\left(\begin{array}{ccc}u& u& v\\ 1& k1& k\end{array}\right)\text{Tet}\left(\begin{array}{ccc}u& v& u\\ 1& k& k1\end{array}\right)$$
with the standard product,
(13)
$$\underset{j=1}{\overset{6g3}{}}\frac{1}{\theta (u,u,k_j)}\underset{v}{}\text{Tet}\left(\begin{array}{ccc}u& u& u\\ k_{v1}& k_{v2}& k_{v3}\end{array}\right).$$
The product (13) is smaller than a global constant, depending on the $`k_j`$, times $`t^{n(4g2)}`$. It remains to ascertain that the term (12) is not too large. Using the inequality from Proposition 3 we get that, regardless of whether $`u=n`$ and $`u=n+1`$, or $`u=n+1`$ and $`u=n`$, the absolute value of (12) is less than $`[n+2]`$, which is a universal constant times $`t^{2n}`$. As long as the genus of the surface is greater than $`1`$, the full product goes to zero as $`n`$ goes to infinity. So, in the limit, all handle-slides are annihilated. ∎
The case of a surface of genus $`1`$ is slightly different. To get a convergent distribution we need to divide the partial sum $`_{i=0}^n(1)^i[i+1]s_i`$ by $`n`$. The sequence is then Cauchy and defines a distribution on $`K_t(T^2)`$.
The algebra $`K_t(T^2)`$ is very nice for working examples. If $`(p,q)`$ is a pair of integers that are relatively prime there is an obvious skein $`s_{(p,q)}`$ in $`K_t(T^2)`$ corresponding to the $`(p,q)`$ curve on the torus . Define a family of skeins based on $`(p,q)`$ by using the following iterative scheme: $`s_{(p,q)_0}=2s_{(0,0)}`$, that is, twice the empty skein, and $`s_{(p,q)_1}=s_{(p,q)}`$. For $`d>1`$ define:
$$s_{(p,q)_d}=s_{(p,q)}s_{(p,q)_{d1}}s_{(p,q)_{d2}}.$$
Finally, if $`d=\text{gcd}\{p,q\}`$, let
$$s_{(p,q)}=s_{(p/d,q/d)_d}.$$
Using this notation the product in $`K_t(T^2)`$ is given by
(14)
$$s_{(p,q)}s_{(u,v)}=t^{\left|\begin{array}{cc}p& q\\ u& v\end{array}\right|}s_{(p+u,q+v)}+t^{\left|\begin{array}{cc}p& q\\ u& v\end{array}\right|}s_{(pu,qv)}.$$
The formula (14) is proven in .
There is a map
$$\mu :K_t(T^2)\mathrm{}H_1(T^2;Z_2)$$
introduced in . Let
$$\mu \left(\underset{(p,q)}{}a_{(p,q)}s_{(p,q)}\right)=a_{(0,0)}\mathrm{}+\underset{(p,q)(0,0)}{}a_{(p,q)}[(p,q)],$$
where $`[(p,q)]`$ is the $`Z_2`$–homology class in $`H_1(T^2;Z_2)`$ corresponding to $`d=\text{gcd}\{p,q\}`$ copies of a $`(p/d,q/d)`$ curve on the torus. The map $`\mu `$ has as its kernel the submodule of all commutators. Hence any linear functional on the five dimensional space that is the image of $`\mu `$ is a trace. It is easy to check that there is a three dimensional family of traces that are invariant under diffeomorphism. In this set up
$$𝒴\left(\underset{(p,q)}{}a_{(p,q)}s_{(p,q)}\right)=a_{(0,0)}.$$
This is the same trace as the one induced from the inclusion of $`K_t(T^2)`$ into the non-commutative torus .
Towards uniqueness of the Yang-Mills measure, it should be normalized, just as the symplectic measure on moduli space needs to be normalized. It should also be invariant under diffeomorphism, and be local. Locality is made up by two rules. One for cutting a surface along an arc and one for removing a point from a closed surface. If we formalize our rules correctly, we get the following:
###### Theorem 3.
The Yang-Mills measure is the unique, local, diffeomorphism invariant trace on $`K_t(\mathrm{\Sigma }_g)`$ up to normalization.∎
## 5. Roots of Unity
Fusion no longer holds in $`K_t(M)`$ when $`t`$ is a root of unity. However, when $`t=e^{\frac{\pi i}{2r}}`$ then one can take a quotient, where an appropriate form of the fusion identity is true. This can be done by setting any skein containing the $`(r1)`$-st Jones-Wenzl idempotent equal to zero. The quotient is denoted $`K_{r,f}(M)`$. The reduced skein is a central object in the construction of quantum invariants of $`3`$-manifolds .
The Yang-Mills measure on a surface with boundary is obtained the same way as for other values of $`t`$. Since $`[r]=0`$, the iterative procedure for finding a skein in the annulus that annihilates handle-slides terminates, to yield
$$\underset{i=0}{\overset{r2}{}}(1)^i[i+1]^i.$$
There is an induced trace,
$$𝒴:K_{r,f}(\mathrm{\Sigma }_g),$$
constructed the same way as for other $`t`$ except that there is no need to take a limit because the formula is a finite sum.
Notice that $`\mathrm{\Sigma }_g`$ is the boundary of a handlebody $`H_g`$ (it doesn’t make any difference which one). There is an action of $`K_{r,f}(\mathrm{\Sigma }_g)`$ on $`K_{r,f}(H_g)`$ given by gluing skeins in $`\mathrm{\Sigma }_g\times I`$ into a collar of the boundary of $`H_g`$. The action gives a map
$$\varphi :K_{r,f}(\mathrm{\Sigma }_g)\mathrm{End}(K_{r,f}(H_g)).$$
As we are working at a root of unity, $`K_{r,f}(H_g)`$ is a finite dimensional vector space. Denote its dimension by $`d`$, and let $`\omega =𝒴(\mathrm{})=_{i=0}^{r2}\frac{1}{[i+1]^{2g2}}`$. The Yang-Mills measure is:
$$𝒴(\alpha )=\frac{\omega }{d}\mathrm{tr}(\varphi (\alpha )).$$
From the map $`\varphi `$ is injective and onto. Hence we can identify $`K_{r,f}(\mathrm{\Sigma }_g)`$ with $`\mathrm{End}(K_{r,f}(H_g))`$. The Yang-Mills measure is zero on commutators. Thus it factors through
$$\mathrm{End}(K_{r,f}(H_g))/[\mathrm{End}(K_{r,f}(H_g)),\mathrm{End}(K_{r,f}(H_g))].$$
This quotient is a $`1`$-dimensional vector space. Hence any two linear functionals that factor through this quotient are equal if they agree on the identity matrix. The trace also vanishes on commutators, thus it factors through the commutator quotient. The normalization in the formula causes the two induced linear functionals to be the same.
Next we address the cases of $`t=\pm 1`$. Since the formula for the measure of a spine is an even function of $`t`$, we only need to consider one value. The value $`t=1`$ is more convenient as the correspondence between $`K_1(F)`$ and the $`SU(2)`$-characters of $`\pi _1(F)`$ is simpler. The skein of the disjoint union of curves $`c_i`$ corresponds to the function that sends the representation $`\rho `$ to
$$\underset{i}{}\mathrm{tr}(\rho (c_i)).$$
###### Theorem 4.
The Yang-Mills measure is well defined on $`K_{\pm 1}(\mathrm{\Sigma }_g)`$ for $`g>1`$. Let $`s_c`$ be an admissibly colored trivalent spine of $`\mathrm{\Sigma }_g`$. If $`t_n`$, with $`|t_n|1`$, is a sequence of complex numbers converging to $`\pm 1`$ then
$$\underset{n\mathrm{}}{lim}𝒴_{t_n}(s_c)=𝒴_{\pm 1}(s_c).$$
###### Proof.
The formulas for working with skeins in $`K_1(F)`$ are the same as the ones for $`|t|1`$ except that quantized integers are replaced by ordinary integers. These formulas are the limits as $`t1`$ of the values we have been using. Revisiting the fundamental estimate (8), we see that,
(15)
$$\frac{|\text{Tet}\left(\begin{array}{ccc}i& i& i\\ k_1& k_2& k_3\end{array}\right)|}{\sqrt{|\theta (i,i,k_1)\theta (i,i,k_2)\theta (i,i,k_3)}|}\sqrt{\frac{\theta (k_1,k_2,k_3)}{(1)^{i+k_3}(k_3+1)(i+1)}}$$
from which we conclude that the right hand side is less than or equal to
$$\frac{C(k_1,k_2,k_3)}{\sqrt{i+1}}.$$
Considering the series for the Yang-Mills measure of a spine, comparison to the p-series implies that it converges as long as the surface has genus greater than $`1`$. Similarly, the Yang-Mills measure is invariant under handle-slides.
The convergence statement follows from the fact that the series that define the Yang-Mills measure at $`t_n`$ converge absolutely, and the terms of the series converge to the terms of the series for the Yang-Mills measure at $`1`$. ∎
For a surface of genus $`1`$ we divide the partial sums, as before, by the number of terms in the sum, and the series then converges.
###### Theorem 5.
The Yang-Mills measure at $`t=1`$ is the symplectic measure on $`(\mathrm{\Sigma }_g)`$.
###### Proof.
Using Weyl orthogonality to compute Witten’s Yang-Mills measure for a surface of area $`\rho `$ yields that its value on the spine $`s_c`$ is given by the series
$$\underset{i=0}{\overset{\mathrm{}}{}}(1)^i(i+1)e^{\rho c_2(i)}\underset{j=1}{\overset{6g3}{}}\frac{1}{\theta (i,i,k_j)}\underset{v}{}\text{Tet}\left(\begin{array}{ccc}i& i& i\\ k_{v1}& k_{v2}& k_{v3}\end{array}\right),$$
where the edges of $`s_c`$ carry colors $`k_i`$, and $`k_{v_i}`$ are the colors of the edges ending at the vertex $`v`$, and $`c_2(i)`$ is the value of the quadratic Casimir operator on the $`(i+1)`$-dimensional irreducible representation of $`SU(2)`$. As both Witten’s series and our series converge absolutely, and Witten’s formula converges term by term to our formula as $`\rho 0`$, the limit of Witten’s Yang-Mills measure is equal to our Yang-Mills measure at $`t=1`$. Finally, Forman showed that the limit as $`\rho 0`$ of Witten’s measure is the symplectic measure on $`(\mathrm{\Sigma }_g)`$, normalized as in . ∎
Suppose now that $`|t|=1`$ and $`t`$ is not a root of unity. Evaluation of the Yang-Mills measure on the empty skein on a surface of genus $`g`$ yields $`_{i=o}^{\mathrm{}}\frac{1}{[i+1]^{2g2}}`$. As $`t`$ is not a root of unity the number $`[i+1]^{2g2}`$ gets arbitrarily close to $`1`$ infinitely often, which means that the series does not converge. Therefore the Yang-Mills measure does not exist away from roots of unity on the unit circle.
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# REFERENCES
Reply to Comment on “Thermal Model for Adaptive Competition in a Market”: In our Letter we introduced a generalization of the minority game (MG) , called the thermal minority game (TMG). One of the main new features was allowing for the stochastic decision making of the agents, controlled by a temperature $`T`$. In the completely deterministic case $`T=0`$, the original MG is recovered.
In their Comment Challet et al. claim that: (i) the equations of our model reduce to Eqs.(2) of ; (ii) Eqs.(2) of lead to the “exact solution of the MG”; (iii) the crossover to a random value of the variance for $`T1`$ found in Figs. 2-3 of is due to short waiting times in the simulations. Remarks (i) and (ii) are incorrect. Point (iii) is true, and highlights even more the crucial role of the temperature in the TMG. The central claim of is that the effects of the temperature in the TMG “can be eliminated by time rescaling” and consequently the behaviour of the TMG is “independent of T”. These statements have no general validity.
Challet et al. obtained their Eqs.(2) of by averaging our equations both over the noise $`\stackrel{}{\eta }`$ and the individual strategy distribution. Averaging over $`\stackrel{}{\eta }`$ is legitimate, since it preserves the full macroscopic dynamics for all $`d`$ and $`T`$. However, averaging over the individual strategies as in is too naive and misses important correlated fluctuations. Besides, this procedure is conceptually wrong: replacing $`R_i^{}`$ by its average is equivalent to allowing the agents (who have $`s`$ fixed strategies available) to play with any strategy formed by a linear combination of the $`s`$ fixed ones. This is not sensible and contrary to the original spirit of the model .
In Fig. 1 we demonstrate the above assertions by comparing the results of simulations on the TMG with those resulting from the equations of Challet et al. Not surprisingly, for $`T0`$ the approximation leading to Eqs.(2) of works well (at $`T=0`$ there is no stochasticity and the average is equal to the best strategy). However, as soon as we turn on $`T`$, this approximation fails to reproduce the correct behaviour. Clearly, Eqs. (2) of miss fluctuation effects and do not describe the behaviour of our system for all $`T0`$.
In Challet et al. approximate the r.h.s. of Eqs.(2) of as the gradient of an effective Hamiltonian $``$ and study the MG by minimizing $``$. This procedure is predicated on the assumption of equilibration of the strategy-use probabilities $`\pi _i^a(t)`$. This assumption is false for $`d<d_c`$ . The consequence is that although their method correctly accounts for the phase $`d>d_c`$, it fails completely for $`d<d_c`$ (see Fig. 1 of ). Thus, to claim as Challet et al. do in (ii), that they have found the “exact solution of the MG” is incorrect and misleading.
Remark (iii) of the Comment is correct: the crossover to a random variance for large $`T`$ observed in is due to finite simulation times. The time required to reach the steady state for $`T1`$ is of order $`NT`$. This is a very interesting observation. It means that in the phase $`d<d_c`$, for any finite temperature larger than a critical value $`T_c1`$ , the performance of the system will be better than the original MG, provided that we wait long enough. In the inset of Fig.1 we show $`\sigma `$ vs. $`T`$ for waiting times much larger than $`NT`$ for the values of $`d`$ of Fig.2 of .
In other words, what Challet et al. have correctly pointed out in their Comment is that any degree of stochasticity above a given threshold makes the TMG perform better than the MG. The claim that “the collective behaviour should be independent of $`T`$” is clearly wrong. The remarkable structure of the TMG with the temperature begs for further investigation.
Andrea Cavagna<sup>1</sup>, Juan P. Garrahan<sup>2</sup>, Irene Giardina<sup>3</sup> and David Sherrington<sup>2</sup>
<sup>1</sup>Department of Physics and Astronomy
The University of Manchester, M13 9PL, UK
<sup>2</sup>Theoretical Physics, University of Oxford
1 Keble Road, Oxford OX1 3NP, UK
<sup>3</sup>Service de Physique Theorique, CEA/Saclay
F-91191 Gif-sur-Yvette Cedex, France
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# A combinatorial property of generic immersions of curves
A divide (called “partage” in \[AC2\]) with $`r`$ branches is the image $`P`$ of a generic relative immersion of the union of $`r,r𝐍,`$ copies of the unit interval $`[0,1]`$ in the unit disk $`D^2`$. A region of a divide $`P`$ is a connected component $`A`$ of $`D^2P`$, that does not meet the boundary $`D^2`$. A segment of $`P`$ is a connected component of the complement of the double points of $`P`$ in $`P`$. A sector of $`P`$ is a connected component of the germ at a double point of a region. We say that a divide is cellular, if $`P`$ is connected and the closure in $`D^2`$ of each region of $`P`$ is contractible. The link $`L(P)S^3`$ of a divide $`P`$ is obtained by a hodographic construction \[AC4-5\]:
$$L(P):=\{(x,u)TD^2xP,uT_xP,x^2+u^2=1\}$$
The link $`L(P)`$ of a divide has a natural orientation. The complement $`S^3L(P)`$ of the link of a connected divide $`P`$ admits a fibration over the circle $`S^1`$, whose restriction to each oriented meridian of each component is a degree $`1`$ map \[AC5\]. The fiber $`F_P`$ is a surface of genus $`\delta (P)`$ with $`r`$ boundary components, where $`\delta (P)`$ is the number of double points of $`P`$. Let $`T_P:F_PF_P`$ be the monodromy diffeomorphism of this fibration. A divide $`P`$ is simple if $`P`$ is connected and has at least one double point, such that there does not exist a relatively embedded copy of $`[0,1]`$ in $`D^2`$, that cuts $`P`$ transversally in one point and separates the double points of $`P`$ in two non-empty sets.
###### Theorem 1
The Lefschetz number of the monodromy $`T_P`$ of a simple, cellular divide is $`0`$.
If the cellular divide $`P`$ is the saddle level of a local real maximal deformation of a real equation of a plane curve singularity \[S1-2,AC2,AC4,G-Z\], the Lefschetz number of the monodromy $`T_P`$ equals the Lefschetz number of the monodromy of the singularity, and therefore equals $`0`$ by \[AC1\]. The divide of Fig. $`1`$ is obtained by a maximal real deformation from the singularity $`x((y^2x^3)^24yx^4x^7)`$ \[G-Z\]. So the Lefschetz number of its monodromy equals $`0`$ although the divide is not cellular.
Fig. 1. Divide for the singularity $`x((y^2x^3)^24yx^4x^7)`$.
The Lefschetz number of the non-cellular divide of Fig. $`2a`$ equals $`2`$. The Lefschetz number of the non-simple divide of Fig. $`2b`$ equals $`1`$.
Fig. 2. a: non-cellular, b: non-simple.
Before giving the proof, we will recall some material about the Seifert form of the link of a connected divide, see \[AC2,AC4, G-Z\]. The connected components of $`D^2P`$ will be signed with $`\pm `$, such that two components have opposite signs if their closures have a segment of $`P`$ in common. Moreover, we choose a basepoint $`b_A`$ in each region $`A`$. The geometric Dynkin diagram $`\mathrm{\Gamma }`$ of the divide $`P`$ is the planar graph in $`D^2`$ obtained as follows. The vertices of $`\mathrm{\Gamma }`$ are the double points of $`P`$ and the basepoints of the regions. The edges of $`\mathrm{\Gamma }`$ are arcs of two species. For each sector at a double point of $`P`$ we draw an arc, that connects the double point with the basepoint of the region to which the sector belongs. Moreover, for each segment of $`P`$, we connect the basepoints of those regions, whose closures meet along the given segment by an arc having one transversal intersection with the segment. Pairs of vertices of $`\mathrm{\Gamma }`$ can be connected by several edges precisely if the divide is non-cellular.
The geometric Dynkin diagram was introduced by Sabir Gusein-Zade \[G-Z\] as system of gradient lines of a morse function on $`D^2`$. We number the vertices of $`\mathrm{\Gamma }`$ with $`1i\mu _P`$ by taking first the basepoints of the $``$ regions, then the double points of $`P`$ and finally the basepoints of the $`+`$ regions. The geometric Dynkin diagram defines the Seifert form $`S`$ of the fibered link $`L(P)`$. On the vector space spanned by the set of vertices $`V`$ of $`\mathrm{\Gamma }`$ we define an upper triangular nilpotent endomorphism $`N:=(n_{ij})`$ by setting the matrix coefficient $`n_{ij}`$ for $`i<j`$ equal to the number of edges of $`\mathrm{\Gamma }`$ that connect $`i`$ and $`j`$.
From the tricollaring of $`\mathrm{\Gamma }`$ by $`\{,.,+\}`$ one can conclude for any divide a basic fact, namely that the equality $`N^3=0`$ holds. A divide is a slalom divide of a rooted tree or of a disk-wide-web if and only if the equality $`N^2=0`$ holds, see \[AC3,AC6\].
The Seifert form is represented by $`Id+N`$. The monodromy action $`T_{}`$ on the first homology of $`F_P`$ is represented by the matrix $`{}_{}{}^{t}(Id+N)_{}^{1}(Id+N)`$. For the Lefschetz number $`\mathrm{\Lambda }_P`$ of $`T`$ this yields the expression
$$1\text{Tr}(Id^tN+N^tNN+^tN^2N)=1\mu _P+\text{Tr}(^tNN)\text{Tr}(^tN^2N)$$
For instance for a slalom divide we have
$$\mathrm{\Lambda }=1\mu _P+\text{Tr}(^tNN)=1\mu _P+(\mu _P1)=0.$$
Proof of theorem $`1`$: For a cellular divide we have $`n_{ij}^2=n_{ij}`$, hence
$$\text{Tr}(^tNN)=\underset{ij}{}n_{ij}^2=\underset{ij}{}n_{ij}=e_P$$
where $`e_P`$ is the number of edges of the geometric Dynkin diagram $`\mathrm{\Gamma }`$. The term $`\text{Tr}(^tN^2N)`$ has for a cellular divide also a combinatorial interpretation. We call flag of the divide $`P`$ a pair $`(a,b)`$ of edges of $`\mathrm{\Gamma }`$ such that the edge $`a`$ connects a basepoint of a $``$ region with a double point of $`P`$, while the edge $`b`$ connects this double point with the basepoint of a $`+`$ region. Calling the number of flags $`f_P`$ we have for a cellular divide $`\text{Tr}(^tN^2N)=f_P`$ and it remains to prove: $`\mu _Pe_P+f_P=1.`$
Fig. $`3`$. Body of a divide with its partial triangulation.
We call body $`B_P`$ of the divide $`P`$ the union of the set of its double points with the closure of the union of its regions. For a simple divide $`P`$ the Euler-Poincaré characteristic $`\chi (B_P)`$ equals $`1`$. Each flag $`(a,b)`$ of $`P`$ defines a triangle $`\mathrm{\Delta }_{(a,b)}=(a,b,ab)B_P`$ in the the body, where $`ab`$ is the edge of $`\mathrm{\Gamma }`$ which connects the non common endpoints of $`a`$ and $`b`$. Let $`SB_P`$ be the union of the set of double points of $`P`$, of the edges of $`\mathrm{\Gamma }`$ and of the triangles associated to the flags. This system of vertices, edges and triangles is a triangulation of $`S`$. The body $`B_P`$ collapses onto $`S`$, see Fig $`3`$. Hence we conclude
$$\mu _Pe_P+f_P=\chi (S)=\chi (B_P)=1,$$
which completes the proof.
I like to thank Marc Burger for his remark, that $`\mu _Pe_P+f_P`$ should be interpreted as Euler-Poincaré characteristic.
Remark: The property $`N^3=0`$ that, as we have seen here above, holds for any divide, indicates that the computation of the traces of the iterates of the monodromy is related to the random walk on the geometric Dynkin diagram $`\mathrm{\Gamma }`$. We ask for an expression of $`\text{Tr}(T^k)`$ in terms of the generating function of the random walk on the Dynkin diagram $`\mathrm{\Gamma }`$, see \[D-J\].
\[AC1\] Norbert A’Campo, Le nombre de Lefschetz d’une monodromie, Nederl. Akad. Wetensch. Proc. Ser. A {76}=Indag. Math. 35 (1973), 113–118,
\[AC2\] Norbert A’Campo, Le Groupe de Monodromie du Déploiement des Singularités Isolées de Courbes Planes I, Math. Ann. 213 (1975) 1–32.
\[AC3\] Norbert A’Campo, Le Groupe de Monodromie du Déploiement des Singularités Isolées de Courbes Planes II, Actes du Congrès International des Mathématiciens, tome 1, 395–404, Vancouver, 1974.
\[AC4\] Norbert A’Campo, Real deformations and complex topology of plane curve singularities, Annales de la Faculté des Sciences de Toulouse 8 (1999), 1, 5–23.
\[AC5\] Norbert A’Campo, Generic immersions of curves, knots, monodromy and gordian number, Publ. Math. I.H.E.S. 88 (1998), 151-169, (1999).
\[AC6\] Norbert A’Campo, Planar trees, slalom curves and hyperbolic knots, Publ. Math. I.H.E.S. 88 (1998), 171-180, (1999).
\[D-J\] E.B. Dynkin, A.A. Juschkewitsch, Saetze und Aufgaben ueber Markoffsche Prozesse, aus dem Russischen uebersetzt von Klaus Schürger: Teoremy i zadachi o prochessakh Markova, Moskva 1967, Berlin : Springer, 1969 Heidelberger Taschenb”ucher.
\[G-Z\] S. M. Gusein-Zade, Matrices d’intersections pour certaines singularités de fonctions de 2 variables, Funkcional. Anal. i Prilozen 8, (1974), 11–15.
\[S1\] Charlotte Angas Scott, On the Higher Singularities of Plane Curves, Amer. J. Math. 14,(1892) 301–325.
\[S2\] Charlotte Angas Scott, The Nature and Effect of Singularities of Plane Algebraic Curves, Amer. J. Math. 15, (1893) 221–243.
Universitaet Basel, Rheinsprung 21, CH-4051 Basel.
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# Untitled Document
hep-th/0005006 IASSNS-HEP-00/36 TIFR/TH/00-20
Noncommutative Tachyons
Keshav Dasgupta<sup>a,</sup><sup>1</sup> keshav@sns.ias.edu, Sunil Mukhi<sup>b,</sup><sup>2</sup> mukhi@tifr.res.in and Govindan Rajesh<sup>a,</sup><sup>3</sup> rajesh@sns.ias.edu
<sup>a</sup>School of Natural Sciences
Institute for Advanced Study
Princeton, NJ 08540, USA
<sup>b</sup>Tata Institute of Fundamental Research
Homi Bhabha Road, Mumbai 400 005, India
When unstable Dp-branes in type II string theory are placed in a B-field, the resulting tachyonic world-volume theory becomes noncommutative. We argue that for large noncommutativity parameter, condensation of the tachyon as a noncommutative soliton leads to new decay modes of the Dp-brane into (p-2)-brane configurations, which we interpret as suitably smeared BPS D(p-1)-branes. Some of these configurations are metastable. We discuss various generalizations of this decay process.
4/00
1. Introduction
The decay modes of unstable D-branes and brane-antibrane pairs have been extensively studied in the last couple of years (for reviews see for example Refs.\[1,,2,,3\]). Some of the most basic decay modes found so far are: annihilation into vacuum\[4,,5,,6,,7\], annihilation via kink condensation into a brane of codimension one\[8,,9,,10,,11\], and annihilation via vortex condensation to a brane of codimension two\[9,,12\]. Condensation of higher-codimension topological solitons has also been studied\[12,,11\]. Some of these decay modes correspond to stable solitons, and in this case the end-products are stable branes, while in other cases the decay modes correspond to unstable solitons and lead to unstable branes. More exotic decay modes are possible when unstable branes or brane-antibrane pairs are either wrapped on homology cycles of nontrivial manifolds\[13,,14,,15\] or suspended between other branes in brane constructions.
In this note, we identify a class of novel processes in which unstable D-branes decay via condensation of a noncommutative soliton. Such solitons are finite-energy classical solutions in noncommutative scalar field theories with large noncommutativity parameter $`\theta _{ij}`$. Physically, this situation can be realised by turning on a suitable constant NS-NS $`B`$-field along the brane world-volume.
Since noncommutative solitons carry no topological charge and are sometimes metastable, the same will be true of the decay products of an unstable brane when such a soliton condenses on its worldvolume. We will interpret these decay products by examining what Ramond-Ramond charge they carry locally.
While this work was nearing completion, we became aware of a forthcoming paper which addresses related questions, mainly in the context of the bosonic string.
2. Noncommutative Tachyons and Brane Decay
Consider an infinitely extended, unstable D2-brane in type IIB string theory. We place it in a $`B_{NS,NS}`$ field with a nonzero component only along the brane world-volume directions, namely $`B_{\mathrm{1\hspace{0.17em}2}}`$. Consider the brane world-volume theory. The light modes are a real tachyon, a massless U(1) gauge field and some massless fermions.
Let us set all fields on the brane, except the tachyon, to 0. The spacetime backgrounds are also as simple as possible: a flat metric and the constant $`B_{\mathrm{1\hspace{0.17em}2}}`$ field. Following arguments in Ref., the effect of the $`B`$-field can be described by making the tachyonic field theory on the D2-brane world-volume into a noncommutative field theory with $``$ product given by
$$f(x)g(x)=e^{\frac{i}{2}\theta (_1_2^{}_2_1^{})}f(x)g(x^{})|_{x=x^{}}$$
where the noncommutativity parameter $`\theta _{ij}`$, given in terms of $`B`$ by
$$\theta _{ij}=(2\pi \alpha ^{})^2\left(\frac{1}{g+2\pi \alpha ^{}B}B\frac{1}{g2\pi \alpha ^{}B}\right)_{ij}$$
has the single component $`\theta _{\mathrm{1\hspace{0.17em}2}}=\theta `$.
Let $`V(T)`$ denote the universal tachyon potential on an unstable D-brane. The action for such a non-BPS D-brane is given by \[21,,22,,23\]:
$$\mu _p^{(u)}d^{p+1}\sigma V(T)\sqrt{det(g_{\mu \nu }+2\pi \alpha ^{}_{\mu \nu }+2\pi \alpha ^{}_\mu T_\nu T)}$$
where $`g_{\mu \nu }`$ is the closed string metric and $`_{\mu \nu }`$ is the linear combination of $`F_{\mu \nu }`$ and $`B_{\mu \nu }`$. We claim that in the presence of a B-field, the tachyonic part of the D2-brane worldvolume action is of the form
$$S=d^2xdt(^\mu T_\mu TV(T))$$
where the notation $`V(T)`$ represents the universal tachyon potential with all products replaced by $``$ products.
To argue this, note that the replacement of $`V(T)`$ by $`V(T)`$ apparently violates universality of the tachyon potential. However, precisely because this universality holds only in the zero-momentum sector, we can expect a violation of universality that vanishes at zero-momentum. This is true of the $``$ product, which reduces to the ordinary product on constant functions. Thus, noncommutativity and universality of the tachyon potential together amount to a derivation of the above action<sup>1</sup> We are grateful to K. Hori for discussions on this point..
For unstable D-branes in type II string theory, $`V(T)`$ is known to be an even function of $`T`$ and is believed to have a double-well shape. Let it take its minima at $`T=\pm T_0`$. At the minimum, the negative potential energy localized on the brane should cancel the unstable D2-brane tension $`\mu _2^{(u)}`$, consistent with the decay of the D2-brane into vacuum. This gives rise to the beautiful equation:
$$\mu _2^{(u)}+V(T=T_0)=0$$
Besides decay into the vacuum, the simplest decay mode of the unstable D2-brane is via kink condensation in this double-well potential, leading to a stable D-string. The D-string charge arises from the Chern-Simons coupling\[11,,21,,24\]:
$$\frac{1}{2T_0}𝑑TB_{RR}$$
where $`B_{RR}`$ is the Ramond-Ramond 2-form in type IIB. For a kink solution along say $`x^2`$, we have $`𝑑x^2_2T=2T_0`$ and we are left with unit coupling to $`B_{RR}`$ and hence unit D-string charge. An anti-kink would produce an anti-D-string.
We work in the limit of large $`\theta `$. It is possible to achieve this limit, while keeping the open string metric fixed, by a suitable choice of scaling for $`\alpha ^{}`$, $`B`$, and the closed string metric $`g`$, for example, the limit
$$\alpha ^{}\sqrt{ϵ},B\sqrt{ϵ},gϵ^2$$
with $`ϵ0`$. Following Ref., it is convenient to scale the coordinates by $`x^i\sqrt{\theta }x^i`$, and one ends up with the action
$$S=d^2xdt(^\mu T_\mu T\theta V(T))$$
where now the $``$ product is as in Eqn.(2.1) but with $`\theta =1`$.
At infinite $`\theta `$, we can look for classical solutions by just solving $`V(T)/T=0`$. For a fairly arbitrary tachyon potential (but in particular, without a linear term), an infinite set of non-constant solitonic solutions of this equation was found in Ref.. In particular, one solution is provided by starting with the function $`\varphi _0(x)2e^{r^2}`$ and writing
$$T(x)=\lambda _i\varphi _0(x)$$
where $`\lambda _i`$ is any nonzero minimum of the ordinary potential. This works by virtue of the fact that (with $`\theta =1`$), we have $`\varphi _0\varphi _0=\varphi _0`$.
Note that this soliton decays to zero at infinity. Although this was not explicitly stated, the analysis and examples of Ref. dealt entirely with potentials having a quadratic minimum at the origin, hence the soliton does in fact tend asymptotically to the vacuum. Our case is just the opposite: a tachyonic potential has a quadratic maximum at the origin. One can of course go from one to the other by shifting $`T(x)`$ by a specific constant, though general shifts of $`T(x)`$ are not allowed — since they introduce a linear term in $`V(T)`$, which in turn invalidates the solutions of Ref.. As we will now see, a small modification of the example above by such a constant shift will give rise to an interesting physical phenomenon with unstable branes. Suppose we first take the classical solution:
$$T(x)=T_0(1\varphi _0(x))$$
Note that, like $`\varphi _0(x)`$, $`(1\varphi _0(x))`$ also squares to itself under the $``$ product. However, it varies from $`+1`$ at spatial infinity to $`1`$ at the origin. Thus the configuration $`T(x)`$ above is a soliton that interpolates from $`T_0`$ at spatial infinity to $`T_0`$ at the origin. One can of course take the negative of this solution and it interpolates in the opposite way.
A key insight into the physical meaning of this process can be obtained by looking at the total energy of this soliton. The energy density of the soliton is easily evaluated using $`(1\varphi _0)(1\varphi _0)=(1\varphi _0)`$:
$$V\left(T_0(1\varphi _0(x))\right)=\left(1\varphi _0(x)\right)V(T_0)$$
Hence the total energy density on the D2-brane, with the noncommutative soliton excited, is
$$\theta (\mu _2^{(u)}+\left(1\varphi _0(x)\right)V(T_0))$$
Using Sen’s conjecture for unstable brane decay, Eqn.(2.1), this works out to be $`\theta \varphi _0(x)\mu _2^{(u)}`$. This has to be interpreted as the energy of the decay product. We see that this decay product is an object whose energy is localised very close to the origin, so it can be considered an exotic 0-brane. It is planar and its energy distribution falls off exponentially with the distance away from the origin. Indeed, it is a D0-brane in a sense, since a string that was ending on the original D2-brane can continue to end on regions where there is finite D2-brane tension. Hence we interpret the noncommutative soliton above as describing the decay of a planar unstable D2-brane into a localised configuration at the origin. We will argue in the next section that this configuration is really a smeared D-string.
It is noteworthy that the energy of the noncommutative soliton can be calculated exactly without a detailed knowledge of the tachyon potential. This is due to the fact, pointed out in Ref., that the soliton solution depends on no details of the potential except its value at the minimum. Remarkably, at the present stage of understanding of string theory, this is the only thing about the tachyon potential on D-branes that we do know reliably, thanks to Eqn.(2.1).
It is clearly important to analyse various properties of this D0-brane including its stability (in the presence of large noncommutativity). The above solution will turn out to be classically unstable, as we will argue in Section 4.
There is, in fact, another classical solution which is most interesting from the point of view of stability. This is given by:
$$T(r)=T_0(12\varphi _0(r))$$
That these are classical solutions follows from the fact that they can equivalently be written
$$T(r)=T_0(1\varphi _0(r))+(T_0)\varphi _0(r)$$
in which form it is evident that they are the superposition<sup>2</sup> Normally, one would not expect superpositions of solutions to be solutions in a non-linear theory. However, the orthogonality properties of the $`\varphi _n`$ under the star product allow us to superpose orthogonal solutions. of two noncommutative solitons, one of which asymptotes to 0 and the other to $`T_0`$. Another way of seeing that these are solutions is that although the function $`(12\varphi _0(r))`$ does not square to itself under the $``$ product (rather it squares to 1), its $`\mathrm{𝚘𝚍𝚍}`$ powers are all equal to itself, and that is sufficient because the tachyon potential is even.
At the origin this solution tends to $`3T_0`$. We will see later that this solution is metastable. The energy of this solution is
$$\begin{array}{cc}\hfill V(T_0(12\varphi _0(r)))& =(1\varphi _0(r))V(T_0)+\varphi _0(r)V(T_0)\hfill \\ & =V(T_0)\hfill \end{array}$$
This is degenerate with the energy of the vacuum solution. Thus, in this decay mode the original D2-brane has decayed into a nontrivial configuration that has the same energy as the vacuum.
Before turning to the issue of stability, we look at some generalizations. A complete set of radially symmetric solutions to the equation $`\varphi \varphi =\varphi `$ was worked out in . They satisfy the relation:
$$\begin{array}{cc}\hfill \varphi _n(r)\varphi _m(r)& =\delta _{mn}\varphi _n(r)\hfill \\ \hfill \underset{n=0}{\overset{\mathrm{}}{}}\varphi _n(r)& =1\hfill \end{array}$$
These functions may be found as follows (the following constitutes a quick alternate derivation to the one in Ref.). Define the generating function
$$\varphi (r,z)=\underset{n=0}{\overset{\mathrm{}}{}}\varphi _n(r)z^n$$
in terms of which the above equations amount to
$$\begin{array}{cc}\hfill \varphi (r,z)\varphi (r,z^{})& =\varphi (r,zz^{})\hfill \\ \hfill \varphi (r,z=1)=1& \end{array}$$
The first of these equations may be solved by going to Fourier space in $`x`$ (recall that above, $`r=|x|`$). One finds that, with a radially symmetric ansatz, the solution is Gaussian in Fourier space and hence also in coordinate space. Inserting this ansatz and solving the above equations, one ends up with the generating function:
$$\varphi (r,z)=\frac{2}{1+z}e^{\frac{z1}{z+1}r^2}$$
Expanding in powers of $`z`$ one finds:
$$\varphi _n(r)=2(1)^ne^{r^2}L_n(2r^2)$$
where $`L_n`$ are the Laguerre polynomials, in agreement with the solution in Ref..
We can now write the obvious generalization of Eqn.(2.1), $`T_n(r)=T_0(1\varphi _n(r))`$, noting that $`(1\varphi _n(r))`$ squares to itself under the star product. This soliton is also a D0-brane, but is $`\sqrt{n}`$ times larger than the original one. While for even $`n`$ the solution interpolates between one vacuum $`T=T_0`$ at infinity to the other vacuum $`T=T_0`$ at the origin, for odd $`n`$ the situation is quite different. From $`\varphi _n(r=0)=2(1)^n`$, we see that the solutions for odd $`n`$ interpolate between the vacuum $`T=T_0`$ at infinity and a non-vacuum configuration $`T=3T_0`$ at the origin. There is, of course, no requirement that the soliton should go to a vacuum at the origin.
It is also easy to see that $`T(r)=T_0(1f(r))`$, with
$$f(r)=\underset{iI}{}\varphi _i(r)$$
where the sum runs over a finite index set $`I`$ of the $`\varphi _i`$’s, each counted exactly once, is also a classical solution, since such a function $`f(r)`$ also squares to itself under the star product. Functions of the form $`(1f(r))`$ with $`f(r)`$ as above are, in fact, the most general functions that square to themselves and have the asymptotic behaviour of a solitonic solution for a tachyonic potential, but as we will see in a moment, these are not the most general extrema of the noncommutative tachyon potential.
It is interesting to note the relation
$$(1\varphi _n)=\underset{mn}{}\varphi _m$$
so that the above solitons may be regarded as the superposition of an infinite set of the noncommutative solitons of . Note that since $`\varphi _n`$ vanishes at infinity, the asymptotics of the functions $`(1\varphi _n)`$ and $`\varphi _n`$ are very different. There is, however, no contradiction with (2.1) because the RHS is an infinite sum.
It is also easy to write down a generalization of Eqn.(2.1). Given that the tachyon potential is an even function, we can find a class of solutions with zero total energy using functions $`f(r)`$ that satisfy $`f(r)f(r)=1`$. Using any such function we can define a solution $`T(r)=T_0f(r)`$ that has $`V(T(r))=V(T_0)`$ and hence vanishing total energy.
Expanding such functions over the $`\varphi _n(r)`$ we have:
$$f(r)=(1+\underset{iI}{}\lambda _i\varphi _i(r))$$
where again $`I`$ is a finite index set of distinct elements. One sees that
$$f(r)f(r)=1+\underset{iI}{}\lambda _i(\lambda _i+2)\varphi _i(r)$$
from which it is clear that each $`\lambda `$ must be equal to 0 or $`2`$. Thus the classical solutions
$$T(r)=T_0(12\underset{iI}{}\varphi _i(r))$$
all have zero energy and will turn out to be metastable.
Finally, we can superpose the two kinds of solutions above to obtain the most general solitonic solution (upto an overall choice of sign, of course) for a double well potential
$$T_r=T_0(1\underset{iI}{}\varphi _i)+(T_0)(\underset{jJ}{}\varphi _j)$$
where $`I`$ and $`J`$ are finite index sets with $`JI`$. This has positive energy when $`J`$ is a proper subset of $`I`$, reducing to (2.1) when $`J`$ is the empty set. When $`J=I`$, on the other hand it reduces to (2.1), and has zero total energy.
More precisely, the total energy of such solitons does not exactly vanish, but for any finite $`\theta `$, it is suppressed relative to normal D-brane energies by a factor of $`\frac{1}{\theta }`$. At finite $`\theta `$, kinetic terms would contribute a small mass correction, as a result of which the energy of these solitons (following analogous arguments in Ref.) is generically $`V(T_0)+𝒪(\frac{1}{\theta })`$. When we rescale coordinates to absorb a $`\sqrt{\theta }`$ in them, the total energy has a factor $`\theta `$ multiplying it, as in Eqn.(2.1). In these coordinates the mass of the above solitons is therefore finite, though the tensions of the standard stable and unstable D-branes of the theory are proportional to $`\theta `$. If we revert to the original coordinates then the standard D-branes have finite tension and the solitons described above have total energies of order $`\frac{1}{\theta }`$.
The extension of the results in this section to unstable Dp-branes for various $`p>2`$ (odd in type IIA and even in type IIB) is straightforward, for the case of noncommutativity only along two spatial directions. More interesting extensions to higher branes will be discussed in a subsequent section.
3. Noncommutative Soliton as a Smeared D-string
In this section we consider the D-brane charge associated with the above solitonic branes. Although the noncommutative solitons are not topological, and hence the exotic 0-branes of the previous section cannot carry a global RR charge, we can gain some insight into their nature by looking at what RR charge they carry locally.
As described for one special case in Eqn.(2.1) above, there is a Chern-Simons coupling on the non-BPS $`p`$-branes\[11,,21,,24\]
$$\frac{1}{2T_0}𝑑TC_p$$
where $`C_p`$ is the Ramond-Ramond $`p`$-form, and $`T`$ is the tachyon. In the limit when the noncommutativity parameter $`\theta \mathrm{}`$, we need only keep terms in the action which scale like $`\theta `$ (after coordinate redefinition). However, the Chern-Simons terms above, being topological, do not pick up any such factor, so in the $`\theta \mathrm{}`$ limit, they might appear to be irrelevant. This, however, is not the case. When $`T`$ is a kink along one of the world-volume directions, we expect the solution to correspond to a $`(p1)`$-brane, and therefore there must exist a coupling $`C_p`$ on its world-volume. The only source for such a coupling is the Chern-Simons coupling above, which can therefore not be neglected, even in the limit $`\theta \mathrm{}`$.
Now consider, as before, the unstable Type IIB D2-brane in a background $`B_{NS}`$ field, and let us consider the tachyon background $`T=T_0(1f)`$, where $`T_0`$ is the minimum of the tachyon potential, and $`f`$ is, say, the Gaussian solution $`\varphi _0`$. Since $`T_0`$ is constant, (3.1) yields
$$\frac{1}{2T_0}T_0𝑑\varphi _0B=\frac{1}{2}\varphi _0^{}(r)B_{t\omega }𝑑rdtd\omega $$
where $`B`$ is the Ramond-Ramond $`2`$-form, so that the soliton locally carries D-string charge. In fact, since, $`\varphi _0^{}(r)𝑑r=\varphi _0(0)=2`$, we get precisely $`B_{t\omega }𝑑td\omega `$, implying that the soliton carries unit D-string charge, and can in fact be interpreted as a D-string winding along the angular coordinate<sup>3</sup> We will use $`\omega `$ to refer to the angular variable of the polar coordinates, reserving the more conventional $`\theta `$ for the noncommutativity parameter. $`\omega `$. The function $`f(r)`$ can be interpreted as the distribution function for the D-string, so that our soliton describes a D-string smeared over the radial coordinate $`r`$, winding around the origin.
Such a D-string would not normally be stable, but would collapse by contracting to a point at the origin. This is equivalent to saying there is no global RR charge carried by this configuration. However, the above argument indicates that there is a local RR charge and enables an interpretation of the exotic decay product in terms of D-strings. To the extent that noncommutativity stabilizes some of the exotic 0-branes, which is discussed in a subsequent section, it will be equally true that such a smeared, winding D-string is stabilized.
Note that if we choose $`T=T_0(1\varphi _0)`$, we get an anti-D-string instead. It is also easy to verify that, if we choose $`f=\varphi _n`$, the soliton carries D-string charge $`(1)^n`$. Again, if $`f=\varphi _0+\varphi _1`$ (so that $`T=T_0(1f)`$ is still a classical solution) the total D-string charge vanishes. However, because the two functions $`\varphi _0`$ and $`\varphi _1`$ do not add to zero, we should really interpret this soliton as a dipole of D-strings smeared over the $`(r,\omega )`$ plane. This argument extends to the situation when we have $`f=\varphi _i`$ for some finite set of $`\varphi _i`$’s.
Finally, for the solutions that are degenerate with the vacuum, for example $`T(r)=T_0(12\varphi _n(r))`$, we find 2 units of D-string or anti D-string charge locally, depending on whether $`n`$ is even or odd. For $`T(r)=T_0(12\varphi _0(r)2\varphi _1(r))`$, we again obtain a dipole of D-strings. The extension to solitons of the form (2.1) and (2.1) is obvious.
4. Stability of the Decay Products
Among the various classical solutions to the noncommutative tachyon theory, we found the class of solutions $`T(r)=T_0(12_{iI}\varphi _i(r))`$ with energy $`V(T_0)`$. This energy is negative, and is in fact just the same as we would have obtained by choosing the trivial classical solution (valid in both commutative and noncommutative cases) $`T(r)=T_0`$. Adding it to the tension $`\mu _2^{(u)}`$ of the unstable D2-brane, we get a total energy 0. Thus (in the limit of infinite noncommutativity) we find that an unstable D2-brane can decay into a configuration that, while quite different from the vacuum<sup>4</sup> The key difference from the vacuum is that this configuration carries local D-string charge, as described in the previous section., is nevertheless degenerate with it. In this situation one would expect the decay product to be stable.
One should, however, think of this as an approximation. As discussed in Section 2, there will be corrections due to finiteness of the noncommutativity parameter $`\theta `$. These corrections will tend to raise the energy of the decay product slightly above the vacuum, in which case it will be metastable. More information about the shape of the tachyon potential in string theory than is presently available would be needed to estimate the lifetime of this state.
The other solutions we described, such as the solution $`T(r)=T_0(1\varphi _n(r))`$, are instead classically unstable. The easiest way to see this is to rewrite the total energy as follows<sup>5</sup> We wish to thank Rajesh Gopakumar for the following argument.
$$V(T)+\mu _2^{(u)}=\stackrel{~}{V}(T^{})$$
with $`T^{}=TT_0`$. Then, $`\stackrel{~}{V}`$ has degenerate minima at $`0`$ and $`2T_0`$, and a maximum at $`T_0`$, and belongs to the general class of potentials discussed in . From this, it is easy to see that the solution $`T(r)=T_0(1\varphi _n(r))`$ corresponds to $`T^{}(r)=T_0\varphi _n(r)`$, which is unstable, as its energy can be decreased by scaling by a constant near unity. In other words, the energy of $`T(r)=T_0(1\varphi _n(r))+ϵ\varphi _n(r)`$, (for small $`ϵ`$) is lower than $`T(r)`$ by $`𝒪(ϵ^2)`$.
These solutions have a finite energy above the vacuum and will decay into it classically. Hence they should be thought of as describing the decay of an unstable D2-brane into a kind of unstable D0-brane. The situation is somewhat reminiscent of the case of a brane-antibrane pair, which can decay into an unstable brane of one lower dimension by condensing an unstable tachyonic kink.
It is fascinating that, from the discussion in the previous section, the classically stable solitons all carry even D-string charge, while all the solitons with odd D-string charge are unstable. Note of course, that not all even charge solitons are stable. For example, the dipole $`T_0(1\varphi _0\varphi _1)`$, and the charge 2 soliton $`T_0(1\varphi _0\varphi _2)`$ are classically unstable.
5. Noncommutative Decays of Higher $`p`$-branes
We can generalize the argument of Section 2 to higher brane solitons. Consider an unstable D3 brane of Type IIA, with worldvolume along $`(x^0,x^1,x^2,x^3)`$ in the presence of a constant background Neveu-Schwarz B-field along $`(x^1,x^2)`$ directions. We expect the geometry along the $`(x^1,x^2)`$ plane to become noncommutative, while $`(x^0,x^3)`$ remain commutative. We can now take the limit of large noncommutativity parameter $`\theta =\theta _{\mathrm{1\hspace{0.17em}2}}`$, and perform the appropriate rescalings of the coordinates. Now consider the tachyon background $`T=K(x^3)(1f(x^1,x^2))`$, where $`K(x^3)`$ is the kink along the commutative direction $`x_3`$, and $`f`$ is one of the noncommutative solitons $`\varphi _n`$. Then from the coupling
$$\frac{1}{2T_0}𝑑TC$$
where $`C`$ is the Ramond-Ramond 3-form, we get
$$\frac{1}{2T_0}_{D3}((1f)dKKdf)C=_{D2}(1f)C\frac{1}{2T_0}_{D3}K𝑑fC$$
where the first term on the RHS is an integral over the D2-brane obtained from the kink $`K`$. In the absence of $`f`$, we should get therefore get a D2 brane. When $`f=\varphi _n`$, however, from the term $`_{D2}C(1f)`$ it appears that the D2 brane charge is altered by an “irrational” amount, since $`_{D2}f(r)𝑑x^1𝑑x^2=2\pi `$, leading to a apparent contradiction. However, we should really examine the charge density, that is, the charge per unit D2-brane volume, and recall that the volume of the D2 brane is infinite. Thus the 1 in $`(1f)`$ still contributes 1 unit of D-brane charge but since $`_{D2}f(r)𝑑x^1𝑑x^2`$ is finite, it does not affect the charge density, so it can be ignored.
We also have a second term in (5.1) :
$$\frac{1}{2T_0}K𝑑fC$$
As in the D2 brane case above, this gives
$$\frac{1}{T_0}K(x^3)C_{0\omega 3}𝑑x^0d\omega dx^3$$
Since the kink $`K`$ goes from $`T_0`$ to $`T_0`$ as $`x^3`$ goes from $`\mathrm{}`$ to $`\mathrm{}`$, we see that we get zero total D2 brane charge. However, we really have a dipole of D2 branes, with the anti-brane at $`x^3<0`$ (and the radial coordinate $`r=\sqrt{(x^1)^2+(x^2)^2}`$ ) and the brane at $`x^3>0`$ (and $`r`$).
Finally, from the coupling in the non-BPS D3 brane action:
$$\frac{1}{2T_0}𝑑TAB_{NS}$$
where $`A`$ is the Ramond-Ramond 1-form, and $`B_{NS}`$ is the background Neveu-Schwarz B-field, we might expect the solitons to carry D0-brane charge. However, since the Neveu-Schwarz background is not quantized, we cannot in general expect integer (or even rational) D-brane charge from (5.1). We expect the resolution to this puzzle to be along the lines of , with the D-brane charge from (5.1) being cancelled by bulk terms.
The above exercise can be repeated with the metastable solutions of the type $`K(12\varphi _n(r))`$ along the noncommutative directions, with similar conclusions.
6. Conclusions
It is pleasing that the field-theoretic study of noncommutative solitons initiated in Ref. has such an elegant application to the scalar field theory of tachyons on the world-volume of unstable D-branes in superstring theory. A more detailed understanding of the physical significance of this decay process, of the stability of the end products, and of modifications due to finite rather than infinite $`\theta `$ is clearly desirable.
In particular, consider the smeared D-string configurations that describe metastable decay products. It would be desirable to explain, from the spacetime (as opposed to brane worldvolume) point of view, why such configurations are rendered metastable by a suitable constant B-field<sup>6</sup> Steve Gubser has suggested that these configurations may be spinning D-strings stabilised by their centrifugal force..
It would be interesting to extend our results to Dp-branes of other dimensions, and to the noncommutative gauge theory that also resides on unstable D-branes. One could also consider the complex tachyon on a brane-antibrane pair, which couples to a worldvolume gauge field, unlike the real noncommutative tachyon studied in this paper. In such a case, one might combine noncommutative solitons with orthogonal tachyonic vortices in the noncommutative scalar-gauge theory that resides on brane-antibrane pairs. Finally, the physical effects of noncommutativity along, rather than transverse to, a topologically stable kink or vortex remain to be explored.
Acknowledgements:
We would like to thank Ori Ganor, Rajesh Gopakumar, Steve Gubser, Jeff Harvey, Joydeep Majumdar, Shiraz Minwalla, Nemani Suryanarayana and Sandip Trivedi for helpful discussions. We are particularly grateful to Rajesh Gopakumar for reading and helping us correct a preliminary version of the manuscript. S.M. would like to thank Nati Seiberg and the Institute for Advanced Study, Princeton, where this work was initiated, for hospitality. The work of K.D. was supported by DOE grant No. DE-FG02-90ER40542, and that of G.R. by NSF grant No. NSF-DMS-9627351.
References
relax A. Sen, ‘‘Non-BPS States and Branes in String Theory’’, hep-th/9904207. relax A. Lerda and R. Russo, ‘‘Stable Non-BPS States in String Theory: A Pedagogical Review’’, hep-th/9905006. relax O. Bergman and M. Gaberdiel, ‘‘Non-BPS Dirichlet Branes’’, hep-th/9908126. relax T. Banks and L. Susskind, ‘‘Brane--Antibrane Forces’’, hep-th/9511194. relax M. Srednicki, ‘‘IIB or Not IIB’’, hep-th/9807138; JHEP 08 (1998) 005. relax A. Sen, ‘‘Stable Non-BPS Bound States of BPS D-branes’’, hep-th/9805019; JHEP 08 (1998) 010. relax A. Sen, ‘‘Tachyon Condensation on the Brane Anti-Brane System’’, hep-th/9805170; JHEP 08 (1998) 012. relax O. Bergman and M. Gaberdiel, ‘‘Stable Non-BPS D-Particle’’, hep-th/9806155; Phys. Lett. B441 (1998) 133. relax A. Sen, ‘‘SO(32) Spinors of Type I and other Brane- Antibrane Pair’’, hep-th/9808141; JHEP 09 (1998) 023. relax A. Sen, ‘‘Type I D-Particle and its Interactions’’, hep-th/9809111; JHEP 10 (1998) 021. relax P. Horava, ‘‘Type IIA D-Branes, K-Theory and Matrix Theory’’, hep-th/9812135; Adv. Theor. Math. Phys. 2 (1999) 1373. relax E. Witten, ‘‘D-branes and K-theory’’, hep-th/9810188; JHEP 12 (1998) 019. relax A. Sen, ‘‘BPS D-Branes on Nonsupersymmetric Cycles’’, hep-th/9812031; JHEP 12 (1998) 021. relax J. Majumder and A. Sen, ‘‘ ‘Blowing Up’ D-Branes on Nonsupersymmetric Cycles’’, hep-th/9906109; JHEP 09 (1999) 004. relax M. Mihailescu, K. Oh and R. Tatar, ‘‘Non-BPS Branes on a Calabi-Yau Threefold and Bose-Fermi Degeneracy’’, hep-th/9910249. relax S. Mukhi, N. V. Suryanarayana and D. Tong, ‘‘Brane- Antibrane Constructions’’, hep-th/0001066; JHEP 03 (2000) 015. relax R. Gopakumar, S. Minwalla and A. Strominger, ‘‘Noncommutative Solitons’’, hep-th/0003160. relax J. Harvey, P. Kraus, F. Larsen, E. Martinec, to appear. relax N. Seiberg and E. Witten, ‘‘String Theory and Noncommutative Geometry’’, hep-th/9908142; JHEP 09 (1999) 032. relax A. Sen, ‘‘Universality of the Tachyon Potential’’, hep-th/9911116. relax A. Sen, ‘‘Supersymmetric World-Volume Action For Non-BPS D-Branes’’, hep-th/9909062; JHEP 10 (1999) 008. relax E. A. Bergshoeff, M. de Roo, T. C. de Wit, E. Eyras and S. Panda, ‘‘T-Duality and Actions for Non-BPS D-Branes’’, hep-th/0003221. relax M. R. Garousi, ‘‘Tachyon Coupling on Non-BPS D-Branes and Dirac-Born-Infeld Action’’, hep-th/0003122. relax M. Billò, B. Craps and F. Roose, ‘‘Ramond-Ramond Coupling of Non-BPS D-Branes’’, hep-th/9905157; JHEP 06 (1999) 033. relax W. Taylor, ‘‘D2 Branes in B fields’’, hep-th/0004141.
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# Magnetic-Field-Driven Superconductor-Insulator-Type Transition in Graphite
## Abstract
A magnetic-field-driven transition from metallic- to semiconducting-type behavior in the basal-plane resistance takes place in highly oriented pyrolytic graphite at a field $`H_c1`$kOe applied along the hexagonal c-axis. The analysis of the data reveals a striking similarity between this transition and that measured in thin-film superconductors and Si MOSFET’s. However, in contrast to those materials, the transition in graphite is observable at almost two orders of magnitude higher temperatures.
, , , , , and thanks: On leave from Instituto de Fisica, Unicamp, 13083-970 Campinas, Sao Paulo, Brasil.
The anomalous carrier density, temperature, magnetic and electric field dependence of the resistivity observed in various two-dimensional (2D) electron and hole systems as well as in superconducting thin films remains still without an accepted explanation. Transport measurements in Si MOSFET’s with large enough carrier density reveal that the resistivity drops by a large factor as the temperature decreases below a temperature $`T\frac{1}{3}T_F`$ ($`T_F`$ is the Fermi temperature) indicative of a clear metallic behavior. Taking into account the observed magnetic field suppression of the metallic behavior , the possible existence of a superconducting phase in two dimensions in Si MOSFET’s has been suggested based mainly on the field-induced superconductor-insulator (SI) transition scaling analysis .
Recent studies of highly oriented pyrolitic graphite (HOPG) suggest the occurrence of high-temperature superconducting correlations in this material . Noting that graphite is a quasi-2D semimetal with a low carrier density (electrons and holes) so that Coulomb interaction effects can be important (the concentration of the majority carriers is $`2\times 10^{18}`$cm<sup>-3</sup> and of the minority carriers is $`6\times 10^{16}`$ cm<sup>-3</sup> ) it seems reasonable to ask whether some of the ideas regarding superconductivity in MOSFET’s may be also applicable to graphite.
In this work we have studied both magnetoresistance and magnetization of HOPG samples with the basal plane resistivity $`\rho `$ ranging from $`\rho (T=300`$K)$`=5\mu \mathrm{\Omega }`$cm to $`135\mu \mathrm{\Omega }`$cm. Our results demonstrate that a field $`H=H_c1`$kOe applied parallel to the hexagonal c-axis, induces a transition from metallic- (d$`R/\mathrm{d}T>0,R`$ is the measured resistance) to semiconducting-type (d$`R/\mathrm{d}T<0`$) behavior. The analysis of the transport data, in particular the obtained scaling reveals a striking similarity to that measured in 2D superconductors and Si MOSFET’s . Here we concentrate mainly on the data obtained for the sample with the lowest resistivity (sample HOPG-3).
HOPG samples were synthesized at temperatures $`T2700\mathrm{}3000^\mathrm{o}`$C and pressure $`P10\mathrm{}30`$MPa at the Research Institute “Graphite” (Moscow) and characterized by means of x-ray diffraction (the analysis of the rocking curves FWHM indicate for the different samples: HOPG-1: $`1.4^o`$, HOPG-2: $`1.2^o`$, HOPG-3: $`0.5^o`$), SEM, STM, and spectrographic analysis techniques . Dc magnetization $`M(H,T)`$ measurements were performed with the SQUID magnetometer MPMS7 (Quantum Design). The dc resistance $`R(H,T)`$ was measured in the temperature interval 2 K - 300 K in applied magnetic field $`0H9`$T using a standard 4-probe method with both current and voltage leads patterned on the sample surface. Resistivity values were obtained with a lead geometry that assures an uniformly distributed current over the sample cross section. The amplitude of the applied current was between 10 $`\mu `$A and 1 mA at which no Joule heating was detected. In all transport measurements the current was inverted in order to eliminate thermoelectric effects.
Figure 1 shows the normalized resistance $`R(T)/R(300`$K) measured for samples with $`\rho (300`$K)$`=45\mu \mathrm{\Omega }`$cm (sample HOPG-1), $`\rho (300`$ K) = 135 $`\mu \mathrm{\Omega }`$cm (sample HOPG-2), and $`\rho (300`$ K) = 5 $`\mu \mathrm{\Omega }`$cm (sample HOPG-3). As can be seen from Fig. 1, at $`T>200`$K all three samples show a semiconducting-type temperature dependence. At lower temperatures $`R(T)`$ obtained for the less resistive samples (HOPG-3 and HOPG-1) crosses over to a metallic behavior (d$`R/\mathrm{d}T>0)`$ below $`T200`$K and $`T50`$K, respectively. However, in the sample HOPG-2 with the highest resistivity, d$`R/\mathrm{d}T<0`$ down to $`T30`$K below which $`R(T)`$ saturates. We stress that the $`R(T)`$-curves shown in Fig. 1 are not specific to our samples: both metallic- and nonmetallic-type $`R(T)`$ generally occur in HOPG depending on the heat treatment . Figure 1 also illustrates that a weak ($`1`$kOe) magnetic field $`H||`$c-axis suppresses the metallic-type state in the HOPG-1 sample.
A pronounced magnetic-field-induced transition from metallic- to semiconducting-type $`R(T)`$ is observed in sample HOPG-3, as shown in Fig. 2(a). The $`R(T)`$ data plotted in Fig. 2(a) for various applied magnetic fields were obtained from isothermal magnetoresistance $`R(H)`$ measurements. Figure 3 shows several $`R(H)`$ isotherms in the vicinity of the “crossover” field $`H_c=1140`$Oe at which the $`R(H)`$ data obtained at temperatures 50 K, 100 K and 200 K intersect.
Noting that the results of Fig. 2(a) resemble very much the resistance behavior in the vicinity of the magnetic-field-driven metal-insulator (MI) and SI transitions we apply here the same scaling approach to the transition measured in HOPG. According to the scaling analysis, the resistance in the critical regime of the quantum SI transition is given by $`R(\delta ,T)=R_cf(|\delta |/T^{1/z\nu })`$, where $`R_c`$ is the resistance at the transition, $`f(|\delta |/T^{1/z\nu })`$ the scaling function such that $`f(0)=1`$, $`z`$ and $`\nu `$ are critical exponents, and $`\delta `$ is the deviation of a variable parameter from its critical value. With $`\delta =HH_c`$, we have plotted in Fig. 2(b) $`R`$ vs. $`|\delta |/T^{1/\alpha }`$, where $`\alpha =0.65\pm 0.05`$ was obtained from the log-log plot of (d$`R/\mathrm{d}H)|_{H_c}`$ vs. $`T^1`$. The collapse of the resistance data into two distinct branches, below and above the critical field $`H_c`$, is obtained in the temperature range 50 K - 200 K, see Fig. 2(b). At $`T<20`$K, where the resistance saturation is apparent, see Fig. 2(a), a clear deviation from the scaling is evident, reminiscent of the behavior observed in amorphous Mo-Ge films where still unknown dissipation effects lead to the saturation of the resistance at low temperatures. Interestingly, the obtained value of the exponent $`\alpha =0.65\pm 0.05`$ for HOPG coincides with that found in the scaling analysis of both the magnetic-field-tuned MI-type transition in Si MOSFETS’s ($`\alpha =0.6\pm 0.1`$ ) and the SI transition in ultrathin a-Bi films ($`\alpha =0.7\pm 0.2`$ ).
We would like to stress the similarities in the transport properties of HOPG and those of a 2D electron gas as in Si MOSFET’s: (1) The occurrence of metallic or semiconducting behavior in graphite at $`H=0`$ (see Fig. 1) is very sensitive to the carrier density which can be changed by, e.g., annealing . (2) The characteristic temperature $`T^{}`$ ($`100`$ K for sample HOPG-3) below which $`\rho (T)`$ strongly decreases is a considerable fraction of the Fermi temperature $`T_F250`$K . (3) There is a clear suppression of the conducting phase by a magnetic field. (4) There exists a critical scaling with a remarkably similar exponent $`\alpha `$.
There is, however, an important difference between the results in HOPG and those obtained in Si MOSFET’s. Whereas the magnetoresistance in Si MOSFET’s is independent of the direction of the applied magnetic field , $`\mathrm{\Delta }R/R`$ measured in HOPG is two orders of magnitude larger for $`H||`$c-axis as compared to the perpendicular orientation. In other words, while a spin-polarization mechanism seems responsible for the transition in Si MOSFET’s, orbital effects are dominant and influence the transition in HOPG.
We note that the scaling shown in Fig. 2(b) can be related to the field-tuned quantum transition if the critical regime extends to $`T200`$K, at least. Although this may appear unlikely, the special characteristic of graphite must be taken into account: Because of the low density and extremely small effective masses of the carriers in graphite, as well as its quasi two-dimensional nature, one expects a strong enhancement of quantum effects. Figure 4 shows the experimental evidence for quantum oscillations due to de Haas-van Alphen (DHVA) effect at $`T=5`$K (a), 100 K (b) and 300 K (c). The occurrence of DHVA oscillations at low fields arising from the minority carriers was already reported in the literature . It is interesting to note, however, that these oscillations become particularly clear at $`HH_c1`$kOe where the SI-type transition takes place, see Fig. 2.
The following points should be considered regarding the origin of possible superconducting correlations in HOPG. (a) We would like to emphasize the potential importance of the minority carriers in HOPG. The effective mass of the minority carriers $`m^{}0.004m_0`$ is 10 to 15 times smaller than the effective mass of the majority holes ($`m^{}0.04m_0`$) and electrons ($`m^{}0.06m_0`$) ($`m_0`$ is the free-electron mass). This may lead to superconducting instabilities in the electron-hole liquid as discussed in Refs.. (b) On the other hand, the coexistence and interplay of superconducting and ferromagnetic states may indicate the relevance of p-wave pairing mediated by ferromagnetic spin fluctuations .
To conclude, the present work provides evidence that at least some of the physical concepts proposed to describe the magnetic-field-tuned SI-type transition in various 2D systems should also be applicable to graphite. The “reentrant” metallic-type behavior observed in the regime of lowest Landau level quantization for both HOPG and Si MOSFET’s further suggests a deep similarity of the physical processes operating in these systems.
We are grateful to A. A. Abrikosov, A. Gerber, I. Ya. Korenblit, I. Luk’yanchuk, A. C. Mota, S. Moehlecke, K. A. Müller, A. Shelankov, M. Sigrist, Z. Tesanovic, and M. Ziese for informative discussions. The authors thank A. S. Kotosonov for providing the samples and V. V. Lemanov for collaboration. This work is supported by the Deutsche Forschungsgemeinschaft under DFG IK 24/B1-1 (project H), and was partially supported by the DAAD and CAPES Proc. No. 077/99 and CNPq proc. No. 301216/93-2. F.M. is supported by the German-Israeli Foundation under G-553-191.14/97.
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# First correlation between compact object and circumstellar disk in the Be/X-ray binaries
## 1 Introduction
The Be/X-ray binaries are the major subclass of massive X-ray binary systems in which a neutron star accretes material from the wind of an early type Be star. The Be stars are known to exhibit emission in the Balmer lines and infrared excess, which are attributed to the presence of cool circumstellar disk. Correlation between the changes of the Be circumstellar envelope and the emission of the compact object can be expected, as a result of the compact object interaction with the surrounding matter. However, no clear correlation has been detected till now – only loose correlations between the optical/infrared properties of the Be circumstellar disks and the X-ray emission of the neutron star have been reported to exist (e.g. Corbet et al. 1985; Coe et al. 1994; Negueruela et al. 1998).
LS I+61$^∘$303 (V615~Cas, GT~0236+610) is a radio emitting X-ray binary which exhibits radio outbursts every $`26.5`$d. The radio outburst peak and the outburst phase are known to vary over a time scale of $``$4 yr (Gregory et al. 1989; Gregory, 1999). Hereafter, we will use the latest values reported and we will refer to these radio periods as $`P_1=26.4917`$ d and $`P_2=1584`$ d. Phase zero for both has been set at JD2443366.775 (Gregory, Peracaula & Taylor, 1999). The $`26.5`$ d period is believed to be the orbital period. The $``$4 yr modulation has been discovered on the basis of continued radio monitoring. Both relativistic jet precession or cyclic variability in the Be star envelope have been proposed as a possible origin of the long term modulation (Paredes, 1987; Gregory et al. 1989), with the second interpretation being the most likely one. This suggestion is supported by the fact that the H$`\alpha `$ emission line varies on the same (4 yr) time scale (Zamanov et al. 1999). However, these authors were not able to derive what is the connection between the radio and H$`\alpha `$ parameters.
In this letter we report an intriguing correlation between the synchrotron non-thermal radio emission, associated with the compact star, and the Be circumstellar disk visible in the $`H\alpha `$ emission line. This is the first clear connection between the Be circumstellar disk variability and the emission from the neutron star in the Be/X-ray binaries.
## 2 $`H\alpha `$ and radio observations
The new $`H\alpha `$ spectroscopic observations used in this paper are obtained with the 2-m RCC telescope of the Bulgarian National Astronomical Observatory ‘Rozhen’ during the last two years. They are analyzed together with the previously published data (Paredes et al. 1994; Zamanov et al. 1999). The H$`\alpha `$ parameters which vary with the $``$4 yr modulation are the equivalent width of H$`\alpha `$, EW(H$`\alpha `$), and the distance between the peaks, $`\mathrm{\Delta }V_{\mathrm{peak}}`$ (Zamanov et al. 1999).
The radio observations during the previous 6 years were retrieved from the Green Bank Interferometer, which is a facility of the USA National Science Foundation operated by NRAO in support of the NASA High Energy Astrophysics program. From this data, we measured the radio outburst peak flux density and the beginning of every outburst. As onset of the outburst, we adopted the time when the radio flux density achieved a value equal to 1/3 of the maximum flux observed for every orbital period. Whenever possible, the onset of the outburst was measured separately from the 2.25 GHz and 8.3 GHz observations and the average value is used in the analysis.
In terms of the Ejector-Propeller model of LS I+61303, the outburst onset better represents the time of the transition of the neutron star from propeller to ejector and the appearance of the expanding radio emitting plasmon (Zamanov, 1995).
## 3 Relation between circumstellar disk and radio outbursts
The radio and $`H\alpha `$ observations obtained during the last years are plotted in Fig.1. It is clearly visible that all these parameters vary over a time scale of $``$4 yr. Following the predictions of Gregory (1999) the radio outburst peak achieved a new maximum at the end of 1999. Big changes in the radio outburst peak flux from values less than 0.150 Jy at JD2451000 to values more than 0.250 Jy at JD2451400 are observed (Fig.1). At the same time, the parameters of the Be circumstellar disk $`EW(H\alpha )`$ and $`\mathrm{\Delta }V_{\mathrm{peak}}`$ do not exhibit extreme values in their behaviour. This suggests that the radio peak flux density and the $`H\alpha `$ parameters do not vary in phase nor in anti-phase.
To understand this behaviour we plotted all the data folded on the $`P_2`$=1584 d period. Every parameter was fitted separately with a cosine function. Our attempts to use more complicated functions (e.g. three term truncated Fourier fit containing sine and cosine terms with half and quarter orbital periods) showed that, with this scatter of the points, the significance of other terms is negligible and we used simple cosine waves in the form: $`y=A+Bcos(2\pi (\varphi +\varphi _0))`$. The best fit parameters are listed in Table 1 and plotted as solid lines in Fig.2.
From Table 1 and from Fig.2 the following correlations can be seen:
a) The H$`\alpha `$ parameters, EW($`H\alpha `$) and $`\mathrm{\Delta }V_{\mathrm{peak}}`$, vary in anti-phase. This is a relationship well known for the Be stars, reflecting the fact that the the size of the disk grows as the EW(H$`\alpha `$) increases (Hanushik, Kozok & Kaizer, 1988).
b) At the same time, a strange phase shift can be seen between the radio outburst peak flux and the H$`\alpha `$ parameters. This shift is about 0.25 in phase. So, if the parameters really vary as simple cosine functions on average, this means that the radio peak varies with the first derivative of the $`H\alpha `$ emission line parameters.
c) The beginning of the radio outburst is in phase, or only slightly shifted, relative to the H$`\alpha `$ variability, at least on the ascending branch.
Ray et al. (1997) have shown that the outburst peak is slowly shifting in phase during the time interval from JD2449400 to JD2450100. In the same interval the outburst peak flux achieved maximum values and began to decline. This is one more confirmation that the variability of the outburst peak and the time of the outburst are phase shifted one to another over the 1584 day period.
Although the scatter of the points in Fig.2 is considerable, all evidences (the cosine fit, the behaviour observed by Ray et al. (1997), the behaviour during the last radio outburst maximum) show that the radio outburst peak is shifted relatively to the behaviour of the circumstellar disk, although the start of the radio outburst is in phase (or almost in phase) with the circumstellar disk changes.
## 4 Discussion
A successful modeling of the radio outburst of LS I+61303 is based on the synchrotron radiation from relativistic particles injected into an expanding plasmon (Paredes et al. 1991). The genesis of the plasmon can be a result of the transition of the neutron star from propeller to ejector state (Zamanov, 1995), or in other words from accretion onto the magnetosphere to ”young radio pulsar” every orbital period.
In such a picture, the start of the outburst will correspond to the moment when the neutron star emerges from the denser parts of the circumstellar disk. Therefore, the bigger the disk the later the outburst can be expected. This can be seen on Fig.2. At phases 0–0.25 we observe increase of the EW(H$`\alpha `$), decrease of the $`\mathrm{\Delta }V_{\mathrm{peak}}`$ and slow shift of the onset of the outburst. The behaviour of the beginning of the outburst on the increasing branch (phases 0–0.25) is stable, but on the decreasing branch (phases 0.25–0.5) the scatter of the points is considerably bigger. The stable behaviour at phases 0.25–0.5 is observed twice at about JD 2449800 and JD 2451400 (Fig.1) so it is unlikely to be a data artifact. Probably this is a result that, during the disk build-up, the increase of the material of the Be disk is feeded only from one source - the B star equatorial region. In contrast the disk-decline can be in two directions - accretion onto the B star or slow dissipation outwards. The behaviour of the start of the outburst points that the disk build up is a stable process and the disk-decline is a more complicated and probably unstable process, or may be it suggests formation of structures like the double disk observed in X Persei (Tarasov & Roche, 1995).
In context of the propeller-ejector transition the surrounding matter will basically influence the expansion velocity of the plasmon, affecting in this way the intensity of the radio outbursts. The remaining plasmon physical parameters (initial magnetic field, injection rate of relativistic particles, etc.) are not expected to vary significantly from one to another outburst. The expanding plasmon calculations predict that there will be weaker outbursts for higher expansion velocities. By expanding faster, the energy losses of the electrons due to the adiabatic expansion are more important and less electron energy is available to be radiated. In addition, the faster decrease of the magnetic field will also contribute to less synchrotron radiation being produced. Supposing that the plasmon is a result of the propeller-ejector transition, the expanding plasmon will appear when the neutron star is receding from the periastron and the plasmon will expand outside of the H$`\alpha `$ emitting disk. The size of this disk is about 40-65$`R_{}`$ (Zamanov & Martí 2000) and the apastron separation between components is about 150$`R_{}`$ (Hutchings & Crampton 1981). The enigmatic behaviour of the outburst peak flux density (its phase shift with 0.25 or $``$400 days) indicates that the conditions outside the H$`\alpha `$ disk vary in a different way compared to the changes inside the H$`\alpha `$ emitting disk.
The X-ray emission of LS I+61303 is observed to exhibit maximum every orbital period and the X-ray outburst is shifted relatively to the radio outburst (Taylor et al. 1996, Harrison et al. 2000). In terms of the ejector-propeller model, the X-ray maximum is due to the propeller action and higher mass accretion rate onto the magnetosphere at the periastron passage (Zamanov & Zamanova, 1997). In this sense it will be very interesting to see what is the behaviour of the X-ray maximum observed in the high-energy emission of LS I+61303 over the $``$4 yr modulation.
Another possible origin of the $``$4 yr modulation may be the precession of the Be star. Lipunov & Nazin (1994) have demonstrated that this value is in rough agreement with the expected period ($`10^3`$d) for precession of the B star. The precession of the B star can be expected, because after the supernova kick the neutron star orbital plane may be different from the Be disk plane (e.g. Bradt & Podsiadlowski, 1995). Our attempts to model the behaviour of the outburst phase as a result of the precession are unsuccessful till now but it can be due to of insufficient data sample, because the systematic radio observations cover about 6 yr (1.5 periods) with considerable gaps and scatter inside the data set. In this context, it deserves to be noted that if the H$`\alpha `$ variability is a result of a precessing disk seen at different inclination angles, this will imply an inclination angle $`i>60^{}`$ and a precession angle $`\mathrm{\Delta }i>6^{}`$ (this estimate is obtained using the values from Table 1 and assuming everywhere an optically thick in H$`\alpha `$ disk).
To conclude, the behaviour of LS I+61303 radio and H$`\alpha `$ emission is evidence that the picture of the interaction between the neutron star and the circumstellar disk in the Be/X-ray binaries is not as simple as generally expected. We need long series of observations over different wavelengths to better understand the behaviour of the Be stars and the Be/X-ray binaries.
###### Acknowledgements.
RZ acknowledges support from Dirección General de Relaciones Culturales y Científicas, Spain. JM acknowledges partial support by DGICYT (PB97-0903) and by Junta de Andalucía (Spain).
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# Low-Temperature Magnetic Penetration Depth in 𝑑-Wave Superconductors: Zero-Energy Bound State and Impurity Effects
## I Introduction
The low-temperature behavior of the magnetic penetration length in $`d`$-wave superconductors is in general a great deal more complicated than that of their isotropic $`s`$-wave cousins. The changing sign of the order parameter, according to where one looks on the Fermi surface, entails coherent zero-energy or low-energy bound states in $`d`$-wave superconductors localized at smooth or almost smooth surfaces or interfaces. These bound states feature peculiar low-temperature contributions to the magnetic penetration length(see also) and the zero-bias conductance peak(see also ).
A minimum in the penetration depth of $`YBa_2Cu_3O_{7\delta }`$ films and grain boundary junctions was thus interpreted as evidence for low-energy Andreev bound states. A conventional shielding-current contribution to the Meissner effect would obviously just monotonically reduce the penetration depth when the temperature goes down. On the other hand, a paramagnetic contribution from low-energy bound states increases the penetration depth. The interplay of these two effects amounts to a minimum in the penetration depth as a function of the temperature. The characteristic temperature $`T_{m0}`$ of this anomaly is shown to be the order $`\sqrt{(\xi _0/\lambda _0)}T_cT_c`$ if the broadening $`\gamma `$ of the bound states is sufficiently small. At this temperature region the bound state contribution to the penetration depth competes with the low-temperature correction from shielding supercurrents to its zero-temperature value.
An alternative explanation of an upturn in the penetration depth is possible in compounds whose bulk paramagnetic properties grow when the temperature goes down, like in the electron-doped cuprate superconductor $`Nd_{1.85}Ce_{0.15}CuO_{4y}`$. There the paramagnetism arises from $`Nd^{3+}`$ ions. We will not discuss these compounds below.
There is yet another important temperature associated with the magnetic penetration depth, $`T_s(\xi _0/\lambda _0)T_c`$. If a given crystal orientation does not carry quasiparticle Andreev bound states, a nonlocal effect can take over as a correction to the zero temperature penetration depth in the clean limit. Then in other orientations which do admit Andreev states, the bound-state contribution and the spontaneous surface supercurrent in particular, can in turn overwhelm the nonlocal effect. At $`TT_s`$ the bound state paramagnetic contribution to $`\lambda `$ in the clean limit is the order of the total London penetration depth $`\lambda _0`$ from the screening currents. In the absence of sub-dominant pairing channels, a spontaneous surface supercurrent brought about by the bound states may arise below the temperature $`T_s`$ (see also on a similar effect of spontaneous magnetization brought about by low energy interface bound states). Having in mind high-temperature superconducting compounds, we will discuss strong type II superconductors. Then $`(\xi _0/\lambda _0)`$ is easily the order $`0.01`$, and the low temperature range splits up into at least three areas staked out by $`0`$, $`T_s`$ and $`T_{m0}`$ ($`T_s(\xi _0/\lambda _0)T_cT_{m0}\sqrt{(\xi _0/\lambda _0)}T_cT_c`$). Quasiparticle scattering off impurities or surface roughness and inelastic processes may also play an important role if they bring about a broadening $`\gamma `$ of the bound states the order or greater than the characteristic temperature $`T_{m0}`$ ($`T_s`$).
We assume below that nonmagnetic impurities dominate the scattering and the broadening. Nonmagnetic impurities in superconductors with an anisotropic order parameter are known to be pair breaking. They suppress $`T_c`$ analogously to what happens to isotropic superconductors with magnetic impurities. Assuming superconductors always clean within the conventional definition $`\xi l`$, we disregard this kind of effects throughout the article. Even then impurity broadening of Andreev bound states in anisotropically paired superconductors can be significant. Since the broadening removes singularities in the density of states (for instance, $`\delta `$-peaks from quasiparticle bound states) as well as in other related physical quantities, superconductors can be sensitive to extremely small concentrations of impurities. This is analogous to the role of pair breaking and small anisotropy of the gap in the Riedel anomaly in isotropic $`s`$-wave superconductors. The Riedel anomaly is associated with the BCS singularity in the density of states. Pair breaking and small anisotropy of the gap are known to wipe out the BCS singularity in the density of states averaged over the Fermi surface, and control the height of the Riedel peak.
The emphasis of the present work is on the various effects of broadening on the low temperature anomalies of the Meissner effect. The zero-energy pole-like term of what is known as the quasiclassical Green’s function was exploited in the investigation. Broadening is introduced into the pole-like term simply sliding the pole along the imaginary energy axis. With small broadening, relatively simple expressions are found for the penetration length in the two lowest-temperature regions defined above. If $`T_{s(m0)}\gamma T_c`$, the growing $`\gamma `$ can wipe out the low-temperature anomalies. Beginning with the critical broadening $`\gamma _{s(m)}`$, anomalies at $`T_{s(m)}`$ are fully destroyed. It turns out that unitary scatterers need to come with significantly larger scattering rates $`\mathrm{\Gamma }_{s(m)}`$ than Born impurities in order to achieve the critical broadening $`\gamma _{s(m)}`$. This effect is peculiar of the impact of impurities on the Andreev bound states as seen in the local density of states and Josephson critical currents. For this reason, the requirements the mean free path must meet for the low temperature anomalies to show up are sensitive to the strength of the impurity potential and very different in the unitary and the Born limits.
For Born scatterers, the shortest normal-state impurity mean free path $`l`$ which preserves the low temperature upturn at $`TT_m`$ is shown to be $`\lambda _0l`$. This looks quite restrictive although conceivably compatible with the strikingly large low temperature mean free paths in some high-$`T_c`$ compounds. For the spontaneous surface supercurrent in the absence of a subdominant component at the surface, we find the threshold $`\lambda _0^2/\xi _0l`$. This demands extraordinary clean samples not available for the time being. On the other hand, the requirements set by unitary scatterers are much weaker and probably can be met. In this case surface roughness is likely to control the broadening and the experimental observability of the effects.
We also examine what the Andreev bound states do to the nonlinear Meissner effect. At low temperatures $`TT_c`$, the field $`\stackrel{~}{H}_0`$ at which the nonlinear response of the bound states saturates in the clean limit is much weaker than the one from the screening current. Ignoring the broadening, $`\stackrel{~}{H}_0`$ is a linear function of the temperature. With $`TT_s`$, nonlinear effects from the bound states become important already in the Meissner state. Close to the transition to the state with a spontaneous surface supercurrent, a nonlinear term entering into the Landau mean-field free energy is important also in a weak external field. The broadening $`\gamma `$ introduces another field, $`\stackrel{~}{H}_\gamma `$ characterizing the nonlinear consequences of the bound states at $`\pi T\gamma `$. For sufficient broadening $`\pi T_s\gamma `$, we get $`H_{c1}\stackrel{~}{H}_\gamma `$ and the nonlinear terms are shown always to be small in the Meissner state.
## II The upturn in the low-temperature dependence of the penetration depth
Our considerations are based on the quasiclassical matrix Green function which describes quasiparticle excitations in thermal equilibrium. The quasiclassical propagator $`\widehat{g}(𝒑_f,x,\epsilon _n)`$ satisfies Eilenberger’s equations, which have a $`2\times 2`$ particle-hole matrix form
$$[\left(i\epsilon _n+\frac{e}{c}𝐯_f𝐀(𝐑)\right)\widehat{\tau }_3\widehat{\mathrm{\Delta }}(𝒑_f,𝐑)\widehat{\sigma }(𝒑_f,𝐑;\epsilon _n),\widehat{g}(𝒑_f,𝐑;\epsilon _n)]+$$
$`+i𝐯_f\mathbf{}_𝐑\widehat{g}(𝒑_f,𝐑;\epsilon _n)=0,`$ (1)
$`\widehat{g}^2(𝒑_f,𝐑;\epsilon _n)=\pi ^2\widehat{1},`$ (2)
where $`\epsilon _n=(2n+1)\pi T`$ are the Matsubara energies, $`𝒑_f`$ the momentum on the Fermi surface, $`𝐯_f`$ the Fermi velocity, $`𝐀`$ the vector potential, $`\widehat{\mathrm{\Delta }}`$ the order parameter matrix, and $`\widehat{\sigma }`$ the impurity self-energy. A symbol with a ‘hat’ denotes a matrix in the Nambu space.
The propagator $`\widehat{g}`$ and the order parameter matrix $`\widehat{\mathrm{\Delta }}`$ parameterize as
$`\widehat{g}=\left(\begin{array}{cc}g& \hfill f\\ f^+& \hfill g\end{array}\right)\text{and}\widehat{\mathrm{\Delta }}=\left(\begin{array}{cc}0& \mathrm{\Delta }\\ \mathrm{\Delta }^{}& 0\end{array}\right).`$ (7)
The gap function $`\mathrm{\Delta }(𝒑_f,𝐑)`$ is related to the anomalous Green function $`f`$ and must be determined self-consistently. The diagonal part $`g(𝒑_f,𝐑,\epsilon _n)`$ of the full matrix propagator $`\widehat{g}`$ carries information on the electrical current density
$$𝒋(𝐑)=2eTN_f\underset{\epsilon _n}{}<𝒗_fg(𝒑_f,𝐑,\epsilon _n)>_{S_f}.$$
(8)
Here $`N_f`$ is the normal state density of states per spin direction and $`<\mathrm{}>_{S_f}`$ means averaging over quasiparticle states at the Fermi surface.
Let an anisotropic singlet strong type II superconductor occupy the right half-space $`x>0`$. A magnetic field is applied along the $`z`$-axis. The induced supercurrent and the vector potential (in the gauge $`\mathrm{div}𝐀(𝐑)=0`$ and vanishing in the bulk) have only $`y`$-components. The linear current-field relation in general has a nonlocal form, i.e. $`j(x)=_0^+\mathrm{}Q(x,x^{},T)A(x^{})𝑑x^{}`$.
For strongly type II superconductors with nodes in the order parameter, a nonlocal current-field relation can be of importance only at very low temperatures $`TT_s`$. Hence, a study of the penetration depth at low temperatures $`T_sTT_{m0}T_c`$ may be carried out disregarding nonlocal effects. Then a magnetic field enters into Eq.(8) only together with the Matsubara frequencies $`\epsilon _ni\frac{e}{c}\mathrm{v}_{f,y}A(x)`$ in the argument of the Green’s function. The kernel $`Q(x,T)`$ can then be written
$$Q(x,T)=\frac{2ie^2TN_f}{c}\underset{n=\mathrm{}}{\overset{+\mathrm{}}{}}<\mathrm{v}_{f,y}^2(𝒑_f)\frac{g(𝒑_f,x,\epsilon _n)}{\epsilon _n}>_{S_f}.$$
(9)
In the presence of zero-energy surface bound states, the pole-like term in the propagator becomes dominating at temperatures $`TT_c`$. Surface bound states as well as their paramagnetic response are localized on the scale of the coherence length at the surface, however, while the conventional screening current has an avenue of the huge thickness of the penetration depth. That is why the zero-energy bound-state contribution to the penetration depth remains a small low-temperature correction to $`\lambda _0\lambda _0(T=0)`$ at all temperatures $`TT_s`$ (in particular, at $`TT_{m0}`$). The contribution from surface bound states must be viewed together with a low temperature correction from the screening current as small low-temperature imports to the zero temperature London penetration depth $`\lambda _0`$. Then the total kernel of the form $`Q(x,T)=\frac{c}{4\pi \lambda _0^2}+\delta Q(x,T)`$ includes only the lowest order corrections in $`\delta Q(x,T)`$.
Solving the Maxwell equation
$$A^{\prime \prime }(x)\frac{1}{\lambda _0^2}A(x)\frac{4\pi }{c}\delta Q(x,T)A(x)=0$$
(10)
perturbatively with respect to the last term delivers a first order approximation to the vector potential:
$$A(x)=A^{(0)}(0)\left[\mathrm{exp}\left(\frac{x}{\lambda _0}\right)\frac{2\pi \lambda _0}{c}\underset{0}{\overset{+\mathrm{}}{}}𝑑x^{}\mathrm{exp}\left(\frac{|xx^{}|}{\lambda _0}\right)\delta Q(x^{},T)\mathrm{exp}\left(\frac{x^{}}{\lambda _0}\right)\right].$$
(11)
The kernel $`\delta Q(x,T)`$ incorporates only a contribution from the bound states and a low-temperature correction from the screening current.
The penetration depth is defined as $`\lambda =_0^+\mathrm{}H(x)𝑑x/H(0)=A(0)/A^{}(0)`$. Expanding this to first order in $`\delta Q`$ and extracting the low temperature correction from the screening current for the case of a superconductor with a line of nodes
$$\lambda (T)=\lambda _0+a\lambda _0\frac{T}{T_c}\frac{4\pi \lambda _0^2}{c}\underset{0}{\overset{\mathrm{}}{}}Q^{bound}(x,T)𝑑x.$$
(12)
Here $`a`$ is a coefficient of the order of unity which depends on the shape of the Fermi surface and on an angular slope of the order parameter near the nodes. For instance, for a quasi-two-dimensional $`d_{x^2y^2}`$ tetragonal superconductor with a cylindrical Fermi surface (with a principal axis $`z`$) and order parameter $`\mathrm{\Delta }(\varphi )=\mathrm{\Delta }_0\mathrm{cos}(2\varphi 2\alpha )`$, one gets $`a0.32`$.
Kernel $`Q^{bound}(x,T)`$ takes negative values. It is a paramagnetic contribution from zero-energy bound states to Eq.(9). One obtains $`Q^{bound}(x,T)`$ from Eq.(9) substituting instead of the full expression for $`g(𝒑_f,x,\epsilon _n)`$ only its singular part (pole-like term) $`g_s(𝒑_f,x,\epsilon _n)`$. Associated with zero energy surface bound states, this term vanishes in the bulk on the scale of the coherence length $`\xi _0`$. It has longer tails only towards the nodes. Node contributions do not dominate, however, in the following expressions. The presence of zero-energy surface bound states is crucial in the reasoning. All sectors of the Fermi surface associated with a sign change of the order parameter in a quasiparticle reflection from the surface, contribute significantly to the results. This allows us to neglect, to a good accuracy, the factor $`\mathrm{exp}(2x/\lambda _0)`$ under the integral sign in Eq.(12).
The analytic expression for the pole-like term has been found in the clean limit and for a smooth surface in Ref. :
$$g_s(𝒑_f,x,\epsilon _n)=\frac{2\pi i}{\epsilon _n}\frac{|\stackrel{~}{\mathrm{\Delta }}(𝒑_f,0)||\stackrel{~}{\mathrm{\Delta }}(\underset{¯}{𝒑}_f,0)|}{|\stackrel{~}{\mathrm{\Delta }}(𝒑_f,0)+|\stackrel{~}{\mathrm{\Delta }}(\underset{¯}{𝒑}_f,0)|}\mathrm{\Theta }(𝒑_f)\mathrm{exp}\left(\frac{2}{|\mathrm{v}_{f,x}(𝒑_f)|}_0^x|\mathrm{\Delta }(𝐩_f,x^{})|𝑑x^{}\right).$$
(13)
The effective surface order parameter $`|\stackrel{~}{\mathrm{\Delta }}(𝒑_f,0)|`$ introduced in Eq.(13), is defined
$$\frac{1}{|\stackrel{~}{\mathrm{\Delta }}(𝒑_f,0)|}=\frac{2}{|\mathrm{v}_{f,x}(𝒑_f)|}\underset{0}{\overset{\mathrm{}}{}}\mathrm{exp}\left(\frac{2}{|\mathrm{v}_{f,x}(𝒑_f)|}\underset{0}{\overset{x}{}}|\mathrm{\Delta }(𝒑_f,x^{})|𝑑x^{}\right)𝑑x.$$
(14)
Here we distinguish between incoming $`𝒑_f`$ and outgoing $`\underset{¯}{𝒑}_f`$ quasiparticle momenta in a reflection event. For specular reflection, the momentum parallel to the interface is conserved. Function $`\mathrm{\Theta }(𝒑_f)`$ is equal to unity where zero energy bound states occur on the Fermi surface (i.e. where the order parameter in the bulk taken for incoming $`𝒑_f`$ and outgoing $`\underset{¯}{𝒑}_f`$ momentum directions have opposite signs), and vanishes elsewhere.
Substituting Eq.(13) in Eq.(9), one can easily sum over the Matsubara frequencies. Integration over the space coordinate $`x`$ in Eq.(12) then yields the penetration depth:
$$\lambda (T)=\lambda _0+a\frac{T}{T_c}\lambda _0+\frac{\pi ^2e^2N_f\lambda _0^2}{c^2T}\mathrm{v}_{f,y}^2(𝒑_f)|\mathrm{v}_{f,x}(𝒑_f)|\mathrm{\Theta }(𝒑_f)_{S_f},T_c\frac{\xi _0}{\lambda _0}TT_c.$$
(15)
For a three dimensional superconductor with a spherical Fermi surface one has the relation $`\lambda _0^2=3c^2/(8\pi e^2\mathrm{v}_f^2N_f)`$. Then the coefficient in front of the third term in Eq.(15) is $`\frac{3\pi }{8T\mathrm{v}_f^2}`$. Analogously, for a simple model of a quasi-two-dimensional superconductor with a cylindrical Fermi surface, $`\lambda _0^2=c^2/(4\pi e^2\mathrm{v}_f^2N_f)`$ and the coefficient $`\frac{\pi }{4T\mathrm{v}_f^2}`$.
In particular, for a $`d_{x^2y^2}`$-wave superconductor with a cylindrical Fermi surface, we get from Eq.(15)
$$\lambda (T)=\lambda _0+a\frac{T}{T_c}\lambda _0+\frac{\mathrm{v}_f}{6T}\left||\mathrm{sin}^3\beta ||cos^3\beta |\right|,T_c\frac{\xi _0}{\lambda _0}TT_c,$$
(16)
where $`\beta =\alpha +(\pi /4)`$ is the angle between the surface normal and the direction to a node of the order parameter, while $`\alpha `$ is the angle between the surface normal and the crystalline $`a`$-axis along its positive lobe.
We note that the correction from zero energy bound states to the penetration depth (the third term in Eq.(15)) has a quite universal form. It is independent both of the spatial profile of the order parameter near a surface and its particular anisotropic structure (basis functions). Therefore, this correction depends only on the type of pairing, which determines regions on the Fermi surface with opposite signs of the order parameter. For example, expression (16) is valid irrespective of a particular form of a momentum direction dependence of the basis function for a $`d`$-wave order parameter of given symmetry.
The ratio of a supercurrent density at the surface $`j_s^{bound}(x=0,T)`$ to the one $`j_s(x_{scr},T)`$ at a characteristic distance $`x_{scr}`$ ($`\xi _0x_{scr}\lambda _0`$) from the surface can be estimated for a clean superconductor at $`TT_c`$ and a smooth surface as $`|j_s^{bound}(x=0,T)/j_s(x_{scr},T)|4\pi \lambda _0^2|Q^{bound}(x=0,T)|/cT_c/T`$. This verifies that at low temperatures $`TT_c`$ the paramagnetic current $`j_s^{bound}(x,T)`$ dominates over the shielding current near the surface within a relatively small characteristic scale $`\xi _0`$.
The temperature dependent terms in Eq.(15) behave in very different fashions from each other. They come from the conventional shielding currents and from the zero-energy bound states. Growing with decreasing temperature, the diamagnetic screening currents monotonically reduce the penetration depth. On the other hand, Andreev surface-bound states respond paramagnetically and increase the penetration depth when the temperature goes down. Disregarding the broadening effects, Eq.(15) delivers the following estimate for the field of the low-temperature minimum of the penetration depth:
$$T_{m0}=\zeta \sqrt{\frac{\xi _0}{\lambda _0}}T_c,$$
(17)
where $`\zeta `$ is of the order of unity for crystalline orientations with sufficient amount of momentum directions admitting zero-energy bound states. Otherwise $`\zeta `$ is a small quantity. For a $`d`$-wave superconductor $`\zeta \left||\mathrm{sin}^3\beta ||cos^3\beta |\right|^{1/2}`$ and vanishes for $`\beta =45^{}`$ (i.e. for $`\alpha =0`$), when there are no zero energy bound states.
Broadening of the bound states can substantially modify the conditions for the presence of a minimum in the low temperature dependence of the penetration depth. For a small broadening $`\gamma (𝒑_f)T_c`$ we simply replace the factor $`\frac{1}{\epsilon _n}`$ in the expression Eq.(13) for the pole-like term with $`\frac{1}{\left[\epsilon _n+\gamma \left(𝒑_f\right)\mathrm{sgn}\left(\epsilon _n\right)\right]}`$. Taking into account the broadening Eq.(15) is generalized to the following form:
$$\lambda (T)=\lambda _0+a\frac{T}{T_c}\lambda _0+\frac{2e^2N_f\lambda _0^2}{c^2T}\mathrm{v}_{f,y}^2(𝒑_f)|\mathrm{v}_{f,x}(𝒑_f)|\mathrm{\Theta }(𝒑_f)\psi ^{}\left(\frac{1}{2}+\frac{\gamma (𝒑_f)}{2\pi T}\right)_{S_f}.$$
(18)
Here and below $`\psi (x)`$ is the digamma function and $`\psi ^{}(x)`$ \- its derivative.
Eq.(18) is a reasonable representation of the role of a broadening in the low temperature anomaly of the penetration depth. The minimum lies at $`T_{m0}1.8\sqrt{\xi _0/\lambda _0}T_c`$ for momentum independent broadening in a $`d_{x^2y^2}`$-superconductor in the clean limit $`\gamma \pi T`$ with the orientation $`\alpha =45^{}`$. With increasing broadening it drifts to lower temperatures (becoming less pronounced at the same time) till $`T_{m\gamma }0.4\sqrt{\xi _0/\lambda _0}T_c`$ at $`\gamma 0.96T_{m0}`$, where it evaporates. As an example, the low-temperature correction to the penetration depth is shown in Fig. 1 in the vicinity of $`T_{m\gamma }`$ for various values of the momentum independent broadening.
There are various contributions to the broadening of the bound states associated, in particular, with surface roughness, nonmagnetic and magnetic impurities and inelastic scattering. We now pin-point the origin of the broadening, assuming that nonmagnetic impurities dominate the scene.
With Born scatterers $`\gamma _b\sqrt{\frac{T_c}{\tau }}`$ (see Ref. ) and the coefficient of the order of unity can be estimated within the simple model of spatially constant order parameter. Then we easily get the shortest normal-state impurity mean free path $`l`$ which admits a low-temperature upturn: $`\lambda _0l`$. In high-temperature superconductors one should distinguish between $`l`$ and the actual mean free path in the normal state at $`T_c`$ incorporating significant contributions from inelastic processes. Impurity scattering dominates there at low temperatures already in the superconducting state where the collapse of inelastic scattering takes place. For instance, below 20 K in $`YBa_2Cu_3O_{7\delta }`$ there is a regime of extremely long and weakly temperature dependent quasiparticle scattering times usually interpreted as due to feeble impurity scattering in high-purity samples.
For scatterers with sufficient strength of impurity potential there are practically no restrictions on the impurity scattering rate in contrast to what was found above for Born impurities. For unitary scatterers with scattering rates $`\mathrm{\Gamma }_uT_c`$ the broadening of the zero-energy bound states is exponentially small: $`\gamma _u=B\sqrt{\mathrm{\Delta }_0\mathrm{\Gamma }_u}\mathrm{exp}\left(b\mathrm{\Delta }_0/\mathrm{\Gamma }_u\right)`$. A scattering rate $`\mathrm{\Gamma }_u`$ which leads to a given broadening $`\gamma _u`$ is almost independent of a constant coefficient $`B`$ in the pre-exponential factor, while it is sensitive to the model dependent parameter $`b`$ in the argument of the exponential function. Within the simple model considered in Ref. , one gets $`b1`$.
For temperatures $`T\sqrt{\mathrm{\Gamma }_u\mathrm{\Delta }_0}`$, the share of the penetration depth from the shielding currents must be modified for unitary scatterers. This leads instead of the linear term in Eqs.(12), (15), (16), (18) to a quadratic low temperature correction of the form $`\lambda _0T^2/\left(\mathrm{\Gamma }^{1/2}\mathrm{\Delta }_0^{3/2}\right)`$ to within a factor the order of unity. Correspondingly, $`T_{m0}`$ given in Eq.(17) is valid for unitary scatterers only if $`T_{m0}>\sqrt{\mathrm{\Gamma }_u\mathrm{\Delta }_0}`$, which sets an upper limit on the scattering rate: $`\mathrm{\Gamma }_u<\frac{\xi _0}{\lambda _0}\mathrm{\Delta }_0`$.
The $`T^2`$-term instead of the linear one in Eq.(15) delivers an estimate for the location $`T_{md}`$ of the low temperature minimum of the penetration depth modified by unitary scatterers:
$$T_{md}\left(\frac{\xi _0}{\lambda _0}\right)^{1/3}\left(\mathrm{\Gamma }_u\mathrm{\Delta }_0^5\right)^{1/6},\mathrm{\Gamma }_u>\frac{\xi _0}{\lambda _0}\mathrm{\Delta }_0.$$
(19)
This expression replaces Eq.(17) in the case of unitary scatterers with the scattering rate $`\mathrm{\Gamma }_u>\frac{\xi _0}{\lambda _0}\mathrm{\Delta }_0`$. The minimum slowly drifts to higher temperatures with increasing $`\mathrm{\Gamma }_u`$. It does not melt away at any $`\mathrm{\Gamma }_uT_c`$. The normal-state impurity mean free path must just be large on the scale of the coherence length.
We conclude that observation of the low temperature upturn of the penetration depth in samples with $`l<\lambda _0`$ is evidence for both Andreev bound states and a sufficiently large strength of the bulk impurity potential in the superconducting compounds. For unitary impurities one needs to take into account the broadening that arises from surface roughness which then very probably controls the total broadening. The same effect with Born scatterers demands the normal-state impurity mean free path larger than the London penetration depth.
## III Zero energy bound states and spontaneous surface current
Throughout this section the broadening of the zero-energy bound states is assumed small. We look at a clean $`d`$-wave superconductor with a smooth surface. Its crystal-to-surface orientation shall admit zero-energy surface bound states and feature an upturn in the penetration depth. Below the upturn temperature $`T_{m0}`$, imagine a great deal of space for $`\lambda `$ to grow, firstly as described by the perturbative result Eq.(15). Then a second order phase transition occurs at $`TT_sT_{m0}`$ into a state which carries a spontaneous surface supercurrent. We shall find an analytic expression for the transition temperature and discuss the impact of impurities on the effect. The transition implies the absence of subdominant channels activated at low temperatures close to the surface on account of the presumably large surface pair breaking in the dominant component of the order parameter. Otherwise a spontaneous current can arise at higher temperatures.
There is experimental evidence for a phase transition on the (110) surface in $`YBa_2Cu_3O_{7\delta }`$ at $`T=7K(\xi _0/\lambda _0)T_c`$. It was interpreted as associated with an activated near surface subdominant channel of the order parameter. For some other crystal-to-surface orientations, however, a subdominant component can be not present near a surface. Zero energy bound states can still arise for a noticeable part of quasiparticle trajectories. Our theoretical study is relevant to these cases.
In order to find an equation for the transition temperature, one has to admit a paramagnetic contribution to the penetration depth at least as large as the diamagnetic one. Then a perturbation treatment of the preceding section is not adequate. In this context we develop an approach based on the integral form of Eq.(10) and take into account only the terms in $`\delta Q(x,T)`$ brought about by the bound states. In other words, a contribution only from the pole-like term Eq.(13) needs to be taken into account in Eqs.(9) for the kernel which enters into Eq.(10). The kernel $`\delta Q(x,T)`$ varies on the characteristic scale $`\xi _0`$ and is associated in the clean limit with large contributions to the magnetic field at the surface at temperatures $`T(\xi _0/\lambda _0)T_c`$. We therefore disregard the nonlocal temperature correction from the Meissner current to $`\delta Q(x,T)`$.
We transform Eq.(10) into the integral form
$$A(x)=\left[A(0)\frac{2\pi \lambda _0}{c}\underset{0}{\overset{+\mathrm{}}{}}𝑑x^{}Q^{bound}(x^{},T)A(x^{})\left(e^{\frac{x^{}}{\lambda _0}}e^{\frac{x^{}}{\lambda _0}}\right)\right]e^{\frac{x}{\lambda _0}}$$
$$\frac{2\pi \lambda _0}{c}\underset{x}{\overset{+\mathrm{}}{}}𝑑x^{}Q^{bound}(x^{},T)A(x^{})\left(e^{\frac{xx^{}}{\lambda _0}}e^{\frac{xx^{}}{\lambda _0}}\right).$$
(20)
The two terms on the right hand side of this equation obey very different scales. The first decays exponentially in the depth on the scale $`\lambda _0`$ while the last term vanishes for $`x\xi _0`$ along with the kernel $`Q^{bound}(x,T)`$. The kernel $`Q^{bound}(x=0,T)`$ can be estimated (see preceding section) for a clean superconductor and a smooth surface as $`2\pi \lambda _0^2Q^{bound}(x=0,T)/cT_c/T`$. Then, in accordance with Eq.(20), the approximate formula $`\left(1\frac{A\left(x_{scr}\right)}{A\left(0\right)}\right)\xi _0^2T_c/\lambda _0^2T`$ is established for a relative deviation of the vector potential $`A(x_{scr})`$ taken at the distance $`x_{scr}`$ ($`\xi x_{scr}\lambda _0`$) from its value $`A(0)`$ at the surface. The deviation reflecting the bound state contribution to the vector potential turns out to be small at all temperatures $`T\left(\xi _0^2/\lambda _0^2\right)T_c`$, in particular, for $`TT_s\left(\xi _0/\lambda _0\right)T_c`$. Varying on the scale $`\xi _0`$, small terms in the expression for the vector potential at temperatures $`TT_s`$ are of importance only when differentiating $`A(x)`$. After that they can already noticeably contribute to the expression for the magnetic field.
Indeed, a spatial differentiation of Eq.(20) leads to
$$H(x)=\frac{1}{\lambda _0}\left[A(0)\frac{2\pi \lambda _0}{c}\underset{0}{\overset{+\mathrm{}}{}}𝑑x^{}Q^{bound}(x^{},T)A(x^{})\left(e^{\frac{x^{}}{\lambda _0}}e^{\frac{x^{}}{\lambda _0}}\right)\right]e^{\frac{x}{\lambda _0}}$$
$$\frac{2\pi }{c}\underset{x}{\overset{+\mathrm{}}{}}𝑑x^{}Q^{bound}(x^{},T)A(x^{})\left(e^{\frac{xx^{}}{\lambda _0}}+e^{\frac{xx^{}}{\lambda _0}}\right).$$
(21)
The second term in the square brackets remains negligibly small $`\left(\xi _0^2T_c/\lambda _0^2T\right)A(0)A(0)`$ as compared with $`A(0)`$ for $`T\left(\xi _0^2/\lambda _0^2\right)T_c`$. The last term of Eq.(21) is the order of $`\left(\xi _0T_c/\lambda _0T\right)(A(0)/\lambda _0)`$. For a deviation of the magnetic field at $`x=x_{scr}`$ from its value at $`x=0`$$`\left(\frac{H\left(x_{scr}\right)}{H\left(0\right)}1\right)\left(\xi _0T_c/\lambda _0T\right)`$. Hence, the bound state contribution to the magnetic field can be viewed as a small perturbation as compared with the shielding contribution unless $`TT_s`$. Considering $`\left(\xi _0^2/\lambda _0^2\right)T_cTT_s`$, we can discard the second term in the square brackets but have to keep track of the last term in Eq.(21). Choosing $`x=0`$ in Eq.(21), the small terms $`x^{}/\lambda _0`$ in the exponential functions under the integral sign can be taken to vanish. For the same reason and within the same accuracy one can treat the vector potential under the integral sign in Eq.(21) as constant in space $`A(0)`$ discarding small terms in the vector potential which vary on the scale $`\xi _0`$. All this results in an explicit relation between $`A(0)`$ and $`H(0)`$ and therefore
$$\lambda =\frac{\lambda _0}{1+{\displaystyle \frac{4\pi \lambda _0}{c}}{\displaystyle \underset{0}{\overset{+\mathrm{}}{}}}Q^{bound}(x,T)𝑑x}.$$
(22)
Proceeding like in the derivation of Eq.(15) above, we find that the paramagnetic (negative) sign of $`Q^{bound}`$ leads to a divergence of $`\lambda `$ at the temperature
$$T_s=\frac{\pi ^2e^2N_f\lambda _0}{c^2}\mathrm{v}_{f,y}^2(𝒑_f)|\mathrm{v}_{f,x}(𝒑_f)|\mathrm{\Theta }(𝒑_f)_{S_f}.$$
(23)
For the model $`d`$-wave superconductor with a cylindrical Fermi surface one gets from Eq.(23)
$$T_s=\frac{\pi \xi _0}{3\lambda _0}\left||\mathrm{sin}^3\beta ||cos^3\beta |\right|T_c,$$
(24)
where $`\xi _0=\mathrm{v}_f/2\pi T_c`$.
The divergence of $`\lambda `$ implies the existence of a nontrivial solution to Eq.(20) in a vanishing external magnetic field. Indeed, if we let $`H(0)=0`$, $`A(0)0`$, then Eq.(21) transforms, with the same approximation as above, into the relation: $`1=\left(4\pi \lambda _0/c\right)_0^+\mathrm{}Q^{bound}(x^{},T)𝑑x^{}`$, which defines the transition temperature $`T_s`$ into a state with a spontaneous surface supercurrent.
The nontrivial solution at $`T_s`$ is a result of interplay between the paramagnetic supercurrent which originates in the zero energy bound states localized within $`\xi _0`$ on the one hand and the diamagnetic supercurrent distributed over the region $`x\lambda _0`$ on the other. The latter compensates for the magnetic field from the bound states at the surface in order to satisfy the boundary conditions in the absence of an external magnetic field. Then $`\underset{0}{\overset{+\mathrm{}}{}}j(x)𝑑x=0`$ always applies being a consequence of the full screening of the spontaneous surface magnetic field in the bulk of a superconductor. Under these conditions the Bloch theorem, in general, admits spontaneous surface currents. The magnetic part of a superconducting free energy $`\frac{1}{8\pi }\underset{0}{\overset{+\mathrm{}}{}}\left[A^2(x)+\frac{4\pi }{c}Q(x,T)A^2(x)\right]𝑑x`$ vanishes at $`T=T_s`$ and becomes negative below $`T_s`$ on account of negative sign of the paramagnetic kernel $`Q^{bound}`$ (Gibbs and Helmholtz free energies coincide in zero external magnetic field). The result is an energetically favorable state with a spontaneous surface supercurrent below $`T_s`$.
The broadening $`\gamma `$ of the bound states modifies Eq.(23):
$$T_s=\frac{2e^2N_f\lambda _0}{c^2}\mathrm{v}_{f,y}^2(𝒑_f)|\mathrm{v}_{f,x}(𝒑_f)|\mathrm{\Theta }(𝒑_f)\psi ^{}\left(\frac{1}{2}+\frac{\gamma (𝒑_f)}{2\pi T_s}\right)_{S_f}.$$
(25)
The broadening prevents the appearance of a spontaneous surface current unless $`\gamma \frac{\xi _0}{\lambda _0}T_c`$. This is a very strong restriction. If Born impurities control the broadening, they admit spontaneous surface supercurrent only with extremely large values of the mean free path $`\lambda _0^2/\xi _0100\lambda _0l`$, unrealistic for high temperature superconductors. Unitary scatterers impose a much weaker restriction $`\mathrm{\Gamma }_u2b\mathrm{\Delta }_0/\mathrm{ln}\left[\frac{\lambda _0^2}{\xi _0^2\mathrm{ln}\left(\lambda _0/\xi _0\right)}\right]0.1T_c`$. Then, however, surface roughness probably dominates the broadening and can destroy the state with a spontaneous surface supercurrent.
## IV Nonlinear Meissner effect from low energy bound states
It is important in the derivation of Eq.(22) that the kernel $`Q^{bound}`$ varies much faster in space than the screening currents. Then contributions of the paramagnetic current carried by surface Andreev-bound states at temperatures $`(\xi _0/\lambda _0)^2T_cT`$, can result in significant spatial variations of the magnetic field near the surface while in the weakly spatially dependent vector potential. This leads to Eq.(22) on the basis of the local current-field relation.
It is straightforward to show within the same framework that a nonlinear penetration depth $`\lambda _{nl}(T,H)`$ incorporating contributions both from screening currents and from zero-energy bound states is described as
$$\lambda _{nl}(T,H)=\frac{\lambda _{nl}^{scr}(T,H_{scr})}{1+{\displaystyle \frac{4\pi \lambda _{nl}^{scr}(T,H_{scr})}{c}}{\displaystyle \underset{0}{\overset{+\mathrm{}}{}}}Q_{nl}^{bound}(x,T)𝑑x},$$
(26)
where $`\lambda _{nl}^{scr}(T,H_{scr})`$ is a contribution from screening supercurrents to $`\lambda _{nl}(T,H)`$, taken at an effective value of the field $`H_{scr}=H(0)\frac{4\pi }{c}\underset{0}{\overset{+\mathrm{}}{}}H(x)𝑑x\underset{0}{\overset{+\mathrm{}}{}}Q_{nl}^{bound}(x^{},T)𝑑x^{}`$. Here $`H(0)`$ is the external magnetic field. The second term describes the field of the zero-energy bound states inside the superconductor at distances $`xx_{scr}`$ ($`\xi _0x_{scr}\lambda _0`$), as can be seen in Eq.(21). A paramagnetic response of zero-energy bound states ($`Q_{nl}^{bound}<0`$) increases the field to be screened by diamagnetic supercurrents ($`H_{scr}H(x_{scr})>H(0)`$). This leads, in general, to more pronounced nonlinear terms in $`\lambda _{nl}^{scr}(T,H_{scr})`$ as compared to disregarding the contribution from zero-energy bound states. In the case of spontaneous surface supercurrent $`H_{scr}`$ differs from zero even in the absence of an external field. We assume the condition $`H_{scr}<H_{c1}`$ for the Meissner state to be stable in the magnetic field on account of a paramagnetic influence of the bound states.
Small nonlocal low temperature corrections to the penetration depth from screening currents can be taken into account in Eq.(26) as perturbations to $`\lambda _{nl}^{scr}(T,H)`$. For a nonlocal current-field relation a penetration depth $`\lambda _{nl}(T,H)`$ is actually a functional of the spatial profile of the magnetic field.
Nonlinear corrections from the shielding supercurrent to the Meissner effect can be given in terms of the dimensionless ratio $`\rho =(H/H_0)`$, where $`H_0`$ is usually the order of the thermodynamic critical field $`\mathrm{\Phi }_0/(\lambda _0\xi _0)`$. Hence, they are always small in strong type II superconductors in the Meissner state. In isotropic $`s`$-wave superconductors, the first nonlinear correction to the penetration depth $`\rho ^2`$. In superconductors with nodes in the order parameter (for instance, $`d`$-wave ) a term linear in $`\rho `$ can arise for particular crystal orientations at low temperatures. The linear term, however, is quite sensitive to nonlocal effects and the impurity influence, in particluar, at sufficiently large strength of impurity potentials.
A nonlinearity in the magnetic response of low energy Andreev surface bound states has, in general, a very different field scale $`\stackrel{~}{H}_0`$. In a clean limit $`\stackrel{~}{H}_0(T)=(\mathrm{\Phi }_0T/\lambda v_f)`$, where $`\lambda `$ is determined by Eq.(26). At low temperatures $`TT_c`$ one always gets $`\stackrel{~}{H}_0(T)H_0`$. For instance, $`\stackrel{~}{H}_0(T_{m0})\sqrt{\frac{\xi _0}{\lambda _0}}H_00.1H_0`$, $`\stackrel{~}{H}_0(T_s)\frac{\lambda _0}{\lambda }H_{c1}H_{c1}0.01H_0`$. Moreover, at sufficiently low temperatures $`TT_s`$, i.e. well below the transition to the state with a spontaneous surface supercurrent, a paramagnetic response from the bound states may become seriously nonlinear already in the Meissner state $`\stackrel{~}{H}_0(T)H<H_{c1}`$. We will show, however, that the broadening of the bound states introduces a new field scale $`\stackrel{~}{H}_\gamma =\frac{\gamma \lambda _0}{T_c\lambda }H_0`$ coming into play at $`\pi T\gamma `$. For $`\gamma \frac{\xi _0}{\lambda _0}T_c0.01T_c`$ nonlinear corrections from Andreev low-energy bound states to the penetration length turn out always to be small in the Meissner state, even at $`T=0`$.
As a pole-like term Eq.(13) decays exponentially on the scale $`\xi _0`$ for almost all momentum directions admitting bound states, we consider a local nonlinear current-field relation
$$𝒋(x)=2eTN_f\underset{\epsilon _n}{}𝒗_fg_s(𝒑_f,x,\epsilon _ni\frac{e}{c}\mathrm{v}_{f,y}A(x))_{S_f}$$
(27)
for the current via the bound states. One can also set the vector potential $`A(x)`$ in the kernel equal to $`A(x=0)`$. Then we easily generalize the reasoning in the derivation of the third term in Eq.(18). Substituting into Eq.(27) the expression Eq.(13) for the pole-like term with the pole shifted in accordance with the broadening, we find:
$$\underset{0}{\overset{\mathrm{}}{}}Q_{nl}^{bound}(x,T)𝑑x=\frac{ieN_f}{A(0)}\mathrm{v}_{f,y}(𝒑_f)|\mathrm{v}_{f,x}(𝒑_f)|\mathrm{\Theta }(𝒑_f)\psi \left(\frac{1}{2}+\frac{\gamma (𝒑_f)+i{\displaystyle \frac{e}{c}}\mathrm{v}_{f,y}(𝒑_f)A(0)}{2\pi T}\right)_{S_f}.$$
(28)
In Eq.(28) the Fermi surface is assumed symmetric in reflections across the $`xz`$-plane. Then averages over the Fermi surface of odd powers of $`\mathrm{v}_{f,y}`$ vanish, no matter whether they are multiplied by $`|\mathrm{v}_{f,x}(𝒑_f)|\mathrm{\Theta }(𝒑_f)`$ or not. This applies, in particular, to a $`d`$-wave superconductor with a cylindrical Fermi surface whose principal axis $`z`$ is parallel to the boundary for arbitrary orientations of the two other crystal axes $`x_0`$, $`y_0`$.
If $`\frac{e}{c}\mathrm{v}_fA(0)\mathrm{max}(2\pi T,\gamma )`$, one can expand the $`\psi `$-function in Eq.(28) in powers of the small parameter $`\mathrm{min}(H(0)/\stackrel{~}{H}_0(T),H(0)/\stackrel{~}{H}_\gamma )`$. Considering nonlinear corrections to the penetration depth from screening currents $`\mathrm{\Delta }\lambda _{nl}^{scr}`$ and bound states $`\mathrm{\Delta }\lambda _{nl}^b`$ to be small, one can represent them in the first approximation as additive contributions to the total penetration depth $`\lambda _{nl}(T,H)\lambda (T)+\mathrm{\Delta }\lambda _{nl}^{scr}+\mathrm{\Delta }\lambda _{nl}^b`$. The nonlinear correction from screening currents takes the form $`\mathrm{\Delta }\lambda _{nl}^{scr}=\frac{\lambda ^2\left(T\right)}{\lambda _0^2\left(T\right)}\left[\lambda _{nl}^{scr}(T,H(0)\frac{\lambda \left(T\right)}{\lambda _0\left(T\right)})\lambda _0(T)\right]`$. Quantities $`\lambda (T)`$ and $`\lambda _0(T)`$ being the zero-field values of $`\lambda _{nl}(T,H)`$, $`\lambda _{nl}^{scr}(T,H)`$ respectively, satisfy Eq.(22). Bound states renormalize nonlinear response from screening currents already in this approximation. Thus, the explicit analysis of $`\mathrm{\Delta }\lambda _{nl}^{scr}`$ can be done combining the results of the preceding section and Refs. . Apart from too close to the transition temperature $`T_s`$, the nonlinear correction to the penetration depth from the bound states is:
$$\mathrm{\Delta }\lambda _{nl}^b\frac{e^4\lambda ^4(T)N_fH^2(0)}{12\pi ^2c^4T^3}\mathrm{v}_{f,y}^4(𝒑_f)|\mathrm{v}_{f,x}(𝒑_f)|\mathrm{\Theta }(𝒑_f)\psi ^{(3)}\left(\frac{1}{2}+\frac{\gamma (𝒑_f)}{2\pi T}\right)_{S_f}.$$
(29)
In the limit $`\frac{e}{c}\mathrm{v}_fA(0)2\pi T`$ when the argument of the $`\psi `$-function in Eq.(28) is large, we obtain
$$\underset{0}{\overset{\mathrm{}}{}}Q_{nl}^{bound}(x,T)𝑑x=\frac{eN_f}{A(0)}\mathrm{v}_{f,y}(𝒑_f)|\mathrm{v}_{f,x}(𝒑_f)|\mathrm{\Theta }(𝒑_f)\mathrm{arctan}\left(\frac{e\mathrm{v}_{f,y}(𝒑_f)A(0)}{c\gamma (𝒑_f)}\right)_{S_f}.$$
(30)
Then the broadening rather than the temperature fixes the bound state contribution to the penetration depth.
As shown above, there is no state with a spontaneous surface current with $`\gamma T_s`$. Then $`\lambda \lambda _0`$ and $`e\mathrm{v}_fA(0)/c\gamma H(0)/H_{c1}`$. Since $`H(0)<H_{c1}`$ in the Meissner state we estimate $`e\mathrm{v}_fA(0)/c\gamma 1`$ and obtain in this limit from Eq.(30)
$$\mathrm{\Delta }\lambda _\gamma ^b=\frac{4\pi e^2N_f\lambda _0^2}{c^2\gamma }\mathrm{v}_{f,y}^2(𝒑_f)|\mathrm{v}_{f,x}(𝒑_f)|\mathrm{\Theta }(𝒑_f)_{S_f}$$
$$\frac{4\pi e^4\lambda _0^4N_fH^2(0)}{3c^4\gamma ^3}\mathrm{v}_{f,y}^4(𝒑_f)|\mathrm{v}_{f,x}(𝒑_f)|\mathrm{\Theta }(𝒑_f)_{S_f}.$$
(31)
For a given $`\lambda _{nl}^{scr}(T,H)`$, Eq.(26) in general should be solved with respect to $`\lambda _{nl}(T,H)`$ in accordance with Eq.(28), since $`A(0)=\lambda _{nl}(T,H)H(0)`$. This is particularly important close to the transition temperature $`T_s`$, where the Landau theory of second order phase transitions is applicable. Then first nonlinear term turns out to be the order of the zero-field paramagnetic contribution in the denominator in Eq.(26). Ignoring a weak field dependence of $`\lambda _{nl}^{scr}(T,H)`$ stipulated by screening currents, we obtain from Eqs.(26), (28) the following equation for $`\lambda _{nl}(T,H)`$:
$$\left(\frac{T}{T_s}1\right)\kappa \lambda _{nl}(T,H)+\eta H^2\lambda _{nl}^3(T,H)=\lambda _{nl}^{scr}(T_s),$$
(32)
where $`H=H(0)`$,
$$\eta =\frac{e^4\lambda _{nl}^{scr}(T_s)N_f}{12\pi ^2c^4T_s^3}\mathrm{v}_{f,y}^4(𝒑_f)|\mathrm{v}_{f,x}(𝒑_f)|\mathrm{\Theta }(𝒑_f)\psi ^{(3)}\left(\frac{1}{2}+\frac{\gamma (𝒑_f)}{2\pi T_s}\right)_{S_f},\kappa =1\frac{d\stackrel{~}{T}_s(T)}{dT}|_{T=T_s}.$$
$`T_s`$ is described by Eq.(25) and $`\stackrel{~}{T}_s(T)`$ is the result of the substitution $`T_sT`$ on the right hand side of Eq.(25). The broadening is assumed to be sufficiently small for admitting the phase transition.
The role of the order parameter in the phase transition can be played by a surface magnetization $`m_S=\underset{0}{\overset{+\mathrm{}}{}}\left(M(x)M_{\mathrm{}}\right)𝑑x=\frac{1}{c}\underset{0}{\overset{+\mathrm{}}{}}𝑑xxj_s(x)=\frac{1}{4\pi }\underset{0}{\overset{+\mathrm{}}{}}H(x)𝑑x=\frac{1}{4\pi }\lambda _{nl}(T,H)H(0)`$, which, for simplicity, we choose constant in space along a smooth surface. The magnetization $`𝐌`$ enters by the conventional definition $`𝐣=c\mathrm{curl}𝐌`$, and $`M_{\mathrm{}}=H(0)/4\pi `$. Then the Landau free energy per unit surface $`_S`$ which leads to the same equation for $`m_S`$ as implied in Eq.(32) has the form
$$_S=\stackrel{~}{\alpha }\left(\frac{T}{T_s}1\right)m_S^2+\stackrel{~}{\beta }m_S^4m_SH,$$
(33)
where $`\stackrel{~}{\alpha }=2\pi \kappa /\lambda _{nl}^{scr}(T_s)`$, $`\stackrel{~}{\beta }=16\pi ^3\eta /\lambda _{nl}^{scr}(T_s)`$, $`H`$ is the external field. As for a conventional order parameter in a strong field near $`T_s`$ one gets $`m_S(T_s,H)H^{1/3}`$, which entails $`\lambda _{nl}(T_s,H)H^{2/3}`$.
Finally, in the limit of very small broadening $`\gamma \frac{H\left(0\right)}{H_0}T_c`$, Eqs.(26) and (30) give
$$\mathrm{\Delta }\lambda _H^b=\lambda _{nl}(T,H)\lambda _{nl}^{scr}(T,H)=\frac{4\pi eN_f\lambda _0}{c|H(0)|}|\mathrm{v}_{f,y}(𝒑_f)\mathrm{v}_{f,x}(𝒑_f)|\mathrm{\Theta }(𝒑_f)_{S_f}$$
(34)
at temperatures $`\frac{\xi _0^2}{\lambda _0^2}T_cT\frac{H\left(0\right)}{H_0}T_c`$. Since $`\mathrm{\Delta }\lambda _H^b`$ is at least the order of $`\lambda _0`$, we put here $`\lambda _{nl}^{scr}(T,H)\lambda _0`$ disregarding small nonlinear corrections from the screening currents in the Meissner state. The approximate inverse proportionality of the penetration length the magnetic field implies the presence of a spontaneous surface magnetization weakly dependent on $`H`$.
## V Conclusion
We have examined the paramagnetic contribution from surface zero-energy Andreev bound states to the low-temperature penetration length of $`d`$-wave superconductors in the Meissner state. The paramagnetic current is localized within several coherence lengths near the surface and grows larger in the clean limit when the temperature goes down. A broadening of the bound states chokes their contribution and determines their actual role in shaping the penetration length. We found that the upturn in the low temperature penetration depth lies at $`T_{m0}\sqrt{\xi _0/\lambda _0}T_c`$ in the clean limit where the paramagnetic contribution from the bound states can be handled with perturbation theory same as small low-temperature corrections to the penetration depth from the screening current. The minimum broadening capable of straightening out the upturn is $`\gamma T_{mo}`$.
Furthermore, we examined the penetration depth when the bound states must be kept track of beyond perturbation theory. A divergence of $`\lambda (T)`$ was found at the phase transition to a state with spontaneous surface supercurrent. This transition occurs only with smallish broadening, $`\gamma <(\xi _0/\lambda _0)T_c`$. In the clean limit and at low temperatures, there is a nonlinear regime of the paramagnetic current already in magnetic fields substantially weaker than the fields for the nonlinear effects to show up in response of shielding supercurrents. The broadening of the bound states modifies and weakens the nonlinear response.
Specifying an origin of the broadening as associated with nonmagnetic impurity scattering, we obtained restrictions on the mean free path admitting the low temperature anomalies. The conditions turn out to be sensitive to the strength of the impurity potential and very different in the unitary and in the Born limits. The Born impurities are shown to easily prevent the anomalies of the penetration depth taking place at least well below $`T_{m0}`$. By contrast, unitary scatterers with sufficiently small normal-state scattering rate $`\mathrm{\Gamma }_uT_c`$ admit the transition to a state with spontaneous surface supercurrent at $`T_s(\xi _0/\lambda _0)T_c`$. In the latter case, however, surface roughness very probably dominates the broadening and controls the bound state contribution to the low-temperature penetration length.
## Acknowledgments
We thank M. Fogelström for useful discussions. This work was supported by the Academy of Finland, research Grant No. 4385. Yu.S.B. acknowledges the financial support of Russian Foundation for Basic Research under grant No. 99-02-17906.
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# Gauge Field Theory Coherent States (GCS) : II. Peakedness Properties
## 1 Introduction
Quantum General Relativity (QGR) has matured over the past decade to a mathematically well-defined theory of quantum gravity. In contrast to string theory, by definition QGR is a manifestly background independent, diffeomorphism invariant and non-perturbative theory. The obvious advantage is that one will never have to postulate the existence of a non-perturbative extension of the theory, which in string theory has been called the still unknown M(ystery)-Theory.
The disadvantage of a non-perturbative and background independent formulation is, of course, that one is faced with new and interesting mathematical problems so that one cannot just go ahead and “start calculating scattering amplitudes”: As there is no background around which one could perturb, rather the full metric is fluctuating, one is not doing quantum field theory on a spacetime but only on a differential manifold. Once there is no (Minkowski) metric at our disposal, one loses familiar notions such as causality structure, locality, Poincaré group and so forth, in other words, the theory is not a theory to which the Wightman axioms apply. Therefore, one must build an entirely new mathematical apparatus to treat the resulting quantum field theory which is drastically different from the Fock space picture to which particle physicists are used to.
As a consequence, the mathematical formulation of the theory was the main focus of research in the field over the past decade. The main achievements to date are the following (more or less in chronological order) :
* Kinematical Framework
The starting point was the introduction of new field variables for the gravitational field which are better suited to a background independent formulation of the quantum theory than the ones employed until that time. In its original version these variables were complex valued, however, currently their real valued version, considered first in for classical Euclidean gravity and later in for classical Lorentzian gravity, is preferred because to date it seems that it is only with these variables that one can rigorously define the kinematics and dynamics of Euclidean or Lorentzian quantum gravity .
These variables are coordinates for the infinite dimensional phase space of an $`SU\left(2\right)`$ gauge theory subject to further constraints besides the Gauss law, that is, a connection and a canonically conjugate electric field. As such, it is very natural to introduce smeared functions of these variables, specifically Wilson loop and electric flux functions. (Notice that one does not need a metric to define these functions, that is, they are background independent). This had been done for ordinary gauge fields already before in and was then reconsidered for gravity (see e.g. ).
The next step was the choice of a representation of the canonical commutation relations between the electric and magnetic degrees of freedom. This involves the choice of a suitable space of distributional connections and a faithful measure thereon which, as one can show , is $`\sigma `$-additive. The proof that the resulting Hilbert space indeed solves the adjointness relations induced by the reality structure of the classical theory as well as the canonical commutation relations induced by the symplectic structure of the classical theory can be found in . Independently, a second representation of the canonical commutation relations, called the loop representation, had been advocated (see e.g. and especially and references therein) but both representations were shown to be unitarily equivalent in (see also for a different method of proof).
This is then the first major achievement : The theory is based on a rigorously defined kinematical framework.
* Geometrical Operators
The second major achievement concerns the spectra of positive semi-definite, self-adjoint geometrical operators measuring lengths , areas and volumes of curves, surfaces and regions in spacetime. These spectra are pure point (discete) and imply a discrete Planck scale structure. It should be pointed out that the discreteness is, in contrast to other approaches to quantum gravity, not put in by hand but it is a prediction !
* Regularization- and Renormalization Techniques
The third major achievement is that there is a new regularization and renormalization technique for diffeomorphism covariant, density-one-valued operators at our disposal which was successfully tested in model theories . This technique can be applied, in particular, to the standard model coupled to gravity and to the Poincaré generators at spatial infinity . In particular, it works for Lorentzian gravity while all earlier proposals could at best work in the Euclidean context only (see, e.g. and references therein). The algebra of important operators of the resulting quantum field theories was shown to be consistent . Most surprisingly, these operators are UV and IR finite ! Notice that, at least as far as these operators are concerned, this result is stronger than the believed but unproved finiteness of scattering amplitudes order by order in perturbation theory of the five critical string theories, in a sense we claim that the perturbation series converges. The absence of the divergences that usually plague interacting quantum fields propagating on a Minkowski background can be understood intuitively from the diffeomorphism invariance of the theory : “short and long distances are gauge equivalent”. We will elaborate more on this point in future publications.
* Spin Foam Models
After the construction of the densely defined Hamiltonian constraint operator of , a formal, Euclidean functional integral was constructed in and gave rise to the so-called spin foam models (a spin foam is a history of a graph with faces as the history of edges) . Spin foam models are in close connection with causal spin-network evolutions , state sum models and topological quantum field theory, in particular BF theory . To date most results are at a formal level and for the Euclidean version of the theory only but the programme is exciting since it may restore manifest four-dimensional diffeomorphism invariance which in the Hamiltonian formulation is somewhat hidden.
* Finally, the fifth major achievement is the existence of a rigorous and satisfactory framework for the quantum statistical description of black holes which reproduces the Bekenstein-Hawking Entropy-Area relation and applies, in particular, to physical Schwarzschild black holes while stringy black holes so far are under control only for extremal charged black holes.
Summarizing, the work of the past decade has now culminated in a promising starting point for a quantum theory of the gravitational field plus matter and the stage is set to pose and answer physical questions.
The most basic and most important question that one should ask is : Does the theory have classical general relativity as its classical limit ? Notice that even if the answer is negative, the existence of a consistent, interacting, diffeomorphism invariant quantum field theory in four dimensions is already a quite non-trivial result. However, we can claim to have a satisfactory quantum theory of Einstein’s theory only if the answer is positive.
It seems that the most natural framework for deriving the classical limit of a theory is based on coherent states or best approximation states. Coherent states have a long history and an extensive literature exists in a vast range of applications (see e.g. and references therein). It has been pointed out by many (see e.g. ) that they are best suited for the analysis of the semi-classical behaviour of any given system because, among other things, in contrast to the WKB-methods more familiar to physicists they avoid the discussion of the critical turning points and it is much more natural to ask questions which address regions in the classical phase space rather than in configuration and momentum space only.
Surprisingly, the vast majority of coherent states have been constructed for systems with only a finite number of degrees of freedom. This is astonishing because in the course of constructions of (interacting) quantum field theories from given classical ones one is almost always forced to regularize and renormalize the operators in that theory and these are operations which have no classical counterpart. Thus, it would be no surprise if it turned out that the classical limit of such quantum field theories is not the classical field theory that one started from. Just to give an example, even if one could rigorously show that the continuum limit of lattice QCD exists, to the best of the knowledge of the authors it is at present unclear whether the classical limit of that continuum quantum field theory would give us back classical $`SU\left(3\right)`$ Yang-Mills theory coupled to quarks.
This paper is the second one in a series of papers entitled “Gauge Field Theory Coherent States” which are geared at shedding light at these questions. Specifically, we are interested in the question whether the non-perturbative quantization of continuum Lorentzian general relativity in four dimensions with and without matter advertized in has the correct classical limit. In fact we eliminate the criticism stated in and show in that quantum general relativity as presently formulated does admit graviton states which would then presumably also enable us to make contact with results from perturbation theory.
The general outline of our programme was given in where a huge family of coherent states, based on the phase space complexifier method , was introduced. Here we specialize to the “heat kernel family” of coherent states that results by choosing the square of electric flux variables as the complexifier which, upon quantization, becomes a Laplacian. This choice is motivated, on the one hand by the beautiful analysis of Hall who established a unitary transfomation between square integrable functions on a compact gauge group with respect to the Haar measure and square integrable, holomorphic functions on the complexified group with respect to the so-called heat kernel measure. On the other hand, it is convenient since an application of this framework to diffeomorphism invariant gauge theories hs already been started in in .
The original purpose of was to solve the reality conditions of quantum general relativity written in terms of the complex valued Ashtekar connection and therefore the properties of the states that came with that heat kernel transform remained untouched. Moreover, the heat kernel transform of obviously complexifies the real connection but it remained unclear how that complex valued connection is expressed in terms of the coordinates of the real phase space. Without that knowledge there is obviously no interpretation of that complex valued connection possible. In this paper we will fill both of these gaps. Namely, using the classical framework of and the complexifier method of we explicitly construct the complex connection out of the real phase space variables. Secondly, we analyze in detail the semi-classical properties of the coherent states so obtained, most importantly their peakedness properties.
This we do in great detail for the compact gauge groups of rank one, that is, $`U\left(1\right)`$ and $`SU\left(2\right)`$, and sketch how the proofs extend to compact groups of higher rank. Details will appear in the forthcoming paper . Coherent states for Higgs fields are completely analogous to the coherent states constructed here because one can describe them by so-called “point-holonomies” which are a special case of the holonomies considered here. Details and coherent states for fermions are treated in .
As it will become obvious, the states constructed in this paper can serve as a tool to perform error-controlled rigorous approximations in quantum general relativity and quantum gauge theory coupled to quantum gravity and therefore as a starting point for numerical canonical quantum general relativity and numerical canonical quantum gauge theory coupled to quantum gravity.
The present article is organized as follows :
Section two is an account of the relevant notions and techniques of non-perturbative classical and quantum general relativity.
Section three explicitly derives the particular complexification of the real phase space of gauge theories or real general relativity based on heat kernel generators as complexifiers. This section depends on the recently constructed theory of symplectic manifolds of quantum general relativity and quantum gauge theory labelled by graphs .
Section four introduces the heat kernel family of gauge-non-invariant states for a general gauge theory without fermions in any spacetime dimension and we prove that they satisfy all the properties that one is used to from the classical harmonic oscillator coherent states. That is, these states are labelled by a classical connection and a classical electric field (a point in phase space) and we show that these states are peaked on these values in the connection-, momentum- and Segal-Bargmann representation. Furthermore, we show that the system of states is overcomplete, saturates the unquenched Heisenberg uncertainty bound with respect to certain complexified holonomy operators and that each state labelled by a point in phase space can be associated with a phase space cell with a volume whose size is controlled by $`\mathrm{}^d`$. We do all this for the gauge group $`SU\left(2\right)`$ and point out how to generalize to an arbitrary compact gauge group.
In section five the analysis of section four 3 is generalized to the gauge invariant heat kernel family. The proofs follow essentially from the proofs derived in section four by employing the group averaging method of refined algebraic quantization (RAQ) . However, the results stated in section five are somewhat less complete than those for section four due to the difficulty to do the group averaging explicitly which makes it hard to establish sharp peakedness. Fortunately, the results of section four are completely sufficient in order to study the semi-classical behaviour of the theory.
Finally in Appendix A we repeat our analysis for the technically much simpler case of $`G=U\left(1\right)`$ and in Appendix B we display the peakedness properties of the states constructed in the configuration and Bargmann-Segal representation graphically, both for $`SU\left(2\right)`$ and $`U\left(1\right)`$. All graphics have been obtained by means of Mathematica and the admittedly large amount of plots is justified by the fact that, to the best of our knowledge, the behaviour of these states has not been studied numerically before.
## 2 Kinematical Structure of Diffeomorphism Invariant Quantum Gauge Theories
In this section we will recall the main ingredients of the mathematical formulation of (Lorentzian) diffeomorphism invariant classical and quantum field theories of connections with local degrees of freedom in any dimension and for any compact gauge group. See and references therein for more details. Also, in this section we will take all quantities to be dimensionless for simplicity, the incoporation of dimensionful parameters will be discussed in the next section.
### 2.1 Classical Theory
Let $`G`$ be a compact gauge group, $`\mathrm{\Sigma }`$ a $`D`$dimensional manifold admitting a principal $`G`$bundle with connection over $`\mathrm{\Sigma }`$. Let us denote the pull-back to $`\mathrm{\Sigma }`$ of the connection by local sections by $`A_a^i`$ where $`a,b,c,..=1,..,D`$ denote tensorial indices and $`i,j,k,..=1,..,dim\left(G\right)`$ denote indices for the Lie algebra of $`G`$. Likewise, consider a vector bundle of electric fields, whose projection to $`\mathrm{\Sigma }`$ is a Lie algebra valued vector density of weight one. We will denote the set of generators of the rank $`N1`$ Lie algebra of $`G`$ by $`\tau _i`$ which are normalized according to $`\text{tr}\left(\tau _i\tau _j\right)=N\delta _{ij}`$ and $`[\tau _i,\tau _j]=2f_{ij}^k\tau _k`$ defines the structure constants of $`Lie\left(G\right)`$.
Let $`F_i^a`$ be a Lie algebra valued vector density test field of weight one and let $`f_a^i`$ be a Lie algebra valued covector test field. We consider the smeared quantities
$$F\left(A\right):=_\mathrm{\Sigma }d^DxF_i^aA_a^i\text{ and }E\left(f\right):=_\mathrm{\Sigma }d^DxE_i^af_a^i$$
(2.1)
While both are diffeomorphism covariant, it is only the latter which is gauge covariant, one reason to consider the singular smearings discussed below. The choice of the space of pairs of test fields $`(F,f)𝒮`$ depends on the boundary conditions on the space of connections and electric fields which in turn depends on the topology of $`\mathrm{\Sigma }`$ and will not be specified in what follows.
Consider the set $`M`$ of all pairs of smooth functions $`(A,E)`$ on $`\mathrm{\Sigma }`$ such that (2.1) is well defined for any $`(F,f)𝒮`$. We define a topology on $`M`$ through the following globally defined metric :
$`d_{\rho ,\sigma }[(A,E),(A^{},E^{})]`$
$`:=`$ $`\sqrt{{\displaystyle \frac{1}{N}}{\displaystyle _\mathrm{\Sigma }}d^Dx\left[\sqrt{det\left(\rho \right)}\rho ^{ab}\text{tr}\left(\left[A_aA_a^{}\right]\left[A_bA_b^{}\right]\right)+{\displaystyle \frac{[\sigma _{ab}\text{tr}\left([E^aE^a][E^bE^b]\right)}{\sqrt{det\left(\sigma \right)}}}\right]}`$
where $`\rho _{ab},\sigma _{ab}`$ are fiducial metrics on $`\mathrm{\Sigma }`$ of everywhere Euclidean signature. Their fall-off behaviour has to be suited to the boundary conditions of the fields $`A,E`$ at spatial infinity. Notice that the metric (2.1) on $`M`$ is gauge invariant. It can be used in the usual way to equip $`M`$ with the structure of a smooth, infinite dimensional differential manifold modelled on a Banach (in fact Hilbert) space $``$ where $`𝒮\times 𝒮`$. (It is the weighted Sobolev space $`H_{0,\rho }^2\times H_{0,\sigma ^1}^2`$ in the notation of ).
Finally, we equip $`M`$ with the structure of an infinite dimensional symplectic manifold through the following strong (in the sense of ) symplectic structure
$$\mathrm{\Omega }((f,F),(f^{},F^{}))_m:=_\mathrm{\Sigma }d^Dx\left[F_i^af_a^iF_i^af_a^i\right]\left(x\right)$$
(2.3)
for any $`(f,F),(f^{},F^{})`$. We have abused the notation by identifying the tangent space to $`M`$ at $`m`$ with $``$. To prove that $`\mathrm{\Omega }`$ is a strong symplectic structure one uses standard Banach space techniques. Computing the Hamiltonian vector fields (with respect to $`\mathrm{\Omega }`$) of the functions $`E\left(f\right),F\left(A\right)`$ we obtain the following elementary Poisson brackets
$$\{E\left(f\right),E\left(f^{}\right)\}=\{F\left(A\right),F^{}\left(A\right)\}=0,\{E\left(f\right),A\left(F\right)\}=F\left(f\right)$$
(2.4)
As a first step towards quantization of the symplectic manifold $`(M,\mathrm{\Omega })`$ one must choose a polarization. As usual in gauge theories, we will use connections as the configuration variables and electric fields as canonically conjugate momenta. As a second step one must decide on a complete set of coordinates of $`M`$ which are to become the elementary quantum operators. The analysis just outlined suggests to use the coordinates $`E\left(f\right),F\left(A\right)`$. However, the well-known immediate problem is that these coordinates are not gauge covariant. Thus, we proceed as follows :
Let $`\mathrm{\Gamma }_0^\omega `$ be the set of all piecewise analytic, finite, oriented graphs $`\gamma `$ embedded into $`\mathrm{\Sigma }`$ and denote by $`E\left(\gamma \right)`$ and $`V\left(\gamma \right)`$ respectively its sets of oriented edges $`e`$ and vertices $`v`$ respectively. Here finite means that $`E\left(\gamma \right)`$ is a finite set. (One can extend the framework to $`\mathrm{\Gamma }_0^{\mathrm{}}`$, the restriction to webs of the set of piecewise smooth graphs but the description becomes more complicated and we refrain from doing this here). It is possible to consider the set $`\mathrm{\Gamma }_\sigma ^\omega `$ of piecewise analytic, infinite graphs with an additional regularity property but for the purpose of this paper it will be sufficient to stick to $`\mathrm{\Gamma }_0^\omega `$. The subscript <sub>0</sub> as usual denotes “of compact support” while <sub>σ</sub> denotes “$`\sigma `$-finite”.
We denote by $`h_e\left(A\right)`$ the holonomy of $`A`$ along $`e`$ and say that a function $`f`$ on $`𝒜`$ is cylindrical with respect to $`\gamma `$ if there exists a function $`f_\gamma `$ on $`G^{\left|E\left(\gamma \right)\right|}`$ such that $`f=p_\gamma ^{}f_\gamma =f_\gamma p_\gamma `$ where $`p_\gamma \left(A\right)=\left\{h_e\left(A\right)\right\}_{eE\left(\gamma \right)}`$. Holonomies are invariant under reparameterizations of the edge and in this article we assume that the edges are always analyticity preserving diffeomorphic images from $`[0,1]`$ to a one-dimensional submanifold of $`\mathrm{\Sigma }`$. Gauge transformations are functions $`g:\mathrm{\Sigma }G;xg\left(x\right)`$ and they act on holonomies as $`h_eg\left(e\left(0\right)\right)h_eg\left(e\left(1\right)\right)^1`$.
Next, given a graph $`\gamma `$ we choose a polyhedronal decomposition $`P_\gamma `$ of $`\mathrm{\Sigma }`$ dual to $`\gamma `$. The precise definition of a dual polyhedronal decomposition can be found in but for the purposes of the present paper it is sufficient to know that $`P_\gamma `$ assigns to each edge $`e`$ of $`\gamma `$ an open “face” $`S_e`$ (a polyhedron of codimension one embedded into $`\mathrm{\Sigma }`$) with the following properties :
(1) the surfaces $`S_e`$ are mutually non-intersecting,
(2) only the edge $`e`$ intersects $`S_e`$, the intersection is transversal and consists only of one point which is an interiour point of both $`e`$ and $`S_e`$,
(3) $`S_e`$ carries the orientation which agrees with the orientation of $`e`$.
Furthermore, we choose a system $`\mathrm{\Pi }_\gamma `$ of paths $`\rho _e\left(x\right)S_e,xS_e,eE\left(\gamma \right)`$ connecting the intersection point $`p_e=eS_e`$ with $`x`$. The paths vary smoothly with $`x`$ and the triples $`(\gamma ,P_\gamma ,\mathrm{\Pi }_\gamma )`$ have the property that if $`\gamma ,\gamma ^{}`$ are diffeomorphic, so are $`P_\gamma ,P_\gamma ^{}`$ and $`\mathrm{\Pi }_\gamma ,\mathrm{\Pi }_\gamma ^{}`$.
With these structures we define the following function on $`(M,\mathrm{\Omega })`$
$$P_i^e(A,E):=\frac{1}{N}\text{tr}\left(\tau _ih_e(0,1/2)\left[_{S_e}h_{\rho _e\left(x\right)}E\left(x\right)h_{\rho _e\left(x\right)}^1\right]h_e(0,1/2)^1\right)$$
(2.5)
where $`h_e(s,t)`$ denotes the holonomy of $`A`$ along $`e`$ between the parameter values $`s<t`$, $``$ denotes the Hodge dual, that is, $`E`$ is a $`\left(D1\right)`$form on $`\mathrm{\Sigma }`$, $`E^a:=E_i^a\tau _i`$ and we have chosen a parameterization of $`e`$ such that $`p_e=e\left(1/2\right)`$.
Notice that in contrast to similar variables used earlier in the literature the function $`P_i^e`$ is gauge covariant. Namely, under gauge transformations it transforms as $`P^eg\left(e\left(0\right)\right)P^eg\left(e\left(0\right)\right)^1`$, the price to pay being that $`P^e`$ depends on both $`A`$ and $`E`$ and not only on $`E`$. The idea is therefore to use the variables $`h_e,P_i^e`$ for all possible graphs $`\gamma `$ as the coordinates of $`M`$.
The problem with the functions $`h_e\left(A\right)`$ and $`P_i^e(A,E)`$ on $`M`$ is that they are not differentiable on $`M`$, that is, $`Dh_e,DP_i^e`$ are nowhere bounded operators on $``$ as one can easily see. The reason for this is, of course, that these are functions on $`M`$ which are not properly smeared with functions from $`𝒮`$, rather they are smeared with distributional test functions with support on $`e`$ or $`S_e`$ respectively. Nevertheless one would like to base the quantization of the theory on these functions as basic variables because of their gauge and diffeomorphism covariance. Indeed, under diffeomorphisms $`h_eh_{\phi ^1\left(e\right)},P_j^eP_j^{\phi ^1\left(e\right)}`$ where we abuse notation since $`P^e`$ depends also explicitly on the $`S_e,\rho _e`$, see for details. We proceed as follows.
###### Definition 2.1
By $`\overline{M}_\gamma `$ we denote the direct product $`[G\times Lie(G)]^{|E(\gamma )|}`$. The subset of $`\overline{M}_\gamma `$ of pairs $`(h_e(A),P_i^e(A,E))_{eE(\gamma )}`$ as $`(A,E)`$ varies over $`M`$ will be denoted by $`(\overline{M}_\gamma )_{|M}`$. We have a corresponding map $`p_\gamma :M\overline{M}_\gamma `$ which maps $`M`$ onto $`(\overline{M}_\gamma )_{|M}`$.
Notice that the set $`\left(\overline{M}_\gamma \right)_{|M}`$ is in general a proper subset of $`M_\gamma `$, depending on the boundary conditions on $`(A,E)`$, the topology of $`\mathrm{\Sigma }`$ and the “size” of $`e,S_e`$. For instance, in the limit of $`e,S_eeS_e`$ but holding the number of edges fixed, $`\left(\overline{M}_\gamma \right)_{|M}`$ will consist of only one point in $`\overline{M}_\gamma `$. This follows from the smoothness of the $`(A,E)`$.
We equip a subset $`M_\gamma `$ of $`\overline{M}_\gamma `$ with the structure of a differentiable manifold modelled on the Banach space $`_\gamma =\text{ }\mathrm{R}^{2dim\left(G\right)\left|E\left(\gamma \right)\right|}`$ by using the natural direct product manifold structure of $`\left[G\times Lie\left(G\right)\right]^{\left|E\left(\gamma \right)\right|}`$. While $`\overline{M}_\gamma `$ is a kind of distributional phase space, $`M_\gamma `$ satisfies appropriate regularity properties similar to $`M`$.
In order to proceed and to give $`M_\gamma `$ a symplectic structure derived from $`(M,\mathrm{\Omega })`$ one must regularize the elementary functions $`h_e,P_i^e`$ by writing them as limits (in which the regulator vanishes) of functions which can be expressed in terms of the $`F\left(A\right),E\left(f\right)`$. Then one can compute their Poisson brackets with respect to the symplectic structure $`\mathrm{\Omega }`$ at finite regulator and then take the limit pointwise on $`M`$. The result is the following well-defined strong symplectic structure $`\mathrm{\Omega }_\gamma `$ on $`M_\gamma `$.
$`\{h_e,h_e^{}\}_\gamma `$ $`=`$ $`0`$
$`\{P_i^e,h_e^{}\}_\gamma `$ $`=`$ $`\delta _e^{}^e{\displaystyle \frac{\tau _i}{2}}h_e`$
$`\{P_i^e,P_j^e^{}\}_\gamma `$ $`=`$ $`\delta ^{ee^{}}f_{ij}^kP_k^e`$ (2.6)
Since $`\mathrm{\Omega }_\gamma `$ is obviously block diagonal, each block standing for one copy of $`G\times Lie\left(G\right)`$, to check that $`\mathrm{\Omega }_\gamma `$ is non-degenerate and closed reduces to doing it for each factor together with an appeal to well-known Hilbert space techniques to establish that $`\mathrm{\Omega }_\gamma `$ is a surjection of $`_\gamma `$. This is done in where it is shown that each copy is isomorphic with the cotangent bundle $`T^{}G`$ equipped with the symplectic structure (2.1) (choose $`e=e^{}`$ and delete the label $`e`$).
Now that we have managed to assign to each graph $`\gamma `$ a symplectic manifold $`(M_\gamma ,\mathrm{\Omega }_\gamma )`$ we can quantize it by using geometric quantization. This can be done in a well-defined way because the relations (2.1) show that the corresponding operators are non-distributional. This is therefore a clean starting point for the regularization of any operator of quantum gauge field theory which can always be written in terms of the $`\widehat{h}_e,\widehat{P}^e,eE\left(\gamma \right)`$ if we apply this operator to a function which depends only on the $`h_e,eE\left(\gamma \right)`$.
As an example , recall that $`(M_\gamma ,\mathrm{\Omega }_\gamma )`$ is subject to a coisotropic constraint, the Gauss constraint, which in terms of the quantities defined above can be written
$$G\left(\mathrm{\Lambda }\right)=\underset{vV\left(\gamma \right)}{}\mathrm{\Lambda }^i\left[\underset{eE\left(\gamma \right),e\left(0\right)=v}{}P_i^e\underset{eE\left(\gamma \right),e\left(1\right)=v}{}O_{ij}\left(h_e\right)P_j^e\right]$$
(2.7)
where the smooth, Lie-algebra valued function of rapid decrease $`\mathrm{\Lambda }`$ is a test function on $`\mathrm{\Sigma }`$ enforcing the local constraint
$$G_i\left(v\right)=\underset{eE\left(\gamma \right),e\left(0\right)=v}{}P_i^e\underset{eE\left(\gamma \right),e\left(1\right)=v}{}O_{ij}\left(h_e\right)P_j^e$$
(2.8)
where $`O_{ij}\left(h\right)=\text{tr}\left(h\tau _ih^1\tau _j\right)/N`$. Since $`G\left(\mathrm{\Lambda }\right)`$ is coisotropic, specifically
$$\{G\left(\mathrm{\Lambda }\right),G\left(\mathrm{\Lambda }^{}\right)\}=G\left([\mathrm{\Lambda },\mathrm{\Lambda }^{}]\right)$$
(2.9)
the dimension of the physical configuration space equals half the dimension of $`M_\gamma `$ (which is $`Edim\left(G\right)`$) minus $`Vdim\left(G\right)`$, the number of constraints. The question is what $`(M_\gamma ,\mathrm{\Omega }_\gamma )`$ has to do with $`M,\mathrm{\Omega }`$. In it is shown that there exists a partial order $``$ on the set $``$ of triples $`l=(\gamma ,P_\gamma ,\mathrm{\Pi }_\gamma )`$. In particular, $`\gamma \gamma ^{}`$ means $`\gamma \gamma ^{}`$ and $``$ is a directed set so that one can form a generalized projective limit $`M_{\mathrm{}}`$ of the $`M_\gamma `$ (we abuse notation in displaying the dependence of $`M_\gamma `$ on $`\gamma `$ only rather than on $`l`$). For this one verifies that the family of symplectic structures $`\mathrm{\Omega }_\gamma `$ is self-consistent in the sense that if $`(\gamma ,P_\gamma ,\mathrm{\Pi }_\gamma )(\gamma ^{},P_\gamma ^{},\mathrm{\Pi }_\gamma ^{})`$ then $`p_{\gamma ^{}\gamma }^{}\{f,g\}_\gamma =\{p_{\gamma ^{}\gamma }^{}f,p_{\gamma ^{}\gamma }^{}g\}_\gamma ^{}`$ for any $`f,gC^{\mathrm{}}\left(M_\gamma \right)`$ and $`p_{\gamma ^{}\gamma }:M_\gamma ^{}M_\gamma `$ is a system of natural projections, more precisely, of (non-invertible) symplectomorphisms.
Now, via the maps $`p_\gamma `$ of definition 2.1 we can identify $`M`$ with a subset of $`M_{\mathrm{}}`$. Moreover, in it is shown that there is a generalized projective sequence $`(\gamma _n,P_{\gamma _n},\mathrm{\Pi }_{\gamma _n})`$ such that $`lim_n\mathrm{}p_{\gamma _n}^{}\mathrm{\Omega }_{\gamma _n}=\mathrm{\Omega }`$ pointwise in $`M`$. This displays $`(M,\mathrm{\Omega })`$ as embedded into a generalized projective limit of the $`(M_\gamma ,\mathrm{\Omega }_\gamma )`$, intuitively speaking, as $`\gamma `$ fills all of $`\mathrm{\Sigma }`$, we recover $`(M,\mathrm{\Omega })`$ from the $`(M_\gamma ,\mathrm{\Omega }_\gamma )`$. Of course, this works with $`\mathrm{\Gamma }_0^\omega `$ only if $`\mathrm{\Sigma }`$ is compact, otherwise we need the extension to $`\mathrm{\Gamma }_\sigma ^\omega `$.
It follows that quantization of $`(M,\mathrm{\Omega })`$, and conversely taking the classical limit, can be studied purely in terms of $`(M_\gamma ,\mathrm{\Omega }_\gamma )`$ for all $`\gamma `$. The quantum kinematical framework for this will be given in the next subsection.
### 2.2 Quantum Theory
Let us denote the set of all smooth connections by $`𝒜`$. This is our classical configuration space and we will choose for its coordinates the holonomies $`h_e\left(A\right),e\gamma ,\gamma \mathrm{\Gamma }_0^\omega `$. $`𝒜`$ is naturally equipped with a metric topology induced by (2.1).
Recall the notion of a function cylindrical over a graph from the previous subsection. A particularly useful set of cylindrical functions are the so-called spin-netwok functions . A spin-network function is labelled by a graph $`\gamma `$, a set of non-trivial irreducible representations $`\stackrel{}{\pi }=\left\{\pi _e\right\}_{eE\left(\gamma \right)}`$ (choose from each equivalence class of equivalent representations once and for all a fixed representant), one for each edge of $`\gamma `$, and a set $`\stackrel{}{c}=\left\{c_v\right\}_{vV\left(\gamma \right)}`$ of contraction matrices, one for each vertex of $`\gamma `$, which contract the indices of the tensor product $`_{eE\left(\gamma \right)}\pi _e\left(h_e\right)`$ in such a way that the resulting function is gauge invariant. We denote spin-network functions as $`T_I`$ where $`I=\{\gamma ,\stackrel{}{\pi },\stackrel{}{c}\}`$ is a compound label. One can show that these functions are linearly independent. From now on we denote by $`\stackrel{~}{\mathrm{\Phi }}_\gamma `$ finite linear combinations of spin-network functions over $`\gamma `$, by $`\mathrm{\Phi }_\gamma `$ the finite linear combinations of elements from any possible $`\stackrel{~}{\mathrm{\Phi }}_\gamma ^{},\gamma ^{}\gamma `$ a subgraph of $`\gamma `$ and by $`\mathrm{\Phi }`$ the finite linear combinations of spin-network functions over an arbitrary collection of graphs. Clearly $`\stackrel{~}{\mathrm{\Phi }}_\gamma `$ is a subspace of $`\mathrm{\Phi }_\gamma `$. To express this distinction we will say that functions in $`\stackrel{~}{\mathrm{\Phi }}_\gamma `$ are labelled by the “coloured graphs” $`\gamma `$ while functions in $`\mathrm{\Phi }_\gamma `$ are labelled simply by graphs $`\gamma `$ where we abuse notation by using the same symbol $`\gamma `$.
The set $`\mathrm{\Phi }`$ of finite linear combinations of spin-network functions forms an Abelian algebra of functions on $`𝒜`$. By completing it with respect to the sup-norm topology it becomes an Abelian C algebra $``$ (here the compactness of $`G`$ is crucial). The spectrum $`\overline{𝒜}`$ of this algebra, that is, the set of all algebraic homomorphisms $`\text{ }\mathrm{C}`$ is called the quantum configuration space. This space is equipped with the Gel’fand topology, that is, the space of continuous functions $`C^0\left(\overline{𝒜}\right)`$ on $`\overline{𝒜}`$ is given by the Gel’fand transforms of elements of $``$. Recall that the Gel’fand transform is given by $`\stackrel{~}{f}\left(\overline{A}\right):=\overline{A}\left(f\right)\overline{A}\overline{𝒜}`$. It is a general result that $`\overline{𝒜}`$ with this topology is a compact Hausdorff space. Obviously, the elements of $`𝒜`$ are contained in $`\overline{𝒜}`$ and one can show that $`𝒜`$ is even dense . Generic elements of $`\overline{𝒜}`$ are, however, distributional.
The idea is now to construct a Hilbert space consisting of square integrable functions on $`\overline{𝒜}`$ with respect to some measure $`\mu `$. Recall that one can define a measure on a locally compact Hausdorff space by prescribing a positive linear functional $`\chi _\mu `$ on the space of continuous functions thereon. The particular measure we choose is given by $`\chi _{\mu _0}\left(\stackrel{~}{T}_I\right)=1`$ if $`I=\{\left\{p\right\},\stackrel{}{0},\stackrel{}{1}\}`$ and $`\chi _{\mu _0}\left(\stackrel{~}{T}_I\right)=0`$ otherwise. Here $`p`$ is any point in $`\mathrm{\Sigma }`$, $`0`$ denotes the trivial representation and $`1`$ the trivial contraction matrix. In other words, (Gel’fand transforms of) spin-network functions play the same role for $`\mu _0`$ as Wick-polynomials do for Gaussian measures and like those they form an orthonormal basis in the Hilbert space $`:=L_2(\overline{𝒜},d\mu _0)`$ obtained by completing their finite linear span $`\mathrm{\Phi }`$.
An equivalent definition of $`\overline{𝒜},\mu _0`$ is as follows :
$`\overline{𝒜}`$ is in one to one correspondence, via the surjective map $`H`$ defined below, with the set $`\overline{𝒜}^{}:=\text{Hom}(𝒳,G)`$ of homomorphisms from the groupoid $`𝒳`$ of composable, holonomically independent, analytical paths into the gauge group. The correspondence is explicitly given by $`\overline{𝒜}\overline{A}H_{\overline{A}}\text{Hom}(𝒳,G)`$ where $`𝒳eH_{\overline{A}}\left(e\right):=\overline{A}\left(h_e\right)=\stackrel{~}{h}_e\left(\overline{A}\right)G`$ and $`\stackrel{~}{h}_e`$ is the Gel’fand transform of the function $`𝒜Ah_e\left(A\right)G`$. Consider now the restriction of $`𝒳`$ to $`𝒳_\gamma `$, the groupoid of composable edges of the graph $`\gamma `$. One can then show that the projective limit of the corresponding cylindrical sets $`\overline{𝒜}_\gamma ^{}:=\text{Hom}(𝒳_\gamma ,G)`$ coincides with $`\overline{𝒜}^{}`$. Moreover, we have $`\left\{\left\{H\left(e\right)\right\}_{eE\left(\gamma \right)};H\overline{𝒜}_\gamma ^{}\right\}=\left\{\left\{H_{\overline{A}}\left(e\right)\right\}_{eE\left(\gamma \right)};\overline{A}\overline{𝒜}\right\}=G^{\left|E\left(\gamma \right)\right|}`$. Let now $`f`$ be a function cylindrical over $`\gamma `$ then
$$\chi _{\mu _0}\left(\stackrel{~}{f}\right)=_{\overline{𝒜}}𝑑\mu _0\left(\overline{A}\right)\stackrel{~}{f}\left(\overline{A}\right)=_{G^{\left|E\left(\gamma \right)\right|}}_{eE\left(\gamma \right)}d\mu _H\left(h_e\right)f_\gamma \left(\left\{h_e\right\}_{eE\left(\gamma \right)}\right)$$
where $`\mu _H`$ is the Haar measure on $`G`$. As usual, $`𝒜`$ turns out to be contained in a measurable subset of $`\overline{𝒜}`$ which has measure zero with respect to $`\mu _0`$.
Let $`\mathrm{\Phi }_\gamma `$, as before, be the finite linear span of spin-network functions over $`\gamma `$ and $`_\gamma `$ its completion with respect to $`\mu _0`$. Clearly, $``$ itself is the completion of the finite linear span $`\mathrm{\Phi }`$ of vectors from the mutually orthogonal $`\stackrel{~}{\mathrm{\Phi }}_\gamma `$. Our basic coordinates of $`M_\gamma `$ are promoted to operators on $``$ with dense domain $`\mathrm{\Phi }`$. As $`h_e`$ is group-valued and $`P^e`$ is real-valued we must check that the adjointness relations coming from these reality conditions as well as the Poisson brackets (2.1) are implemented on our $``$. This turns out to be precisely the case if we choose $`\widehat{h}_e`$ to be a multiplication operator and $`\widehat{P}_j^e=i\mathrm{}\kappa X_j^e/2`$ where $`X_j^e=X_j\left(h_e\right)`$ and $`X^j\left(h\right),hG`$ is the vector field on $`G`$ generating left translations into the $`jth`$ coordinate direction of $`Lie\left(G\right)T_h\left(G\right)`$ (the tangent space of $`G`$ at $`h`$ can be identified with the Lie algebra of $`G`$) and $`\kappa `$ is the coupling constant of the theory. For details see .
## 3 The Heat Kernel Complexifier
The results of this section hold for arbitrary compact, semisimple connected gauge groups and direct products of such with Abelian ones. We will be as explicit as in in order to make this paper self-contained.
As we want to bring in Planck’s constant $`\mathrm{}`$ as a measure of closeness to classical physics, we need to spend a few moments on dimensionalities, see for a general discussion. The dimension of the time coordinate $`x^0`$ is taken to be the same as that of the spatial coordinates $`x^a`$, namely $`\left[x^0\right]=\left[x^a\right]=`$cm<sup>1</sup> which can always be achieved by absorbing an appropriate power of the speed of light into the coupling constant $`\kappa `$ of the theory.
We will take take our connection one-form to be of dimension $`\left[A\right]=`$cm<sup>-1</sup> so that its holonomy is dimensionless. In $`D+1`$ spacetime dimensions the kinetic term of the classical action is given by
$$A_{kin}=\frac{1}{\kappa }_{\text{ }\mathrm{R}}𝑑t_\mathrm{\Sigma }d^DxE_i^a\left(x\right)\dot{A}_a^i\left(x\right)$$
and its dimension is that of an action, that is, $`\left[A_{kin}\right]=\left[\mathrm{}\right]`$. In Yang-Mills theories the electric field is a first derivative of $`A_a^i`$ and thus has dimension $`\left[E_i^a\right]=`$cm<sup>-2</sup>. In general relativity the metric components, the D-beins and also $`\left[E_i^a\right]=`$cm<sup>0</sup> are dimensionfree. It follows that in Yang-Mills (YM) theory the Feinstruktur constant
$$\alpha :=\mathrm{}\kappa $$
(3.1)
has dimension $`\left[\alpha \right]:=`$cm<sup>D-3</sup> and in general relativity (GR) $`\left[\alpha \right]=`$cm<sup>D-1</sup>.
Let now $`\gamma `$ be a graph and consider the symplectic manifold $`(M_\gamma ,\mathrm{\Omega }_\gamma )`$ introduced in section 2.1 with its canonical coordinates $`h_e,P_i^e:eE\left(\gamma \right)`$. The electric flux variable (2.5) then has dimension $`\left[P_i^e\right]=`$cm<sup>D-3</sup> in YM and cm<sup>D-1</sup> in GR respectively and in general let $`\left[P_i^e\right]=`$cm$`^{n_D^{}}`$. Let now $`a`$ be an arbitrary but fixed constant with the dimension of a length, $`\left[a\right]=`$cm<sup>1</sup>, say $`a=1`$cm if $`n_D^{}0`$ and let $`a`$ be dimensionfree otherwise. Then we introduce the dimensionfree quantity
$$p_i^e:=\frac{P_i^e}{a^{n_D}}$$
(3.2)
where $`n_D=n_D^{}`$ if $`n_D^{}0`$ and $`n_D=1`$ otherwise. Notice that a natural choice for a dimensionful constant in general relativity in any $`D`$ would be $`a=1/\sqrt{\left|\mathrm{\Lambda }\right|}`$ where $`\mathrm{\Lambda }`$ is the (supposed to be non-vanishing) cosmological constant.
On the other hand, it is $`E_i^a/\kappa `$ which is canonically conjugate to $`A_a^i`$ rather than $`E_i^a`$ itself, therefore the brackets (2.1) get modified into
$`\{h_e,h_e^{}\}_\gamma `$ $`=`$ $`0`$
$`\{{\displaystyle \frac{P_i^e}{\kappa }},h_e^{}\}_\gamma `$ $`=`$ $`\delta _e^{}^e{\displaystyle \frac{\tau _i}{2}}h_e`$
$`\{{\displaystyle \frac{P_i^e}{\kappa }},{\displaystyle \frac{P_j^e^{}}{\kappa }}\}_\gamma `$ $`=`$ $`\delta ^{ee^{}}f_{ij}^k{\displaystyle \frac{P_k^e}{\kappa }}`$ (3.3)
We are now ready to define the complexifier for the symplectic manifold $`M_\gamma `$, it is given by
$$C_\gamma :=\frac{1}{2\kappa a^{n_D}}\underset{eE\left(\gamma \right)}{}\delta ^{ij}P_i^eP_j^e$$
(3.4)
and since $`C_\gamma `$ is gauge invariant it will pass to the reduced phase space. Using the partial order $``$ of or section 2.1 it is immediately clear that $`C_\gamma `$ defines a self-consistently defined function on the $`M_\gamma `$, that is, for $`\gamma \gamma ^{}`$ we have $`\{p_{\gamma ^{}\gamma }^{}C_\gamma ,p_{\gamma ^{}\gamma }^{}f_\gamma \}_\gamma ^{}=p_{\gamma ^{}\gamma }^{}\{C_\gamma ,f_\gamma \}_\gamma `$ for any $`f_\gamma C^{\mathrm{}}\left(M_\gamma \right)`$.
We can explicitly compute the complexified holonomy and complexified momenta for any compact, semi-simple gauge group $`G`$. Since $`\{P_i^e,C_\gamma \}=0`$ (gauge invariance of $`C_\gamma `$) we have
$`\{h_e,C_\gamma \}_\gamma `$ $`=`$ $`P_i^e{\displaystyle \frac{\tau _i}{2a^{n_D}}}h_e=p_i^e{\displaystyle \frac{\tau _i}{2}}h_e`$
$`\{h_e,C_\gamma \}_{\gamma \left(2\right)}`$ $`=`$ $`{\displaystyle \frac{1}{a^{2n_D}}}P_i^eP_j^e{\displaystyle \frac{\tau _i\tau _j}{4}}h_e=\left(p_j^e{\displaystyle \frac{\tau _j}{2}}\right)^2h_e`$ (3.5)
$`\left(={\displaystyle \frac{p_e^2}{4}}h_e\right)`$
where in the last line we have displayed a simplification that results for $`G=SU\left(2\right)`$ upon using the Clifford algebra relation $`\tau _i\tau _j=\delta _{ij}1_G+f_{ij}^k\tau _k`$ for the Pauli matrices and we define generally $`p^e:=\sqrt{p_j^ep_j^e}`$. In the second line of (3) we have made us of the fact that $`G`$ is semi-simple so that the structure constants are completely skew and so $`\{p_j^e,C_\gamma \}=0`$.
We therefore conclude that the complexification of $`h_e`$ is given by (see for full details)
$`h_e^{\text{ }\mathrm{C}}`$ $`:=`$ $`g_e={\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{i^n}{n!}}\{h_e,C\}_{\left(n\right)}`$ (3.6)
$`=`$ $`\left[{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{i^n}{n!}}\left(p_j^e{\displaystyle \frac{\tau _j}{2}}\right)^n\right]h_e`$
$`=`$ $`e^{i\tau _jp_j^e/2}h_e=:H_eh_e`$
$`\left(=\left[\mathrm{cosh}\left({\displaystyle \frac{p^e}{2}}\right)1_Gi{\displaystyle \frac{p_j^e}{p^e}}\tau _j\mathrm{sinh}\left({\displaystyle \frac{p^e}{2}}\right)\right]h_e\right)`$
and similarly $`P_i^{e\text{ }\mathrm{C}}=P_i^e`$ where we follow the notation of to denote elements of $`G^{\text{ }\mathrm{C}}`$ by $`g`$ while elements of $`G`$ are denoted by $`h`$. In the last line of (3.6) we have again displayed the formula for the special case of $`G=SU\left(2\right)`$. Thus we have established the following.
###### Lemma 3.1
The complexification of the holonomy for compact and semisimple $`G`$ is given directly as a left polar decomposition, where the right unitary factor is the holonomy of the compact gauge group while the left positive definite hermitean factor is just the exponential of $`ip_j^e\tau _j/2`$.
For $`G=U\left(1\right)`$ the generator $`\tau _j/2`$ has to be replaced by the imaginary unit $`i`$.
Notice that (3.6) makes sense since $`p_j^e`$ is dimensionless. Moreover, we have naturally stumbled on the diffeomorphism
$$\mathrm{\Phi }:T^{}\left(G\right)G^{\text{ }\mathrm{C}};(p^j,h)g:=Hh=e^{ip^j\tau _j/2}h.$$
(3.7)
The diffeomorphism (3.7) has a further consequence : $`(T^{}\left(G\right),\omega )`$ is a symplectic manifold while $`G^{\text{ }\mathrm{C}}`$ is a complex manifold. Thus, $`T^{}\left(G\right)`$ is a symplectic manifold with a complex structure which, as one can show ( and references therein), is $`\omega `$-compatible. In fact, $`\omega `$ is just given by (3) with $`P_i^e`$ replaced by $`p_i`$ and the label $`e=e^{}`$ dropped. Therefore, $`T^{}\left(G\right)`$ is in fact a Kähler manifold and a Segal-Bargmann representation (wave functions depending on $`g`$) corresponds to a positive Kähler polarization .
Finally, let us compute the Segal-Bargmann transform corresponding to $`C_\gamma `$ (see for more details). As follows from the previous section, we have in the connection representation (wave functions depending on the $`h_e`$)
$$\widehat{P}_j^e=\frac{i\mathrm{}\kappa }{2}X_j^e\text{ where }X_j^e=X_j\left(h_e\right),$$
(3.8)
and $`X_j\left(h\right)`$ denotes the right invariant vector fields on $`G`$ at $`h`$, that is $`X_j\left(h\right):=\text{tr}\left(\left(\tau _jh\right)^T/h\right)`$. Thus, the coherent state transform is (following the notation of )
$$\widehat{W}_{\gamma t}:=e^{\frac{\widehat{C}_\gamma }{\mathrm{}}}=e^{\frac{t}{2}\mathrm{\Delta }_\gamma }$$
(3.9)
where we have defined the Laplacian on $`\gamma `$ by
$$\mathrm{\Delta }_\gamma =\underset{eE\left(\gamma \right)}{}\mathrm{\Delta }_e,\mathrm{\Delta }_e=\frac{1}{4}\delta ^{ij}X_i^eX_j^e$$
(3.10)
and the heat kernel time parameter has the following interpretation in terms of the fundamental constants of the theory
$$t:=\frac{\mathrm{}\kappa }{a^{n_D}}.$$
(3.11)
Notice that $`a`$ is just a parameter that we have put in by hand to make things dimensionless, for instance, it could be $`1`$cm in quantum general relativity in $`D+1=4`$ spacetime dimensions or $`a=10^5`$ for Yang-Mills in $`D+1=4`$ and thus is “large”. The semiclassical limit $`\mathrm{}0`$ thus corresponds to $`t0`$. That $`t`$ is a tiny positive real number will be crucial in all the estimates that we are going to perform in this and the next paper of this series.
The factor of $`1/4`$ in the definition of $`\mathrm{\Delta }_e`$ relative to $`\left(X_j^e\right)^2`$ is due to the factor of $`1/2`$ in the second Poisson bracket of (3) and it is the same factor which gives $`\mathrm{\Delta }_e`$ the standard spectrum $`j\left(j+1\right);j=0,\frac{1}{2},1,\frac{3}{2},..`$ for the case of $`G=SU\left(2\right)`$.
We can also explicitly compute the quantum operator corresponding to $`g_e`$ in (3.6) for arbitrary $`G`$. We have
$`\widehat{g}_e`$ $`=`$ $`e^{t\mathrm{\Delta }_\gamma /2}\widehat{h}_e^{t\mathrm{\Delta }_\gamma /2}={\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\left(t\right)^n}{2^nn!}}[\widehat{h}_e,\mathrm{\Delta }_e]_{\left(n\right)}`$
$`[\widehat{h}_e,\mathrm{\Delta }_e]`$ $`=`$ $`{\displaystyle \frac{1}{4}}\left(X_e^i\tau _i\widehat{h}_e+\tau _i\widehat{h}_eX_e^i\right)=X_e^i{\displaystyle \frac{\tau _i}{2}}\widehat{h}_e{\displaystyle \frac{\left(\tau _i\right)^2}{4}}\widehat{h}_e`$ (3.12)
$`\left(=\left(X_e^i{\displaystyle \frac{\tau _i}{2}}+{\displaystyle \frac{3}{4}}\right)\widehat{h}_e\right)`$
where the last line is the specialization to $`G=SU\left(2\right)`$. Since $`\mathrm{\Delta }_\gamma `$ commutes with $`X_e^i`$ we immediately find
$`\widehat{g}_e`$ $`=`$ $`e^{t\widehat{X}_e^i\frac{\tau _i}{4}t\frac{\tau _i^2}{8}}\widehat{h}_e=e^{i\widehat{p}_e^j\frac{\tau _j}{2}\frac{t\tau _j^2}{8}}\widehat{h}_e=e^{i\widehat{p}_e^j\frac{\tau _j}{2}}e^{t\frac{\tau _j^2}{8}}\widehat{h}_e`$ (3.13)
$`\left(=e^{\frac{3t}{8}}e^{i\widehat{p}_e^j\frac{\tau _j}{2}}\widehat{h}_e\right)`$
since $`itX_e^j/2=\widehat{p}_j^e`$ and in the third step we used that the matrix $`\tau _j^2`$ commutes with $`\tau _i`$. The last equality holds for $`G=SU\left(2\right)`$ only. Since the $`\widehat{p}_j`$ are not mutually commuting the exponential in (3.13) cannot be defined by the spectral theorem, however, we can define it through Nelson’s analytic vector theorem. Thus, we find precisely the quantization of the polar decomposition (3.6) up to a factor of $`e^{\tau _j^2t/8}`$ which tends to unity linear in $`t0`$. as to be expected. In particular, for $`G=U\left(1\right)`$ we find with $`\tau _j/2`$ replaced by $`i`$
$$\widehat{g}_e=e^{\widehat{p}_e+t/2}\widehat{h}_e=e^{t/2}e^{\widehat{p}_e}\widehat{h}_e$$
(3.14)
Notice that one obtains the first line of (3) from (3.6) if one replaces everywhere $`\{.,.\}`$ by $`[.,.]/\left(i\mathrm{}\right)`$ and phase space functions by operators which holds, of course, by the very construction of the map $`\widehat{W}_t`$ .
## 4 Peakedness Proofs for Gauge-Variant Coherent States
As outlined in the general form of the above transform guarantees immediately that the gauge-variant Coherent States
$$\psi _{\gamma ,\stackrel{}{g}}^t\left(\stackrel{}{h}\right):=\left(\widehat{W}_t\delta _{\gamma \mu _\gamma ,\stackrel{}{h}^{}}\left(\stackrel{}{h}\right)\right)_{|\stackrel{}{h}^{}\stackrel{}{g}}$$
(4.1)
obtained by heat kernel evolution followed by analytic continuation, where $`\stackrel{}{g}=\left\{g_e\right\}_{eE\left(\gamma \right)}`$ and similarly for $`\stackrel{}{h}`$, satisfy a number of desired properties. Here, for completeness we explicitly recall that
$`\delta _{\gamma \mu _\gamma ,\stackrel{}{h^{}}}\left(\stackrel{}{h}\right)={\displaystyle \underset{eE\left(\gamma \right)}{}}\delta _{\mu _H,h_e^{}}\left(h_e\right),`$
$`\delta _{\mu _H,h^{}}\left(h\right)={\displaystyle \underset{\pi }{}}d_\pi \chi _\pi \left(h^{}h^1\right)`$ (4.2)
where $`d\mu _\gamma \left(\stackrel{}{h}\right)=_{eE\left(\gamma \right)}d\mu _H\left(h_e\right)`$ is simply the Haar measure on $`G^E`$, the sum in (4) runs over all distinct irreducible representations $`\pi `$ of $`G`$ (pick once and for all a fixed representant from each equivalence class of those), $`d_\pi =dim\left(\pi \right)`$ is the dimension of the representation space corresponding to $`\pi `$ and $`\chi _\pi (.)=\text{tr}\left(\pi (.)\right)`$ is the character of $`\pi `$ which is a class function and therefore depends only on the equivalence class of $`\pi `$. It follows immediately that therefore the coherent states are explicitly given by
$`\psi _{\gamma ,\stackrel{}{g}}^t\left(\stackrel{}{h}\right)={\displaystyle \underset{eE\left(\gamma \right)}{}}\psi _{g_e}^t\left(h_e\right)`$
$`\psi _g^t\left(h\right)={\displaystyle \underset{\pi }{}}d_\pi e^{\frac{t}{2}\lambda _\pi }\chi _\pi \left(gh^1\right)`$ (4.3)
where $`\lambda _\pi 0,=0`$ only if $`\pi `$ is trivial, is the eigenvalue of the Laplacian in the representation $`\pi `$. For one copy of $`G`$, (4) are precisely the states introduced by Hall who proved various crucial functional analytic properties of these states, in particular that they are entire analytic in $`G^{\text{ }\mathrm{C}}`$ and that heat kernel evolution is densely defined in the Hilbert space $`L_2(G,d\mu _H)`$. Moreover, he proved that the Coherent State Transform
$$\widehat{U}_t:L_2(G,d\mu _H)L_2(G^{\text{ }\mathrm{C}},d\nu _t);f\left(\widehat{U}_tf\right)\left(g\right):=<\overline{\psi _g^t},f>$$
(4.4)
is a unitary transformation between two Hilbert spaces where $`\nu _t`$ is a certain measure to be defined later and $`L_2`$ means a space of square integrable holomorphic functions. This, of course, means that the coherent states so defined satisfy the overcompleteness criterion already.
The product structure of the coherent states, that is, that the coherent state for a graph is just the product over its edges of the coherent states for the edges, is a huge simplification which basically will allow us to reduce all the estimates to just estimates for one copy of $`G`$.
The properties mentioned above are :
* Eigenstates
The coherent states labelled by $`\stackrel{}{g}`$ are simultanous eigenstates for each of the annihilation operators $`\widehat{g}_e^{AB},A,B=1,..,N`$ constructed in the previous section. That is
$$\widehat{g}_e^{AB}\psi _{\gamma ,\stackrel{}{g}}^t=g_e^{AB}\psi _{\gamma ,\stackrel{}{g}}^t$$
(4.5)
* Expectation values
From property (i) it immediately follows that the expectation value of the sum of products of normally ordered functions, that is, the product of any analytic function $`f`$ of the annihilation operators $`\widehat{g}_e^{AB}`$ and any analytic function $`f^{}`$ of the creation operators $`\left(\widehat{g}_e^{AB}\right)^{}`$ in the state $`\psi _{\gamma ,\stackrel{}{g}}^t`$ is given by its classical value at $`\stackrel{}{g},\overline{\stackrel{}{g}}`$. That is,
$$\frac{<\psi _{\gamma ,\stackrel{}{g}}^t,f^{}\left(\stackrel{}{\widehat{g}}^{}\right)f\left(\stackrel{}{\widehat{g}}\right)\psi _{\gamma ,\stackrel{}{g}}^t>}{\psi _{\gamma ,\stackrel{}{g}}^t^2}=f^{}\left(\overline{\stackrel{}{g}}\right)f\left(\stackrel{}{g}\right)$$
(4.6)
* Uncertainty bound
The coherent states automatically saturate, with equal weight (they are unquenched), the uncertainty bound for each pair of self-adjoint operators
$$(\widehat{x}_e^{AB},\widehat{y}_e^{AB}):=(\frac{1}{2}\left(\widehat{g}_e^{AB}+\left(\widehat{g}_e^{AB}\right)^{}\right),\frac{1}{2i}\left(\widehat{g}_e^{AB}\left(\widehat{g}_e^{AB}\right)^{}\right))$$
(4.7)
that is
$$\frac{<\psi _{\gamma ,\stackrel{}{g}}^t,(\widehat{x}_e^{AB}x_e^{AB})^2\psi _{\gamma ,\stackrel{}{g}}^t>}{\psi _{\gamma ,\stackrel{}{g}}^t^2}=\frac{<\psi _{\gamma ,\stackrel{}{g}}^t,(\widehat{y}_e^{AB}y_e^{AB})^2\psi _{\gamma ,\stackrel{}{g}}^t>}{\psi _{\gamma ,\stackrel{}{g}}^t^2}=\frac{1}{2}\frac{|<\psi _{\gamma ,\stackrel{}{g}}^t,[\widehat{x}_e^{AB},\widehat{y}_e^{AB}]\psi _{\gamma ,\stackrel{}{g}}^t>|}{\psi _{\gamma ,\stackrel{}{g}}^t^2}$$
(4.8)
where $`x,y`$ respectively are the expectation values of $`\widehat{x},\widehat{y}`$ respectively. We will compute the actual value of the bound in a later subsection.
These properties are satisfied for any set of coherent states defined by some complexifier $`\widehat{C}`$ which satisfies certain sufficiently strong growth conditions on its eigenvalues (labelled by $`\pi )`$. The peakedness properties that we are after are harder to prove. We will do this in the next subsections for the gauge group $`G=SU\left(2\right)`$. The generalization to an arbirtrary compact gauge group is straightforward but technically difficult and will be displayed in a separate paper . A sketch is contained in section 4.5. In appendix A we also consider the technically much simpler case of $`G=U\left(1\right)`$ and the interested reader is urged to consult that appendix first before looking at the remainder of this section. The graphical supplement to the remaining subsections can be found in Appendix B.
As is obvious from the tensor product structure of our states, it will be completely sufficient to establish the peakedness properties for one copy of $`G`$ only and we can therefore drop the edge lable $`e`$ for the remainder of this section.
### 4.1 Peakedness in the Connection Representation
The coherent states $`\psi _g^t\left(h\right)`$ are defined by the explicit series representation (4) and we are interested in the limit $`t0`$ of the probability distribution (with respect to Haar measure)
$$p_g^t\left(h\right):=\frac{\left|\psi _g^t\left(h\right)\right|^2}{\psi _g^t^2}$$
(4.9)
of which we would like to prove that it is concentrated at $`h=u`$ where $`g=Hu`$ is the polar decomposition of $`gSU\left(2\right)^{\text{ }\mathrm{C}}=SL(2,\text{ }\mathrm{C})`$. As the series in (4) clearly converges worse and worse the smaller $`t`$ gets, the basic tool for all the estimates that follow is the elementary Poisson Summation Formula<sup>1</sup><sup>1</sup>1The authors are indebted to Brian Hall for him pointing out the importance of this formula..
###### Theorem 4.1 (Poisson Summation Formula)
Let $`f`$ be an $`L_1(\text{ }\mathrm{R},dx)`$ function such that the series
$$\varphi \left(y\right)=\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}f\left(y+ns\right)$$
is absolutely and uniformly convergent for $`y[0,s],s>0`$. Then
$$\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}f\left(ns\right)=\frac{2\pi }{s}\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}\stackrel{~}{f}\left(\frac{2\pi n}{s}\right)$$
(4.10)
where $`\stackrel{~}{f}(k):=_{\text{ }\mathrm{R}}\frac{dx}{2\pi }e^{ikx}f(x)`$ is the Fourier transform of $`f`$.
The proof of this theorem can be found in any textbook on Fourier series, see e.g. the classical book by Bochner .
The importance of this remarkable theorem for our purposes is that it converts a slowly converging series $`_nf\left(ns\right)`$ as $`s0`$ into a possibly rapidly converging series $`\frac{1}{s}_n\stackrel{~}{f}\left(2\pi n/s\right)`$ of which in our case almost only the term with $`n=0`$ will be relevant. This is also crucial for numerical approximations as we will see in appendix B.
The way the theorem is stated, it immediately applies to our problem only for the case $`G=U\left(1\right)`$ but one can actually generalize it to any compact gauge group $`G`$ (see e.g. , and references therein). Thus the method of proof displayed below for $`G=SU\left(2\right)`$ can be taken over to the general case.
We begin with the following observation :
$$\psi _g^t\left(h\right)=\psi _H\left(hu^1\right)=\psi _{Huh^1}^t\left(1\right)$$
(4.11)
if $`g=Hu`$ is the polar decomposition of $`g`$. Thus, we see that proving that $`p_g^t\left(h\right)`$ is peaked at $`h=u`$ is equivalent to proving that $`p_H^t\left(h\right)`$ is peaked at $`h=1`$ independently of the positive definite, Hermitean matrix $`H`$. By the same observation and the translation invariance of the Haar measure we see that $`\psi _g^t=\psi _H^t`$. In fact we find
$$\psi _g^t^2=\psi _{H^2}^{2t}\left(1\right)$$
(4.12)
which one shows using the orthogonality relations
$$_G𝑑\mu _H\left(h\right)\overline{\pi \left(h\right)_{mn}}\pi ^{}\left(h\right)_{m^{}n^{}}=\frac{1}{d_\pi }\delta _{\pi \pi ^{}}\delta _{mm^{}}\delta _{nn^{}}$$
(4.13)
(see e.g. , it is also one of the implications of the Peter&Weyl theorem).
So far everything applies to any compact and connected $`G`$. We now specialize to $`G=SU\left(2\right)`$. In this case representations $`\pi _j\left(g\right)_{mn}`$ of dimension $`d_j=2j+1`$ are labelled by half-integral spin quantum numbers $`j=0,\frac{1}{2},1,..`$ and magnetic quantum numbers $`m,n\{j,j+1,..,j\}`$, the eigenvalues of the Laplacian are $`\lambda _j=j\left(j+1\right)`$. In order to compute the character $`\chi _j\left(g\right)=\text{tr}\left(\pi _j\left(g\right)\right),gSL(2,\text{ }\mathrm{C})`$, we need the explicit form of the matrix elements. One finds (see, e.g. )
$$\pi _j\left(g\right)_{mn}=\underset{l}{}\frac{\sqrt{\left(j+m\right)!\left(jm\right)!\left(j+n\right)!\left(jn\right)!}}{\left(jml\right)!\left(j+nl\right)!\left(mn+l\right)!l!}a^{j+nl}d^{jml}b^{mn+l}c^l$$
(4.14)
where the sum extends over all integers for which none of the factorials has negative arguments and
$$g=\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)a,b,c,d\text{ }\mathrm{C},adbc=1.$$
(4.15)
The eigenvalues $`\lambda _1,\lambda _2`$ of $`g`$ follow from the two equations $`det\left(g\right)=\lambda _1\lambda _2=1,\text{tr}\left(g\right)=a+b=\lambda _1+\lambda _2`$ which reveals
$$\lambda _1=\lambda :=x+\sqrt{x^21},\lambda _2=\lambda _1^1=x\sqrt{x^21}\text{ where }x=\frac{a+d}{2}.$$
(4.16)
Since both signs appear in (4.16) there is no ambiguity in taking the square root of the complex number $`x^21`$.
Since the character is a class function invariant under conjugation we can assume $`g`$ to be diagonal in (4.14) in which case the sum over $`l`$ collapses to a single term $`l=0`$ and the sum over $`m`$ becomes a geometric series
$$\chi _j\left(g\right)=\underset{m=j}{\overset{j}{}}a^{j+m}d^{jm}=\underset{m=j}{\overset{j}{}}\lambda ^{2m}=\frac{\lambda ^{2j+1}\lambda ^{\left(2j+1\right)}}{\lambda \lambda ^1}$$
(4.17)
which is invariant under $`\lambda \lambda ^1`$, the action of the Weyl subgroup. Formula (4.17) is a special case of the Weyl character formula .
We can now bring $`\psi _g^t\left(1\right)`$ into a form suitable for the Poisson summation formula
$`\psi _g^t\left(1\right)`$ $`=`$ $`{\displaystyle \underset{j}{}}\left(2j+1\right)e^{\frac{t}{2}j\left(j+1\right)}{\displaystyle \frac{\lambda ^{2j+1}\lambda ^{\left(2j+1\right)}}{\lambda \lambda ^1}}`$ (4.18)
$`=`$ $`{\displaystyle \frac{e^{t/8}}{\lambda \lambda ^1}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}ne^{tn^2/8}\left(\lambda ^n\lambda ^n\right)`$
$`=`$ $`{\displaystyle \frac{e^{t/8}}{\lambda \lambda ^1}}{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}ne^{tn^2/8}\lambda ^n`$
Next we notice that $`\mathrm{ln}\left(\lambda \right)=\text{arcosh}\left(x\right)`$, where the choice of the branch cuts will be defined below, and define $`s:=\sqrt{t}/2,z=\text{arcosh}\left(x\right)/s`$. Then (4.17) can be written as
$$\psi _g^t\left(1\right)=\frac{e^{t/8}}{2s\sqrt{x^21}}\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}\left(ns\right)e^{\left(ns\right)^2/2}e^{\left(ns\right)z}=\frac{e^{t/8}}{2s\sqrt{x^21}}\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}f\left(ns\right)$$
(4.19)
where $`f\left(x\right)=x\mathrm{exp}\left(x^2/2+xz\right)`$. This function certainly satisfies all the conditions for the application of the Poisson summation formula, its Fourier transform is given by
$$\stackrel{~}{f}\left(k\right)=\frac{zik}{\sqrt{2\pi }}e^{\frac{1}{2}\left(k+iz\right)^2}$$
(4.20)
as can be shown by performing a contour integral. Thus we immediately find the desired formula
$`\psi _g^t\left(1\right)`$ $`=`$ $`{\displaystyle \frac{e^{t/8}}{2s\sqrt{x^21}}}{\displaystyle \frac{\sqrt{2\pi }}{s^2}}{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}\left(\text{arcosh}\left(x\right)2\pi in\right)e^{\frac{\left(2\pi n+i\text{arcosh}\left(x\right)\right)^2}{2s^2}}`$ (4.21)
$`=`$ $`{\displaystyle \frac{4\sqrt{2\pi }e^{t/8}}{t^{3/2}}}{\displaystyle \frac{1}{\sqrt{x^21}}}{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}\left(\text{arcosh}\left(x\right)2\pi in\right)e^{2\frac{\left(2\pi n+i\text{arcosh}\left(x\right)\right)^2}{t}}`$
Let likewise $`y=\text{tr}\left(gg^{}\right)/2=\text{tr}\left(H^2\right)/2`$ then
$$\psi _{H^2}^{2t}\left(1\right)=\frac{2\sqrt{\pi }e^{t/4}}{t^{3/2}}\frac{1}{\sqrt{y^21}}\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}\left(\text{arcosh}\left(y\right)2\pi in\right)e^{\frac{\left(2\pi n+i\text{arcosh}\left(y\right)\right)^2}{t}}$$
(4.22)
thus
$$p_H^t\left(h\right)=\frac{\frac{16\sqrt{\pi }}{t^{3/2}}\frac{1}{\left|x^21\right|}\left|_{n=\mathrm{}}^{\mathrm{}}\left(\text{arcosh}\left(x\right)2\pi in\right)e^{2\frac{\left(2\pi n+i\text{arcosh}\left(x\right)\right)^2}{t}}\right|^2}{\frac{1}{\sqrt{y^21}}_{n=\mathrm{}}^{\mathrm{}}\left(\text{arcosh}\left(y\right)2\pi in\right)e^{\frac{\left(2\pi n+i\text{arcosh}\left(y\right)\right)^2}{t}}}$$
(4.23)
Next we observe that upon writing $`H=\mathrm{exp}\left(i\tau _jp^j/2\right)`$ we find $`y=\mathrm{cosh}\left(p\right),p=\sqrt{\left(p^j\right)^2}`$ which allows us to write the probability amplitude as
$$p_H^t\left(h\right)=\frac{\frac{16\sqrt{\pi }}{t^{3/2}}\frac{1}{\left|x^21\right|}\left|_{n=\mathrm{}}^{\mathrm{}}\left(\text{arcosh}\left(x\right)2\pi in\right)e^{2\frac{\left(2\pi n+i\text{arcosh}\left(x\right)\right)^2+p^2/4}{t}}\right|^2}{\frac{1}{\mathrm{sinh}\left(p\right)}_{n=\mathrm{}}^{\mathrm{}}\left(p2\pi in\right)e^{\frac{\left(2\pi n\right)^2+4i\pi np}{t}}}$$
(4.24)
Let us first focus on the denominator $`D_p^t`$ in (4.24) which can be written more explicitly as
$$D_p^t=\frac{p}{\mathrm{sinh}\left(p\right)}\left[1+2\underset{n=1}{\overset{\mathrm{}}{}}e^{4\pi ^2n^2/t}\mathrm{cos}\left(\pi np/t\right)+4\pi \underset{n=1}{\overset{\mathrm{}}{}}ne^{4\pi ^2n^2/t}\frac{\mathrm{sin}\left(4\pi np/t\right)}{p}\right]$$
(4.25)
The term in the square brackets becomes at $`p=0`$ equal to
$$1+2\underset{n=1}{\overset{\mathrm{}}{}}e^{4\pi ^2n^2/t}+16\pi ^2\underset{n=1}{\overset{\mathrm{}}{}}\frac{n^2}{t}e^{4\pi ^2n^2/t}$$
which is still convergent and in fact for $`t0`$ approaches the value $`1`$ exponentially fast with $`t`$. The same is true for $`p0`$ as we show by means of the following lemma.
###### Lemma 4.1
For any complex number $`z`$ we have $`|\mathrm{sin}(z)/z|2\mathrm{cosh}(\mathrm{}(z))<2\mathrm{exp}(|\mathrm{}(z)|)`$.
Proof of Lemma 4.1 :
Let $`z=x+iy`$ then $`\mathrm{sin}\left(z\right)=\mathrm{sin}\left(x\right)\mathrm{cosh}\left(y\right)+i\mathrm{cos}\left(x\right)\mathrm{sinh}\left(y\right)`$. Using that $`\left|x/z\right|,\left|y/z\right|,\left|\mathrm{cos}\left(x\right)\right|1`$ we have
$`\left|{\displaystyle \frac{\mathrm{sin}\left(z\right)}{z}}\right|`$ $``$ $`\left|{\displaystyle \frac{\mathrm{sin}\left(x\right)}{x}}\right|\left|{\displaystyle \frac{x}{z}}\right|\mathrm{cosh}\left(y\right)+\left|\mathrm{cos}\left(x\right)\right|\left|{\displaystyle \frac{\mathrm{sinh}\left(y\right)}{y}}\right|\left|{\displaystyle \frac{y}{z}}\right|`$ (4.26)
$``$ $`{\displaystyle \frac{\mathrm{sin}\left(\left|x\right|\right)}{\left|x\right|}}\mathrm{cosh}\left(y\right)+{\displaystyle \frac{\mathrm{sinh}\left(\left|y\right|\right)}{\left|y\right|}}`$
Now $`\mathrm{sin}\left(x\right)x`$ for all $`x0`$ and employing the Taylor series expansion of $`\mathrm{sinh}\left(y\right)`$ we see that
$$\left|\frac{\mathrm{sinh}\left(y\right)}{y}\right|\underset{n=0}{\overset{\mathrm{}}{}}\frac{y^{2n}}{\left(2n+1\right)!}\underset{n=0}{\overset{\mathrm{}}{}}\frac{y^{2n}}{\left(2n\right)!}=\mathrm{cosh}\left(y\right)$$
(4.27)
which concludes the proof.
$`\mathrm{}`$
With this information at our disposal we can estimate the absolute value of (4.25) as follows
$`\left|D_p^t\right|{\displaystyle \frac{p}{\mathrm{sinh}\left(p\right)}}[1+2{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}e^{4\pi ^2n^2/t}\text{min}_p\left(\mathrm{cos}(4\pi np/t)\right)`$ (4.28)
$`+4\pi {\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}ne^{4\pi ^2n^2/t}\text{min}_p\left({\displaystyle \frac{\mathrm{sin}\left(4\pi np/t\right)}{p}}\right)]`$
$`=`$ $`{\displaystyle \frac{p}{\mathrm{sinh}\left(p\right)}}[12{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}e^{4\pi ^2n^2/t}\text{max}_p\left(\right|\mathrm{cos}(4\pi np/t)\left|\right)`$
$`{\displaystyle \frac{16\pi ^2}{t}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}n^2e^{4\pi ^2n^2/t}\text{max}_p\left(\right|{\displaystyle \frac{\mathrm{sin}\left(4\pi np/t\right)}{4\pi np/t}}\left|\right)]`$
$``$ $`{\displaystyle \frac{p}{\mathrm{sinh}\left(p\right)}}\left[12{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}e^{4\pi ^2n^2/t}{\displaystyle \frac{32\pi ^2}{t}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}n^2e^{4\pi ^2n^2/t}\right]`$
$`=`$ $`{\displaystyle \frac{p}{\mathrm{sinh}\left(p\right)}}\left[1e^{4\pi ^2/t}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}e^{4\pi ^2\left(n^21\right)/t}\left(2+{\displaystyle \frac{32\pi ^2n^2}{t}}\right)\right]`$
$``$ $`{\displaystyle \frac{p}{\mathrm{sinh}\left(p\right)}}\left[1e^{4\pi ^2/t}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}e^{4\pi ^2n^2/t}\left(2+{\displaystyle \frac{32\pi ^2\left(n+1\right)^2}{t}}\right)\right]`$
where in the last step we have used $`\left(n1\right)^2n^21`$ valid for all integers $`n1`$. The series in the last line of (4.28) is certainly still convergent for any $`t>0`$, the dominant term being the one at $`n=0`$ which at $`t0`$ behaves as $`1/t`$. Since $`lim_{t0}e^{a/t}/t^n=0`$ for all $`a>0,n\text{ }\mathrm{Z}`$ we find the first main result.
###### Lemma 4.2
i) There exists a positive constant $`K_t`$ (independent of $`p`$), and exponentially vanishing with $`t0`$ such that $`D_p^t\frac{p}{\mathrm{sinh}(p)}(1K_t)`$ for all $`p0`$.
ii) For the same constant $`K_t`$ it holds that $`D_p^t\frac{p}{\mathrm{sinh}(p)}(1+K_t)`$ for all $`p0`$.
The second part of this lemma is proved by similar methods, one just has to inverse signs in the estimates and replace $`\mathrm{min}\mathrm{max}`$ in the first line of (4.28).
The next step consists in the computation of $`\text{arcosh}\left(x\right)`$. First of all we have with $`h=\mathrm{exp}\left(\theta ^j\tau _j\right),\theta =\sqrt{\left(\theta ^j\right)^2}[0,\pi ],\tau _j=i\sigma _j`$ where $`\sigma _j`$ are the standard Pauli matrices, $`\tau _i\tau _j=\delta _{ij}+ϵ_{ijk}\tau _k`$
$`x={\displaystyle \frac{\text{tr}\left(Hh\right)}{2}}={\displaystyle \frac{\text{tr}\left(\left[\mathrm{cosh}\left(p/2\right)i\frac{p^j}{p}\tau _j\mathrm{sinh}\left(p/2\right)\right]\left[\mathrm{cos}\left(\theta \right)+\frac{\theta ^j}{\theta }\tau _j\mathrm{sin}\left(\theta \right)\right]\right)}{2}}`$ (4.29)
$`=`$ $`\mathrm{cosh}\left(p/2\right)\mathrm{cos}\left(\theta \right)+i\mathrm{sinh}\left(p/2\right)\mathrm{sin}\left(\theta \right)\mathrm{cos}\left(\alpha \right)`$
where $`\mathrm{cos}\left(\alpha \right)=\left(p^j\theta ^j\right)/\left(p\theta \right)[1,1]`$. We wish to write $`x`$ as $`\mathrm{cosh}\left(s+i\varphi \right)`$ for some $`s\text{ }\mathrm{R},\varphi [0,\pi ]`$ and it is a non-trivial question whether this is always possible.
###### Lemma 4.3
For any complex number $`z=R+iI`$ there exist real numbers $`s\text{ }\mathrm{R}`$ and $`\varphi [0,\pi ]`$ such that $`\mathrm{cosh}(s+i\varphi )=z`$. These numbers are uniquely determined except in the case $`I=0,|R|>1`$ in which case the sign of $`s`$ is undetermined.
Proof of Lemma 4.3 :
We will give a constructive proof as we will need the following formulae later on.
We have $`\mathrm{cosh}\left(s+i\varphi \right)=\mathrm{cosh}\left(s\right)\mathrm{cos}\left(\varphi \right)+i\mathrm{sinh}\left(s\right)\mathrm{sin}\left(\varphi \right)`$, thus if the statement of the lemma is true we must have
$$\mathrm{cosh}\left(s\right)\mathrm{cos}\left(\varphi \right)=\mathrm{}\left(z\right)=:R\text{ and }\mathrm{sinh}\left(s\right)\mathrm{sin}\left(\varphi \right)=\mathrm{}\left(z\right)=:I$$
(4.30)
The sign of $`s`$ coincides with that of $`I`$ while $`\varphi \pi /2`$ if $`R0`$ and $`\varphi \pi /2`$ if $`R0`$. Using the trigonometric and hyperbolic relations $`1=\mathrm{cos}^2\left(\varphi \right)+\mathrm{sin}^2\left(\varphi \right)=\mathrm{cosh}^2\left(s\right)\mathrm{sinh}^2\left(s\right)`$ we find after solving a system of quadratic equations unambiguously
$`\mathrm{cosh}^2\left(s\right)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(1+R^2+I^2+\sqrt{\left(1+R^2+I^2\right)^24R^2}\right)`$
$`\mathrm{cos}^2\left(\varphi \right)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(1+R^2+I^2\sqrt{\left(1+R^2+I^2\right)^24R^2}\right)`$ (4.31)
Since $`0\varphi \pi `$ we have $`\mathrm{sin}\left(\varphi \right)0`$ and $`\mathrm{cosh}\left(s\right)1`$ for either sign of $`s`$. Thus the only ambiguity in taking the square root of (4.1) appears in the definition of $`\mathrm{sinh}\left(s\right),\mathrm{cos}\left(\varphi \right)`$. However, in the range of $`s,\varphi `$ that we are considering we find uniquely
$`\mathrm{cosh}\left(s\right)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\sqrt{1+R^2+I^2+\sqrt{\left(1+R^2+I^2\right)^24R^2}}`$
$`\mathrm{sinh}\left(s\right)`$ $`=`$ $`{\displaystyle \frac{\text{sgn}\left(I\right)}{\sqrt{2}}}\sqrt{1+R^2+I^2+\sqrt{\left(1+R^2+I^2\right)^24R^2}}`$
$`\mathrm{cos}\left(\varphi \right)`$ $`=`$ $`{\displaystyle \frac{\text{sgn}\left(R\right)}{\sqrt{2}}}\sqrt{1+R^2+I^2\sqrt{\left(1+R^2+I^2\right)^24R^2}}`$
$`\mathrm{sin}\left(\varphi \right)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\sqrt{1R^2I^2+\sqrt{\left(1+R^2+I^2\right)^24R^2}}`$ (4.32)
where sgn denotes the non-standard step function $`\text{sgn}\left(x\right)=1`$ if $`x0`$ and $`\text{sgn(x)}=1`$ if $`x<0`$. Since the functions $`\mathrm{cos}`$ and $`\mathrm{sinh}`$ respectively are invertible on $`[0,\pi ]`$ and $`\mathrm{R}`$ respectively, above formulae define $`\varphi ,s`$ uniquely. One can explicitly check that the squares of the first and third lines in (4.1) are always greater or smaller than one respectively for any choice of $`R,I`$ and that the arguments of all square roots are non-negative.
We compute
$$\mathrm{cosh}\left(s\right)\mathrm{cos}\left(\varphi \right)+i\mathrm{sinh}\left(s\right)\mathrm{sin}\left(\varphi \right)=\text{sgn}\left(R\right)\left|R\right|+i\text{sgn}\left(I\right)\left|I\right|=z$$
(4.33)
since although sgn is a non-vanishing function, the function $`\text{sgn}\left(x\right)\left|x\right|`$ vanishes anyway at $`x=0`$. This shows that the above choice for $`s,\varphi `$ solves the task to reproduce $`z`$. To see that $`s,\varphi `$ are in fact uniquely determined unless $`\left|R\right|>1,I=0`$ we notice that an ambiguity can possibly arise only through the sign function, that is, if either $`R`$ or $`I`$ vanish.
* $`R=0,I0`$
Then $`\varphi =\pi /2,\text{sgn}\left(s\right)=\text{sgn}\left(I\right)0`$ are uniquely determined.
* $`R0,I=0`$
Then either $`s=0`$ or $`\varphi =0,\pi `$.
Subcase a) : $`\left|R\right|=1`$.
Then $`s=0`$ and $`\varphi =0,\pi `$ if $`R=1,1`$ are uniquely determined.
Subcase b) : $`\left|R\right|<1`$.
If $`s0`$ then necessarily $`\varphi =0,\pi `$ so $`\left|\mathrm{cosh}\left(s\right)\mathrm{cos}\left(\varphi \right)\right|>1>\left|R\right|`$ which is not allowed. Thus $`s`$=0 and $`\varphi `$ is uniquely determined.
Subcase c) : $`\left|R\right|>1`$.
If $`s=0`$ then $`\left|\mathrm{cosh}\left(s\right)\mathrm{cos}\left(\varphi \right)\right|1<\left|R\right|`$ which is not allowed. Thus $`\varphi =0,\pi `$ according to the sign of $`R`$ but the sign of $`s0`$ is ambiguous.
* $`R=I=0`$
Now necessarily $`\varphi =\pi /2,s=0`$ are uniquely determined.
$`\mathrm{}`$
We remark that the undeterminacy of the sign of $`s`$ in the case $`\left|R\right|>1,I=0`$ does not affect us because applied to our situation we have $`R=\mathrm{cosh}\left(p/2\right)\mathrm{cos}\left(\theta \right),I=\mathrm{sinh}\left(p/2\right)\mathrm{sin}\left(\theta \right)\mathrm{cos}\left(\alpha \right)`$ and thus $`\left|R\right|>1`$ means that $`p>0`$ and either $`\theta =0,\pi `$ or $`\alpha =\pi /2`$. In the first case we simply have $`h=\pm 1`$ so $`g=\pm H,x=\pm \mathrm{cosh}\left(p/2\right),\lambda =\pm \mathrm{cosh}\left(p/2\right)+\mathrm{sinh}\left(p/2\right)=\pm \mathrm{exp}\left(\pm p/2\right)`$, thus we fix the signs by $`s:=p/2,\varphi :=\theta =0,\pi `$. In the second case, which is different from the first one only if $`\theta 0,\pi `$, we have $`x=\mathrm{cosh}\left(p/2\right)\mathrm{cos}\left(\theta \right),\lambda =\mathrm{cosh}\left(p/2\right)\mathrm{cos}\left(\theta \right)+\sqrt{\mathrm{cosh}^2\left(p/2\right)\mathrm{cos}^2\left(\theta \right)1}`$ is real-valued because $`\left|R\right|>1`$ and either $`\lambda `$ or $`\lambda `$ is positive if $`\mathrm{cos}\left(\theta \right)`$ is positive or negative respectively. In that case we define $`\varphi =0,\pi `$ respectively and $`s=\text{arcosh}\left(\left|x\right|\right)=\mathrm{ln}\left(\left|x\right|+\sqrt{\left|x\right|^21}\right)>0`$ uniquely so that $`\lambda =e^{s+i\varphi }=\pm e^s`$ is uniquely determined.
Consider now the exponent of the n-th term of the series in the numerator $`N_H^t\left(h\right)`$ of (4.24) given by $`2\left[\left(2\pi n+i\text{arcosh}\left(x\right)\right)^2+p^2/4\right]/t`$ and whose real part is given by $`2\left[\left(2\pi n\varphi \right)^2s^2+p^2/4\right]/t`$. We have the following elementary estimate
$$\left(2\pi n\varphi \right)^2+p^2/4s^24\pi ^2\left(\left|n\right|1\right)^2+\left[\varphi ^2+p^2/4s^2\right]$$
(4.34)
for all $`n0`$ which shows that it is important to know the sign of the function $`p^2/4s^2+\varphi ^2`$. In fact, we would like to show that it is non-negative and vanishing if and only if $`\theta =0`$. The following theorem is the first main theorem of this subsection.
###### Theorem 4.2
For all $`p,\theta ,\alpha `$ it holds that the function
$$\delta ^2(p,\theta ,\alpha ):=p^2/4s^2(p,\theta ,\alpha )+\varphi ^2(p,\theta ,\alpha )\theta ^2$$
(4.35)
is non-negative and zero if and only if either a) $`\varphi =\theta ,|s|=p/2`$ arbitrary and $`|\mathrm{cos}(\alpha )|=1`$ or b) $`\alpha `$ arbitrary and $`s=p=0`$ or $`\varphi =\theta =0,\pi `$.
The proof of this theorem given below is elementary but lengthy, therefore we will break it into several lemmas.
###### Lemma 4.4
The function $`\delta ^2`$ in (4.35) depends on $`\alpha `$ only through $`r:=\mathrm{cos}^2(\alpha )[0,1]`$ and is strictly monotonously decreasing as $`r`$ increases from $`0`$ to $`1`$ for all $`p>0`$ and all $`\theta (0,\pi )`$.
Proof of Lemma 4.4 :
Since $`f`$ depends only on $`\left|s\right|`$ we can determine $`\left|s\right|`$ from $`\mathrm{cosh}\left(s\right)=\mathrm{cosh}\left(\left|s\right|\right)`$ and $`\varphi `$ from $`\mathrm{cosh}\left(\varphi \right)`$. But both formulae in (4.1) depend on $`\alpha `$ only through $`I^2=\mathrm{sinh}^2\left(p\right)\mathrm{sin}^2\left(\theta \right)r`$. Thus, in particular, the sign ambiguity in $`s`$ is irrelevant as far as $`f`$ is concerned.
Lets us define $`\sigma =1+R^2+I^2`$. Then $`\sigma _{,r}=\mathrm{sinh}^2\left(p\right)\mathrm{sin}^2\left(\theta \right)>0`$ for $`p>0,\theta 0,\pi `$. We compute
$`\delta _{,r}^2`$ $`=`$ $`2(\varphi \left[\mathrm{arccos}\left({\displaystyle \frac{\text{sgn}\left(\mathrm{cos}\left(\theta \right)\right)}{\sqrt{2}}}\sqrt{\sigma \sqrt{\sigma ^24R^2}}\right)\right]_{,r}`$ (4.36)
$`\left|s\right|\left[\text{arcosh}\left({\displaystyle \frac{1)}{\sqrt{2}}}\sqrt{\sigma +\sqrt{\sigma ^24R^2}}\right)\right]_{,r})`$
$`=`$ $`2\sigma _{,r}(\varphi {\displaystyle \frac{\text{sgn}\left(\mathrm{cos}\left(\theta \right)\right)}{\sqrt{2}\sqrt{1\frac{1}{2}\left(\sigma \sqrt{\sigma ^24R^2}\right)}}}{\displaystyle \frac{1\frac{\sigma }{\sqrt{\sigma ^24R^2}}}{2\sqrt{\sigma \sqrt{\sigma ^24R^2}}}}`$
$`\left|s\right|{\displaystyle \frac{1}{\sqrt{2}\sqrt{\frac{1}{2}\left(\sigma +\sqrt{\sigma ^24R^2}\right)1}}}{\displaystyle \frac{1+\frac{\sigma }{\sqrt{\sigma ^24R^2}}}{2\sqrt{\sigma +\sqrt{\sigma ^24R^2}}}})`$
$`=`$ $`{\displaystyle \frac{\sigma _{,r}}{\sqrt{\sigma ^24R^2}}}(\varphi {\displaystyle \frac{1}{\sqrt{2\left(\sigma \sqrt{\sigma ^24R^2}\right)}}}\text{sgn}\left(\mathrm{cos}\left(\theta \right)\right)\sqrt{\sigma \sqrt{\sigma ^24R^2}}`$
$`\left|s\right|{\displaystyle \frac{1}{\sqrt{\sigma +\sqrt{\sigma ^24R^2}2}}}\sqrt{\sigma +\sqrt{\sigma ^24R^2}})`$
$`=`$ $`{\displaystyle \frac{\sigma _{,r}}{\sqrt{\sigma ^24R^2}}}(\varphi {\displaystyle \frac{\sqrt{2}\mathrm{cos}\left(\varphi \right)}{\sqrt{2}\sqrt{1\mathrm{cos}^2\left(\varphi \right)}}}|s|{\displaystyle \frac{\sqrt{2}\mathrm{cosh}\left(\left|s\right|\right)}{\sqrt{2}\sqrt{\mathrm{cosh}^2\left(\left|s\right|\right)1}}}`$
$`=`$ $`{\displaystyle \frac{\sigma _{,r}}{\sqrt{\sigma ^24R^2}}}\left(\varphi \text{cot}\left(\varphi \right)\left|s\right|\text{coth}\left(\left|s\right|\right)\right)`$
where in the last step we have observed that $`\left|\mathrm{sin}\left(\varphi \right)\right|=\mathrm{sin}\left(\varphi \right)`$ because $`\varphi [0,\pi ]`$ and that $`\left|\mathrm{sinh}\left(s\right)\right|=\mathrm{sinh}\left(\left|s\right|\right)`$. The last line in (4.36) is evidently non-positive for $`\theta \pi /2`$ (recall that $`\varphi </=/>\pi /2`$ iff $`\theta </=/>\pi /2`$). That this is also true for all of the range of $`\theta `$ follows from the following simple observation.
###### Lemma 4.5
i) The function $`xx\text{cot}(x),x[0,\pi ]`$ is bounded from above by $`1`$ which is reached for $`x=0`$.
ii) The function $`xx\text{coth}(x),x[0,\mathrm{})`$ is bounded from below by $`1`$ which is reached for $`x=0`$.
Proof of Lemma 4.5 :
i)
We simply compute
$$\left(x\text{cot}\left(x\right)\right)^{}=\frac{\mathrm{sin}\left(2x\right)2x}{2\mathrm{sin}^2\left(x\right)}0$$
(4.37)
since $`x\mathrm{sin}\left(x\right)x0`$. Thus the function is monotonously decreasing and therefore its maximum is attained at $`x=0`$ where its value is $`1`$. Notice that the derivative exists even at $`x=0`$.
ii) Likewise we have
$$\left(x\text{coth}\left(x\right)\right)^{}=\frac{\mathrm{sinh}\left(2x\right)2x}{2\mathrm{sinh}^2\left(x\right)}0$$
(4.38)
since $`x\mathrm{sinh}\left(x\right)x0`$. Thus the function is monotonously decreasing and therefore its minimum is attained at $`x=0`$ where its value is $`1`$. Notice that the derivative exists even at $`x=0`$.
$`\mathrm{}`$
Using Lemma 4.5 we conclude that $`\delta _{,r}^20`$ and $`\delta _{,r}^2=0`$ only if either $`\sigma _{,r}=\mathrm{sinh}^2\left(p\right)\mathrm{sin}^2\left(\theta \right)=0`$ or $`\varphi =s=0`$. If $`\varphi =s=0`$ then $`\mathrm{cosh}\left(p/2\right)\mathrm{cos}\left(\theta \right)=1,\mathrm{sinh}\left(p/2\right)\mathrm{sin}\left(\theta \right)\mathrm{cos}\left(\alpha \right)=0`$. Let us exclude the values $`p=0,\theta =0,\pi `$ for the moment, then we find that $`f_{,r}<0`$, except possibly at $`r=0`$ if also $`\mathrm{cosh}\left(p/2\right)\mathrm{cos}\left(\theta \right)=1`$. Thus $`\delta ^2`$ is strictly monotonously decreasing for all $`r>0`$ and since it is a continuous function of $`r`$ it is strictly monotonously decreasing for all $`r[0,1],p>0,\theta (0,\pi )`$.
$`\mathrm{}`$
Proof of Theorem 4.2 :
Using Lemma 4.4 we know that for $`p>0,\theta (0,\pi )`$ the function $`\delta ^2`$ attains its minimum at $`r=1`$ for which we have $`\mathrm{cosh}\left(p/2\right)\mathrm{cos}\left(\theta \right)=\mathrm{cosh}\left(s\right)\mathrm{cos}\left(\varphi \right),\mathrm{sinh}\left(s\right)\mathrm{sin}\left(\varphi \right)=\pm \mathrm{sinh}\left(p/2\right)\mathrm{sin}\left(\theta \right)`$ and thus $`s=\pm p/2,\varphi =\theta `$. Therefore $`\delta ^20`$ for $`p>0,\theta 0,\pi `$ and all $`\alpha [0,\pi ]`$ and $`f=0`$ only at $`\alpha =0,\pi `$.
The remaining cases are $`p=0`$ or $`\theta =0,\pi `$.
Case $`p=0`$ :
Then $`\mathrm{cosh}\left(s\right)\mathrm{cos}\left(\varphi \right)=\mathrm{cos}\left(\theta \right),\mathrm{sinh}\left(s\right)\mathrm{sin}\left(\varphi \right)=0`$. Thus either $`s=0,\varphi =\theta `$ or $`\varphi =0`$ and then $`\mathrm{cosh}\left(s\right)=\mathrm{cos}\left(\theta \right)`$ which is possible only if $`\theta =0,s=0`$. In both of these cases we find $`\delta ^2=0`$.
Case $`\theta =0`$ :
Now $`\mathrm{cosh}\left(s\right)\mathrm{cos}\left(\varphi \right)=\mathrm{cosh}\left(p/2\right),\mathrm{sinh}\left(s\right)\mathrm{sin}\left(\varphi \right)=0`$. Thus either $`\varphi =0,\left|s\right|=p/2`$ or $`s=0`$ and then $`\mathrm{cos}\left(\varphi \right)=\mathrm{cosh}\left(p/2\right)`$ which is possible only if $`p=0,\varphi =0`$. In both cases we find $`\delta ^2=0`$.
Case $`\theta =\pi `$ :
Finally $`\mathrm{cosh}\left(s\right)\mathrm{cos}\left(\varphi \right)=\mathrm{cosh}\left(p/2\right),\mathrm{sinh}\left(s\right)\mathrm{sin}\left(\varphi \right)=0`$. Thus either $`\varphi =\pi ,\left|s\right|=p/2`$ or $`s=0`$ and then $`\mathrm{cos}\left(\varphi \right)=\mathrm{cosh}\left(p/2\right)`$ which is possible only if $`p=0,\varphi =\pi `$. In both cases we find $`\delta ^2=0`$.
Collecting all the results we find $`\delta ^20`$ for all $`p0,\theta [0,\pi ],\alpha [0,\pi ]`$ and $`\delta ^2=0`$ is possible only if either a) $`\alpha =0,\pi `$ while $`\theta ,p`$ can be arbitrary or b) $`p=0`$ or $`\theta =0,\pi `$ while $`\alpha `$ can be arbitrary.
$`\mathrm{}`$
We now come back to the numerator $`N_H^t\left(h\right)`$ in (4.24). The series involved can be transformed into the following expression
$`{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}\left(\text{arcosh}\left(x\right)2\pi in\right)e^{2\frac{\left(2\pi n+i\text{arcosh}\left(x\right)\right)^2+p^2/4}{t}}`$ (4.39)
$`=`$ $`e^{\frac{2}{t}\left(\varphi ^2+p^2/4s^22is\varphi \right)}{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}\left(\text{arcosh}\left(x\right)2\pi in\right)e^{\frac{8\pi ^2n^2}{t}}e^{\frac{8i\pi n}{t}\left(s+i\varphi \right)}`$
$`=`$ $`e^{\frac{2}{t}\left(f+\theta ^22is\varphi \right)}[\text{arcosh}\left(x\right)(1+2{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}e^{\frac{8\pi ^2n^2}{t}}\mathrm{cos}\left({\displaystyle \frac{8\pi n}{t}}\text{arcosh}\left(x\right)\right)`$
$`4\pi {\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}ne^{\frac{8\pi ^2n^2}{t}}\mathrm{sin}\left({\displaystyle \frac{8\pi n}{t}}\text{arcosh}\left(x\right)\right)]`$
$`=`$ $`e^{\frac{2}{t}\left(f+\theta ^22is\varphi \right)}\text{arcosh}\left(x\right)[1+2{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}e^{\frac{8\pi ^2n^2}{t}}\mathrm{cos}\left({\displaystyle \frac{8\pi n}{t}}\text{arcosh}\left(x\right)\right)`$
$`4\pi {\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}ne^{\frac{8\pi ^2n^2}{t}}{\displaystyle \frac{\mathrm{sin}\left(\frac{8\pi n}{t}\text{arcosh}\left(x\right)\right)}{\text{arcosh}\left(x\right)}}]`$
The square bracket expression is certainly regular at $`x=1`$, that is, $`\text{arcosh}\left(x\right)=0`$ and still converges exponentially fast to $`1`$ similarly as for the denominator at $`p=0`$. The same holds at $`\text{arcosh}\left(x\right)0`$. This can be seen as follows : Taking the absolute value of (4.39) we see that it can be estimated from above by (using Lemma 4.1 and Theorem 4.2)
$`\left|\text{4.39}\right|`$ $``$ $`e^{\frac{2}{t}\left(\theta ^2+\delta ^2\right)}\left|\text{arcosh}\left(x\right)\right|[1+2{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}e^{\frac{8\pi ^2n^2}{t}}|\mathrm{cos}\left({\displaystyle \frac{8\pi n}{t}}\text{arcosh}\left(x\right)\right)|`$ (4.40)
$`+32\pi ^2{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{n^2}{t}}e^{\frac{8\pi ^2n^2}{t}}\left|{\displaystyle \frac{\mathrm{sin}\left(\frac{8\pi n}{t}\text{arcosh}\left(x\right)\right)}{8\pi n\text{arcosh}\left(x\right)/t}}\right|]`$
$``$ $`e^{\frac{2}{t}\left(\theta ^2+\delta ^2\right)}\left|\text{arcosh}\left(x\right)\right|[1+2{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}e^{\frac{8\pi ^2n^2}{t}}e^{\frac{8\pi n}{t}\varphi }`$
$`+64\pi ^2{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{n^2}{t}}e^{\frac{8\pi ^2n^2}{t}}e^{8\pi n\varphi /t}]`$
We now consider two cases :
Case (A) : $`0\varphi \left(1c\right)\pi `$ where $`0<c<1`$ will be specified in the course of our derivation.
Case (B) : $`\left(1c\right)\pi \varphi \pi `$.
Turning to Case (A) we can further estimate (4.40) by
$`\left|\text{4.39}\right|`$ $``$ $`e^{\frac{2}{t}\left(\theta ^2+\right)\delta ^2}\left|\text{arcosh}\left(x\right)\right|\left[1+{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}e^{\frac{8\pi ^2n^2}{t}}e^{8\pi ^2\left(1c\right)n/t}\left(2+64\pi ^2{\displaystyle \frac{n^2}{t}}\right)\right]`$ (4.41)
$``$ $`e^{\frac{2}{t}\left(\theta ^2+\delta ^2\right)}\left|\text{arcosh}\left(x\right)\right|\left[1+e^{8\pi ^2c/t}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}e^{\frac{8\pi ^2\left(n^2n\right)}{t}}\left(2+64\pi ^2{\displaystyle \frac{n^2}{t}}\right)\right]`$
$``$ $`e^{\frac{2}{t}\left(\theta ^2+\delta ^2\right)}\left|\text{arcosh}\left(x\right)\right|\left[1+e^{8\pi ^2c/t}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}e^{\frac{8\pi ^2\left(n1\right)^2}{t}}\left(2+64\pi ^2{\displaystyle \frac{n^2}{t}}\right)\right]`$
$``$ $`e^{\frac{2}{t}\left(\theta ^2+\delta ^2\right)}\left|\text{arcosh}\left(x\right)\right|\left[1+e^{8\pi ^2c/t}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}e^{\frac{8\pi ^2n^2}{t}}\left(2+64\pi ^2{\displaystyle \frac{\left(n+1\right)^2}{t}}\right)\right]`$
where in the second step we have dropped the $`n1`$ multiplying $`c`$, in the third we used the estimate $`\left(n1\right)^2n^2n`$ valid for $`n1`$ and in the last we rewrote the series starting at $`n=0`$. The term in the square bracket certainly converges to $`1`$ exponentially fast with $`t0`$ for any $`c>0`$ by an argument already mentioned.
Turning to Case (B) we notice that, as we have to divide the absolute value of the square of (4.40) by $`\left|x^21\right|`$ we need to make sure that $`\text{arcosh}^2\left(x\right)/\left(x^21\right)`$ is bounded at $`x=\pm 1`$. At $`x=1`$ we find $`lim_{x1}\text{arcosh}^2\left(x\right)/\left(x^21\right)=lim_{x1}\text{arcosh}\left(x\right)/\left(x\sqrt{x^21}\right)=lim_{x1}\left(1/\sqrt{x^21}\right)/\left(1/\sqrt{x^21}\right)=1`$ while at $`x\mathrm{}`$ we find $`lim_x\mathrm{}\text{arcosh}^2\left(x\right)/\left(x^21\right)=lim_x\mathrm{}\mathrm{ln}\left(x+\sqrt{x^21}\right)/x=0`$. However at $`x=1`$ we have $`\text{arcosh}\left(x\right)=\pm i\pi `$ and so the expression is in danger to blow up. This is, however not the case. We simply have to write (4.39) in the variable $`\sigma :=\text{arcosh}\left(x\right)i\pi =si\left(\pi \varphi \right)`$ which then becomes
$`{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}\left(\text{arcosh}\left(x\right)2\pi in\right)e^{2\frac{\left(2\pi n+i\text{arcosh}\left(x\right)\right)^2+p^2/4}{t}}`$ (4.42)
$`=`$ $`{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}\left(\sigma i\left(2n1\right)\pi \right)e^{2\frac{\left(\left(2n1\right)\pi +i\sigma \right)^2+p^2/4}{t}}`$
$`=`$ $`e^{\frac{2}{t}\left(p^2/4s^2+\left(\pi \varphi \right)^2+2is\left(\pi \varphi \right)\right)}{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}\left(\sigma i\left(2n1\right)\pi \right)e^{2\frac{\left[\left(2n1\right)\pi \right]^2+2i\left(2n1\right)\sigma \pi }{t}}`$
$`=`$ $`e^{\frac{2}{t}\left(p^2/4s^2+\left(\pi \varphi \right)^2+2is\left(\pi \varphi \right)\right)}[2\sigma {\displaystyle \underset{n=1,\text{odd}}{\overset{\mathrm{}}{}}}e^{2\frac{n^2\pi ^2}{t}}\mathrm{cos}(4n\pi \sigma /t)`$
$`+2\pi {\displaystyle \underset{n=1,\text{odd}}{\overset{\mathrm{}}{}}}ne^{2\frac{n^2\pi ^2}{t}}\mathrm{sin}(4n\pi \sigma /t)]`$
$`=`$ $`e^{\frac{2}{t}\left(p^2/4s^2+\left(\pi \varphi \right)^2+2is\left(\pi \varphi \right)\sigma \right)}\sigma [2{\displaystyle \underset{n=1,\text{odd}}{\overset{\mathrm{}}{}}}e^{2\frac{n^2\pi ^2}{t}}\mathrm{cos}(4n\pi \sigma /t)`$
$`+8\pi ^2{\displaystyle \underset{n=1,\text{odd}}{\overset{\mathrm{}}{}}}{\displaystyle \frac{n^2}{t}}e^{2\frac{n^2\pi ^2}{t}}{\displaystyle \frac{\mathrm{sin}\left(4n\pi \sigma /t\right)}{4n\pi \sigma /t}}]`$
Using again Lemma 4.1, Theorem 4.2 and the fact that $`\left(1c\right)\pi \varphi \pi `$ we can estimate the absolute value of (4.39) as
$`\left|\text{4.39}\right|`$ $``$ $`e^{\frac{2}{t}\left(p^2/4s^2+\left(\pi \varphi \right)^2\right)}\left|\sigma \right|{\displaystyle \underset{n=1,\text{odd}}{\overset{\mathrm{}}{}}}e^{2\frac{n^2\pi ^2}{t}}\left(2+16\pi ^2n^2/t\right)e^{4n\pi \left(\pi \varphi \right)/t}`$ (4.43)
$``$ $`e^{\frac{2}{t}\left(p^2/4s^2+\varphi ^2+\pi ^22\pi \varphi \right)}\left|\sigma \right|{\displaystyle \underset{n=1,\text{odd}}{\overset{\mathrm{}}{}}}e^{2\frac{n^2\pi ^2}{t}}\left(2+16\pi ^2n^2/t\right)e^{4nc\pi ^2/t}`$
$``$ $`e^{\frac{2}{t}\left(\theta ^2+\delta ^2\pi ^2\right)}\left|\sigma \right|{\displaystyle \underset{n=1,\text{odd}}{\overset{\mathrm{}}{}}}e^{2\frac{n^2\left(12c\right)\pi ^2}{t}}\left(2+16\pi ^2n^2/t\right)e^{4c\left(n^2n\right)\pi ^2/t}`$
$``$ $`e^{\frac{2}{t}\left(\theta ^2+\delta ^2\pi ^2+\left(12c\right)\pi ^2\right)}\left|\sigma \right|{\displaystyle \underset{n=1,\text{odd}}{\overset{\mathrm{}}{}}}e^{2\frac{\left(n^21\right)\left(12c\right)\pi ^2}{t}}\left(2+16\pi ^2n^2/t\right)`$
$``$ $`e^{\frac{2}{t}\left(\theta ^2+\delta ^22c\pi ^2\right)}\left|\sigma \right|{\displaystyle \underset{n=1,\text{odd}}{\overset{\mathrm{}}{}}}e^{2\frac{\left(n1\right)^2\left(12c\right)\pi ^2}{t}}\left(2+16\pi ^2n^2/t\right)`$
where in the last line we have assumed $`c<1/2`$ and used the estimate $`\left(n1\right)^2n^21`$ valid for $`n1`$. At this point we choose some $`c<1/2`$ so that $`\theta \pi /2`$. Then we have for any $`0<d<1`$, letting $`n`$ start at $`0`$ in the series,
$`\left|\text{4.39}\right|`$ $``$ $`e^{\frac{2}{t}\left(\left(1d\right)\theta ^2+\delta ^2+\left(d/42c\right)\pi ^2\right)}\left|\sigma \right|{\displaystyle \underset{n=0,\text{even}}{\overset{\mathrm{}}{}}}e^{2\frac{n^2\left(12c\right)\pi ^2}{t}}\left(2+16\pi ^2\left(n+1\right)^2/t\right)`$
$``$ $`\left|\sigma \right|e^{\frac{2}{t}\left(\left(1d\right)\theta ^2+\delta ^2\right)}\left[1+e^{\frac{2}{t}\left(d/42c\right)\pi ^2}{\displaystyle \underset{n=0,\text{even}}{\overset{\mathrm{}}{}}}e^{2\frac{n^2\left(12c\right)\pi ^2}{t}}\left(2+16\pi ^2\left(n+1\right)^2/t\right)\right]`$
where in the second step we have assumed $`2c<d/4`$ in order to isolate the term with $`n=0`$. We see that if we choose $`c<d/8`$ then the term in the square bracket converges exponentially fast to $`1`$ as $`t0`$ and for $`d<1`$ the exponential prefactor decreases exponentially fast to zero as $`t0`$ since $`\theta \pi /2`$.
Let us, for definiteness, choose $`d=1/2,c=1/32`$ which clearly also satisfies $`2c<1`$. Then, putting (4.41) and (4.1) together we have shown :
###### Lemma 4.6
i) For all $`0\varphi 31\pi /32`$ there exists a positive constant $`K_t^{}`$ (independent of $`H,h`$), decaying exponentially fast to $`0`$ as $`t0`$ such that
$$\left|N_H^t\left(h\right)\right|\frac{16\sqrt{\pi }}{t^{3/2}}\frac{\left|\text{arcosh}\left(x\right)\right|^2}{\left|x^21\right|}e^{\frac{4\left(\theta ^2\delta ^2\right)}{t}}\left(1+K_t^{}\right)$$
(4.45)
where $`x=\mathrm{cosh}(p/2)\mathrm{cos}(\theta )+i\mathrm{sinh}(p/2)\mathrm{sin}(\theta )\mathrm{cos}(\alpha )`$.
ii) For all $`31\pi /32\varphi \pi `$ there exists a positive constant $`K_t^{\prime \prime }`$ (independent of $`H,h`$), decaying exponentially fast to $`0`$ as $`t0`$ such that
$$\left|N_H^t\left(h\right)\right|\frac{16\sqrt{\pi }}{t^{3/2}}\frac{\left|\text{arcosh}\left(x\right)i\pi \right|^2}{\left|x^21\right|}e^{\frac{2\left(\theta ^2+2\delta ^2\right)}{t}}\left(1+K_t^{\prime \prime }\right)$$
(4.46)
Finally, combining Lemmata 4.2 and 4.6 we find the following uniform bound for the probability density in position space which is the second main theorem of this subsection.
###### Theorem 4.3
i) For all $`0\varphi 31\pi /32`$ there exist positive constants $`K_t,K_t^{}`$ (independent of $`H,h`$), decaying exponentially fast to $`0`$ as $`t0`$ such that
$$p_H^t\left(h\right)\frac{\frac{16\sqrt{\pi }}{t^{3/2}}\frac{\left|\text{arcosh}\left(x\right)\right|^2}{\left|x^21\right|}e^{\frac{4\left(\theta ^2+\delta ^2\right)}{t}}\left(1+K_t^{}\right)}{\frac{p}{\mathrm{sinh}\left(p\right)}\left(1K_t\right)}$$
(4.47)
where $`x=\mathrm{cosh}(p/)\mathrm{cos}(\theta )+i\mathrm{sinh}(p/)\mathrm{sin}(\theta )\mathrm{cos}(\alpha )`$.
ii) For all $`31\pi /32\varphi \pi `$ there exist positive constants $`K_t^{\prime \prime }`$ (independent of $`H,h`$), decaying exponentially fast to $`0`$ as $`t0`$ such that
$$p_H^t\left(h\right)\frac{\frac{16\sqrt{\pi }}{t^{3/2}}\frac{\left|\text{arcosh}\left(x\right)i\pi \right|^2}{\left|x^21\right|}e^{\frac{2\left(\theta ^2+2\delta ^2\right)}{t}}\left(1+K_t^{\prime \prime }\right)}{\frac{p}{\mathrm{sinh}\left(p\right)}\left(1K_t\right)}$$
(4.48)
Obviously, the bounds are not completely optimal but the remarkable and most important feature is that the bound decays exponentially fast for $`\theta 0`$. At $`\theta =0`$ we have $`\varphi =0`$ and $`s=\pm p/2`$ and the bound is given by the $`p`$-dependent value
$$\frac{\left(p/2\right)^2\mathrm{sinh}\left(p\right)}{p\mathrm{sinh}^2\left(p/2\right)}\frac{16\sqrt{\pi }}{t^{3/2}}\left(1+K_t^{}\right)/\left(1K_t\right)=\frac{p\mathrm{sinh}\left(p\right)}{\mathrm{cosh}\left(p\right)1}\frac{8\sqrt{\pi }}{t^{3/2}}\left(1+K_t^{}\right)/\left(1K_t\right)$$
which defines a rather sharp peak as $`t0`$ and that peak grows linearly with $`p`$. This is in contrast to the harmonic oscillator for which the bound is also Gaussian suppressed in $`xq`$ but it is also independent of $`t,p`$. Of course, this is the effect of the non-Abelian nature of $`G`$ and due to the fact that $`G`$ is not a linear space.
Notice also, that the $`e^{4\delta ^2/t}`$ cannot be dispensed with : For $`\alpha =\theta =\pi /2`$ it can happen that $`s`$ stays bounded while $`p`$ becomes large, in fact $`s=0,\varphi =\pi /2`$ in this case. Since $`\left|\text{arcosh}\left(x\right)\right|^2/\left|x^21\right|=\left[s^2+\varphi ^2\right]/\left[\mathrm{sinh}^2\left(s\right)+\mathrm{sin}^2\left(\varphi \right)\right]`$ we obtain $`\left|\text{arcosh}\left(x\right)\right|^2/\left|x^21\right|=\left(\pi /2\right)^2`$ and it seems that the peak grows exponentially with $`p`$ in this case. However, this is not true : the function $`\delta ^2`$ now takes the value $`p^2/4`$ and so the peak is in fact Gaussian damped with $`p`$ !
### 4.2 Peakedness of the Overlap Function
We compute first the inner product between two coherent states and find
$$<\psi _g^t,\psi _g^{}^t>=\psi _{HH^{}}^{2t}\left(h\right)$$
(4.49)
where $`g=Hu,g^{}=H^{}u^{}`$ are the polar decompositions of $`g,g^{}`$ and $`h=u^1u^{}`$. Our objective is to show that the Overlap Function for these coherent states given by
$$i^t(g,g^{}):=\frac{|<\psi _g^t,\psi _g^{}^t>|^2}{\psi _g^t^2\psi _g^{}^t^2}=\frac{\left[\psi _{HH^{}}^{2t}\left(h\right)\right]^2}{\psi _{H^2}^{2t}\left(1\right)\psi _{\left(H^{}\right)^2}^{2t}\left(1\right)}$$
(4.50)
is peaked at $`g=g^{}`$ which in some sense would mean that the coherent state labelled by $`g`$ represents a neighbourhood (whose size is controlled by $`t`$) of the point $`(p,u)`$ defined by $`g=Hu`$ in the phase space $`T^{}G`$. The existence of a Segal-Bargmann Hilbert space in which wave functions depend on phase space rather than momentum or configuration space will allow us to specify the meaning of this statement precisely in a later subsection.
The idea of proof is to use Theorem 4.3 of the previous subsection. However, in order to do that we must first compute the polar decomposition of $`HH^{}`$ which is not necessarily a Hermitean, positive definite matrix any longer. Using the parameterizations $`H=\mathrm{cosh}\left(p/2\right)i\tau _jp^j\mathrm{sinh}\left(p/2\right)/p,H^{}=\mathrm{cosh}\left(p^{}/2\right)i\tau _jp_j^{}\mathrm{sinh}\left(p^{}/2\right)/p^{}`$ we write $`HH^{}=\stackrel{~}{H}(H,H^{})\stackrel{~}{u}(H,H^{})`$ where $`\stackrel{~}{H}`$ and $`\stackrel{~}{u}`$ are uniquely determined and then have $`\psi _{HH^{}}^{2t}\left(u^1u^{}\right)=\psi _{\stackrel{~}{H}}^{2t}\left(\stackrel{~}{h}\right)`$ where $`\stackrel{~}{h}=u^1u^{}\stackrel{~}{u}^1`$. Suppose then that we can prove that (4.50) is peaked at $`H=H^{}`$ and $`\stackrel{~}{h}=1`$. Then, since $`\stackrel{~}{u}=1`$ and $`\stackrel{~}{H}=H`$ at $`H=H^{}`$, we have automatically shown that $`i^t(g,g^{})`$ is peaked at $`g=g^{}`$. This will be our strategy.
Let then $`\stackrel{~}{H}=\mathrm{exp}\left(i\tau _j\stackrel{~}{p}^j/2\right)`$ and $`\stackrel{~}{h}=\mathrm{exp}\left(\stackrel{~}{\theta }^j\tau _j/2\right)`$. We define as before $`\stackrel{~}{p}=\sqrt{\stackrel{~}{p}^j\stackrel{~}{p}^j},\stackrel{~}{\theta }=\sqrt{\stackrel{~}{\theta }^j\stackrel{~}{\theta }^j},\mathrm{cos}\left(\stackrel{~}{\alpha }\right)=\stackrel{~}{\theta }^j\stackrel{~}{p}^j/\left(\stackrel{~}{\theta }\stackrel{~}{p}\right)`$ and just have to compute $`\stackrel{~}{p}`$ in terms of $`p_j,p_j^{}`$. Defining also $`p=\sqrt{p_jp_j},p^{}=\sqrt{p_j^{}p_j^{}},\mathrm{cos}\left(\beta \right)=p_jp_j^{}/\left(pp^{}\right)`$ we have
$`HH^{}`$ $`=`$ $`\left[\mathrm{cosh}\left(p/2\right)\mathrm{cosh}\left(p^{}/2\right)+\mathrm{cos}\left(\beta \right)\mathrm{sinh}\left(p/2\right)\mathrm{sinh}\left(p^{}/2\right)\right]1_2`$
$`i\tau _j[\mathrm{cosh}(p/2)\mathrm{sinh}(p^{}/2){\displaystyle \frac{p_j^{}}{p^{}}}+\mathrm{cosh}(p^{}/2)\mathrm{sinh}(p/2){\displaystyle \frac{p_j}{p}}`$
$`i{\displaystyle \frac{ϵ_{jkl}p_kp_l^{}}{pp^{}}}\mathrm{sinh}(p/2)\mathrm{sinh}(p^{}/2)]`$
$`\left(HH^{}\right)^{}`$ $`=`$ $`\left[\mathrm{cosh}\left(p/2\right)\mathrm{cosh}\left(p^{}/2\right)+\mathrm{cos}\left(\beta \right)\mathrm{sinh}\left(p/2\right)\mathrm{sinh}\left(p^{}/2\right)\right]1_2`$ (4.51)
$`i\tau _j[\mathrm{cosh}(p/2)\mathrm{sinh}(p^{}/2){\displaystyle \frac{p_j^{}}{p^{}}}+\mathrm{cosh}(p^{}/2)\mathrm{sinh}(p/2){\displaystyle \frac{p_j}{p}}`$
$`+i{\displaystyle \frac{ϵ_{jkl}p_kp_l^{}}{pp^{}}}\mathrm{sinh}(p/2)\mathrm{sinh}(p^{}/2)]`$
Taking the product of these two matrices we find $`\stackrel{~}{H}^2`$ from which we could compute $`\stackrel{~}{p}_j`$ but it turns out that we only need $`\stackrel{~}{p}`$ which we get from the trace
$`\text{tr}\left(HH^{}\left(HH^{}\right)^{}\right)=2[\mathrm{cosh}^2\left(p\right)\mathrm{cosh}^2\left(p^{}\right)+\mathrm{sinh}^2\left(p\right)\mathrm{sinh}^2\left(p^{}\right)+\mathrm{cosh}^2\left(p\right)\mathrm{sinh}^2\left(p^{}\right)`$
$`+\mathrm{cosh}^2\left(p^{}\right)\mathrm{sinh}^2\left(p\right)+4\mathrm{cosh}\left(p\right)\mathrm{cosh}\left(p^{}\right)\mathrm{sinh}\left(p\right)\mathrm{sinh}\left(p^{}\right)\mathrm{cos}\left(\beta \right)]`$ (4.52)
which equals $`2\mathrm{cosh}\left(\stackrel{~}{p}\right)`$. Using hyperbolic identities and addition theorems it is possible to cast (4.2) into the following form
$$\stackrel{~}{p}=\text{arcosh}\left(\left(1+c\right)\mathrm{cosh}^2\left(\frac{p+p^{}}{2}\right)+\left(1c\right)\mathrm{cosh}^2\left(\frac{pp^{}}{2}\right)1\right)$$
(4.53)
where we have used that the $`\mathrm{cosh}`$ function is invertible on the positive real line and $`c=\mathrm{cos}\left(\beta \right)`$ takes values in $`[1,1]`$. The minimum of the argument of (4.53) with respect to $`c`$ at fixed $`p,p^{}`$ is given at $`c=1`$ which is still positive.
Recalling (4.21), (4.22) we find
$`\psi _{\stackrel{~}{H}}^{2t}\left(\stackrel{~}{h}\right)`$ $`=`$ $`{\displaystyle \frac{2\sqrt{\pi }e^{t/4}}{t^{3/2}}}{\displaystyle \frac{1}{\sqrt{x^21}}}{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}\left(\text{arcosh}\left(x\right)2\pi in\right)e^{\frac{\left(2\pi n+i\text{arcosh}\left(x\right)\right)^2}{t}}`$
$`\psi _{H^2}^{2t}\left(1\right)`$ $`=`$ $`{\displaystyle \frac{2\sqrt{\pi }e^{t/4}}{t^{3/2}}}{\displaystyle \frac{1}{\sqrt{y^21}}}{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}\left(\text{arcosh}\left(y\right)2\pi in\right)e^{\frac{\left(2\pi n+i\text{arcosh}\left(y\right)\right)^2}{t}}`$
$`\psi _{\left(H^{}\right)^2}^{2t}\left(1\right)`$ $`=`$ $`{\displaystyle \frac{2\sqrt{\pi }e^{t/4}}{t^{3/2}}}{\displaystyle \frac{1}{\sqrt{\left(y^{}\right)^21}}}{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}\left(\text{arcosh}\left(y^{}\right)2\pi in\right)e^{\frac{\left(2\pi n+i\text{arcosh}\left(y^{}\right)\right)^2}{t}}`$ (4.54)
where $`x=\mathrm{cosh}\left(s+i\varphi \right)=\mathrm{cosh}\left(\stackrel{~}{p}/2\right)\mathrm{cos}\left(\stackrel{~}{\theta }\right)+i\mathrm{sinh}\left(\stackrel{~}{p}/2\right)\mathrm{sin}\left(\stackrel{~}{\theta }\right)\mathrm{cos}\left(\stackrel{~}{\alpha }\right)`$ and $`y=\mathrm{cosh}\left(p\right),y^{}=\mathrm{cosh}\left(p^{}\right)`$. Therefore the overlap function is given by
$`i^t(g,g^{})`$ $`=`$ $`{\displaystyle \frac{1}{\left|x^21\right|}}|{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}(\text{arcosh}\left(x\right)2\pi in)e^{\frac{\left(2\pi n+i\text{arcosh}\left(x\right)\right)^2}{t}}|^2\times `$
$`\times `$ $`{\displaystyle \frac{1}{\frac{1}{\sqrt{y^21}}_{n=\mathrm{}}^{\mathrm{}}\left(\text{arcosh}\left(y\right)2\pi in\right)e^{\frac{\left(2\pi n+i\text{arcosh}\left(y\right)\right)^2}{t}}}}\times `$
$`\times `$ $`{\displaystyle \frac{1}{\frac{1}{\sqrt{\left(y^{}\right)^21}}_{n=\mathrm{}}^{\mathrm{}}\left(\text{arcosh}\left(y^{}\right)2\pi in\right)e^{\frac{\left(2\pi n+i\text{arcosh}\left(y^{}\right)\right)^2}{t}}}}`$
$`=`$ $`e^{\frac{1}{t}\left(p^2+\left(p^{}\right)^2\stackrel{~}{p}^2/2\right)}{\displaystyle \frac{\frac{1}{\left|x^21\right|}\left|_{n=\mathrm{}}^{\mathrm{}}\left(\text{arcosh}\left(x\right)2\pi in\right)e^{\frac{\left(2\pi n+i\text{arcosh}\left(x\right)\right)^2+\stackrel{~}{p}^2/4}{t}}\right|^2}{D_p^tD_p^{}^t}}`$
where $`D_p^t`$ was defined in (4.24). Consider now the exponential in front of the fraction in (4.2).
###### Lemma 4.7
The function
$$\mathrm{\Delta }^2(p,p^{},c):=p^2+\left(p^{}\right)^2\stackrel{~}{p}^2/2$$
(4.56)
is positive definite, vanishing if and only if $`p_j=p_j^{}`$.
Proof of Lemma 4.7 :
Showing that $`f0`$ is equivalent to proving that $`\stackrel{~}{p}\sqrt{2\left[p^2+\left(p^{}\right)^2\right]}`$ or (recall (4.53))
$$\mathrm{cosh}\left(\sqrt{2\left[p^2+\left(p^{}\right)^2\right]}\right)\left(1+c\right)\mathrm{cosh}^2\left(\frac{p+p^{}}{2}\right)+\left(1c\right)\mathrm{cosh}^2\left(\frac{pp^{}}{2}\right)1$$
(4.57)
for any $`p,p^{}0`$ and $`c[1,1]`$. The derivative with respect to $`c`$ of the right hand side of (4.57) is given by $`\mathrm{cosh}^2\left(\frac{p+p^{}}{2}\right)\mathrm{cosh}^2\left(\frac{pp^{}}{2}\right)`$ which is positive unless $`p=p^{}=0`$ in which case the derivative vanishes. However, at $`p=p^{}=0`$ both sides of (4.57) equal $`1`$ so that we are left with the remaining case that not both of $`p,p^{}`$ vanish in which case the right hand side is strictly monotonously increasing with $`c`$. Thus, the right hand side takes its maximum at $`c=1`$. Thus, (4.57) will be true for all $`c`$ given $`p,p^{}`$ if and only if it is true at $`c=1`$ in which case it becomes
$`\mathrm{cosh}\left(\sqrt{2\left[p^2+\left(p^{}\right)^2\right]}\right)2\mathrm{cosh}^2\left({\displaystyle \frac{p+p^{}}{2}}\right)1=\mathrm{cosh}\left(\left(p+p^{}\right)\right)`$ (4.58)
$``$ $`2\left[p^2+\left(p^{}\right)^2\right]\left(p+p^{}\right)^2\left(pp^{}\right)^20`$
Thus, in both cases the inequality is true and becomes an equality only if $`p_j=p_j^{}`$.
$`\mathrm{}`$
Unfortunately it is not possible to prove the more intuitive result $`\mathrm{\Delta }^2\left(p_jp_j^{}\right)^2`$, in fact one can show that the opposite inequality $`\mathrm{\Delta }^2\left(p_jp_j^{}\right)^2`$ holds. Therefore we must live with the function $`\mathrm{\Delta }`$ as a replacement for $`\left(p_jp_j^{}\right)^2`$.
Now consider the remaining factor in (4.2). We see that we can apply Lemma 4.6 to its numerator and Lemma 4.2 to its denominator, the only difference being that we have to replace $`t`$ by $`2t`$ in the final estimate. Therefore we immediately find the main theorem of this subsection.
###### Theorem 4.4
i) For all $`0\varphi 31\pi /32`$ there exist positive constants $`K_t,K_t^{}`$ (independent of $`g,g^{}`$), decaying exponentially fast to $`0`$ as $`t0`$ such that
$$i^t(g,g^{})\frac{\frac{\left|\text{arcosh}\left(x\right)\right|^2}{\left|x^21\right|}e^{\frac{\mathrm{\Delta }^2+2\stackrel{~}{\theta }^2+2\delta ^2}{t}}\left(1+K_t^{}\right)}{\frac{p}{\mathrm{sinh}\left(p\right)}\left(1K_t\right)\frac{p^{}}{\mathrm{sinh}\left(p^{}\right)}\left(1K_t\right)}$$
(4.59)
where $`x=\mathrm{cosh}(\stackrel{~}{p}/2)\mathrm{cos}(\stackrel{~}{\theta })+i\mathrm{sinh}(\stackrel{~}{p}/2)\mathrm{sin}(\stackrel{~}{\theta })\mathrm{cos}(\stackrel{~}{\alpha })`$ and $`\mathrm{\Delta }^2=p^2+(p^{})^2\stackrel{~}{p}^2/2`$.
ii) For all $`31\pi /32\varphi \pi `$ there exist positive constants $`K_tK_t^{\prime \prime }`$ (independent of $`g,g^{}`$), decaying exponentially fast to $`0`$ as $`t0`$ such that
$$i^t(g,g^{})\frac{\frac{\left|\text{arcosh}\left(x\right)i\pi \right|^2}{\left|x^21\right|}e^{\frac{\mathrm{\Delta }^2+\stackrel{~}{\theta }^2+2\delta ^2}{t}}\left(1+K_t^{\prime \prime }\right)}{\frac{p}{\mathrm{sinh}\left(p\right)}\left(1K_t\right)\frac{p^{}}{\mathrm{sinh}\left(p^{}\right)}\left(1K_t\right)}$$
(4.60)
By its very definition, the overlap function is at most unity at $`g=g^{}`$ by the Schwarz inequality and otherwise sharply damped at $`gg^{}`$ as the theorem reveals. In fact, as in the previous section, either $`\delta `$ grows as $`\stackrel{~}{p}^2/4`$ as $`p,p^{}\mathrm{}`$ which leads to Gaussian damping or $`\delta `$ stays bounded in which case $`s`$ grows as $`\stackrel{~}{p}/2`$. In the latter case $`\left|\text{arcosh}\left(x\right)\right|^2/\left|x^21\right|`$ behaves as $`\stackrel{~}{p}^2/\left(4\mathrm{sinh}\left(\stackrel{~}{p}\right)\right)\stackrel{~}{p}^2e^{\stackrel{~}{p}}`$ while the denominator in theorem 4.4 contributes a factor of $`\mathrm{sinh}\left(p\right)\mathrm{sinh}\left(p^{}\right)/\left(pp^{}\right)`$. Now the overlap function is still Gaussian damped due to $`\mathrm{\Delta }0`$ unless $`\stackrel{}{p}=\stackrel{}{p}^{}`$ in which case the two factors just discussed cancel each other as one or both of $`pp^{}`$ get large.
### 4.3 Peakedness in the Electric Field Representation
We first need to define what we even mean by “the electric field representation”.
###### Definition 4.1
i) Let $`|jmn>`$ be the state defined by
$`<h,jmn>:=_Gd\mu _H(h^{})\delta (h,h^{})|jmn>(h^{})=<\delta _h,jmn>=\pi _j(h)_{mn}`$ and let $`\psi L_2(G,d\mu _H)`$ be any state. Then we define the electric field representation of $`\psi `$ by
$$\stackrel{~}{\psi }\left(jmn\right):=<jmn,\psi >$$
(4.61)
that is, the electric field representation of $`\psi `$ is nothing else than its “Fourier coefficients” with respect to the complete orthogonal system $`|jmn>`$ normalized by $`|jmn>^2=1/d_j`$.
ii) The Peter&Weyl theorem guarantees that $`\psi \stackrel{~}{\psi }`$ is a unitary transformation between $`L_2(G,d\mu _H)`$ and the Hilbert space $`\mathrm{}_2`$ of sequences $`(c_{jmn})`$ of complex numbers equipped with the inner product $`<c,c^{}>=_{jmn}d_j\overline{c_{jmn}}c_{jmn}^{}`$.
We have defined the electric field representation for a general compact group $`G`$ where $`j`$ is some discrete label for a complete system of representants from each equivalence class of irreducible representations and $`m,n`$ labels its matrix elements. For $`SU\left(2\right)`$ $`j`$ is a half integral non-negative integer and $`m,n`$ take the $`d_j=2j+1`$ values $`j,j+1,..,j`$.
We easily calculate that
$$\stackrel{~}{\psi }_g^t\left(jmn\right)=e^{tj\left(j+1\right)/2}\pi _j\left(g\right)_{mn}$$
(4.62)
and are interested in the probability amplitude
$$p_g^t\left(jmn\right):=\frac{\left|\stackrel{~}{\psi }_g^t\left(jmn\right)\right|^2}{\psi _g^t^2}$$
(4.63)
for the momentum of the particle to be in the configuration $`jmn`$ in the state $`\psi _g^t`$. The precise relation between the classical numbers $`p^j`$ and the quantum numbers $`jmn`$ will become clear shortly.
We notice the following elementary estimates : Let $`g=Hu=u\left(u^1Hu\right)=uH^{}`$ be the unique left and right polar decompositions of $`g`$. Define $`X_m^{}:=\pi _j\left(H\right)_{mm^{}},\overline{Y_m^{}}:=\pi _j\left(u\right)_{m^{}n},X_m^{}^{}:=\pi _j\left(H^{}\right)_{m^{}n},\overline{Y_m^{}^{}}:=\pi _j\left(u\right)_{mm^{}}`$, then by the Schwarz inequality
$`\left|\pi _j\left(g\right)_{mn}\right|`$ $`=`$ $`\left|{\displaystyle \underset{m^{}}{}}\overline{Y_m^{}}X_m^{}\right|=|<Y,X>|\left|\right|Y\left|\right|\left|\right|X\left|\right|`$
$`\left|\pi _j\left(g\right)_{mn}\right|`$ $`=`$ $`\left|{\displaystyle \underset{m^{}}{}}\overline{Y_m^{}^{}}X_m^{}^{}\right|=|<Y^{},X^{}>|\left|\right|Y^{}\left|\right|\left|\right|X^{}\left|\right|`$ (4.64)
where the inner product is the Hermitean inner product of the $`d_j`$ dimensional representation space corresponding to $`\pi _j`$. But $`Y^2=Y^{}^2=1`$ by the unitarity of $`u`$ while $`X^2=\pi _j\left(H^2\right)_{mm}`$ and $`X^{}^2=\pi _j\left(\left(H^{}\right)^2\right)_{nn}`$ by the hermiticity of $`H,H^{}`$. We summarize this observation in the following Lemma.
###### Lemma 4.8
The matrix elements of $`\pi _j(g)_{mn}`$ have the factorizing bound
$$\left|\pi _j\left(g\right)_{mn}\right|^2\sqrt{\pi _j\left(H^2\right)_{mm}}\sqrt{\pi _j\left(\left(H^{}\right)^2\right)_{nn}}$$
(4.65)
for all $`jm,nj`$ where $`H,H^{}`$ are the left and right polar decompositions of $`g=Hu=uH^{}`$.
This factorization property will be crucial later on when we project the gauge-variant coherent states on a general graph to the gauge invariant subspace of the Hilbert space.
In the considerations that follow we will again specialize to $`G=SU\left(2\right)`$. The following Lemma, recalling (3.14), justifies the name “electric field representation”.
###### Lemma 4.9
Let $`{}_{}{}^{R}p_{e}^{j}:=p_e^j`$ as in section 3 and $`{}_{}{}^{L}p_{e}^{j}:=\frac{1}{2}\text{tr}(h_e\tau _jh_e^1\tau _k)^Rp_e^j`$ (recall (3.4)). Then, dropping the label $`e`$,
$`{}_{}{}^{R}\widehat{p}_{3}^{}|jmn>`$ $`=`$ $`itm|jmn>`$
$`{}_{}{}^{L}\widehat{p}_{3}^{}|jmn>`$ $`=`$ $`itn|jmn>`$
$`(^R\widehat{p}_j)^2|jmn>`$ $`=`$ $`(^L\widehat{p}_j)^2|jmn>=+t^2j(j+1)|jmn>`$ (4.66)
that is, the three operators $`{}_{}{}^{R}\widehat{p}_{3}^{},^L\widehat{p}_3,(^R\widehat{p}_j)^2`$ are simultanuously diagonizable with $`|jmn>`$ as eigenstates. Moreover, the magnetic quantum numbers $`mt,nt`$ have the interpretation of the 3-component of $`{}_{}{}^{R}p_{j}^{}`$ and $`{}_{}{}^{L}p_{j}^{}`$ respectively while for large $`p`$ the quantum number $`jt`$ has the interpretation of the norm of $`{}_{}{}^{R}p_{j}^{}`$ which equals the norm of $`{}_{}{}^{L}p_{j}^{}`$.
Proof of Lemma 4.9 :
The proof follows almost immediately from the fact that $`{}_{}{}^{R}\widehat{p}_{j}^{}=it^RX_j/2,^L\widehat{p}_j=it^LX_j/2`$ where $`{}_{}{}^{R}X,^LX`$ denotes the right or left invariant vector field on $`G`$ which certainly commute with each other and their square gives four times the Laplacian. The eigenvalues displayed can be easily computed from the fact that $`{}_{}{}^{R}X_{j}^{}=\frac{d}{ds}_{|s=0}L_{\mathrm{exp}\left(s\tau _j\right)}`$ and $`{}_{}{}^{L}X_{j}^{}=\frac{d}{ds}_{|s=0}R_{\mathrm{exp}\left(s\tau _j\right)}`$ where $`R_h,L_h`$ denote right and left translation on $`G`$ and from the explicit matrix element formula (4.14) by expanding $`\pi _j\left(\mathrm{exp}\left(s\tau _j\right)g\right)_{mn},\pi _j\left(g\mathrm{exp}\left(s\tau _j\right)\right)_{mn}`$ around $`s=0`$.
$`\mathrm{}`$
We need the following lemma.
###### Lemma 4.10
The functions $`p_g^t(jmn)`$ are bounded as $`j,m,n\mathrm{}`$ with a peak at $`(j+1/2)t=p`$ for all $`m,n`$.
Proof of Lemma 4.10 :
Recalling Lemma 4.2 we have first of all
$$p_g^t\left(jmn\right)\frac{\mathrm{sinh}\left(p\right)}{p}\frac{t^{3/2}e^{t/4}}{2\sqrt{\pi }}\frac{e^{p^2/t}}{1K_t}e^{tj\left(j+1\right)}\left|\pi _j\left(g\right)_{mn}\right|^2$$
(4.67)
for some positive constant $`K_t`$ decaying exponentially to zero as $`t0`$. We have the elementary estimate
$$\left|\pi _j\left(g\right)_{mn}\right|^2\underset{mn}{}\left|\pi _j\left(g\right)_{mn}\right|^2=\chi _j\left(gg^{}\right)=\chi _j\left(H^2\right)=\frac{\mathrm{sinh}\left(\left(2j+1\right)p\right)}{\mathrm{sinh}\left(p\right)}$$
(4.68)
and therefore after simple algebraic manipulations
$$p_g^t\left(jmn\right)\frac{1}{p}\frac{t^{3/2}}{4\sqrt{\pi }}\frac{1}{1K_t}e^{\frac{\left[\left(j+1/2\right)tp\right]^2}{t}}$$
(4.69)
for all $`p1`$, say, and, using again Lemma 4.1, we find
$$p_g^t\left(jmn\right)\left(2j+1\right)\frac{t^{3/2}}{2\sqrt{\pi }}\frac{1}{1K_t}e^{\frac{\left[\left(j+1/2\right)tp\right]^2}{t}}$$
(4.70)
for all $`0p1`$. From these estimates peakedness is obvious at the value claimed.
$`\mathrm{}`$
Up to now all estimates were for general $`p`$. From now on we restrict attention to large $`p`$ (that is, of order unity or larger) as it is of interest in applications to semi-classical approximations. As the probability amplitude is then small, according to the previous lemma, unless $`jtp`$, we can restrict attention to the case that $`j`$ is large in what follows.
The next theorem is the main result of this subsection.
###### Theorem 4.5
The diagonal matrix elements $`\pi _j(H)_{mm},\pi _j(H^{})_{nn}`$ are for large $`p,p^{},j`$ peaked at $`m/j=p_3/p`$ and $`n/j=p_3^{}/p^{}`$ where $`H=\mathrm{exp}(ip_j\tau _j/2),H^{}=\mathrm{exp}(ip_j^{}\tau _j/2)`$. The maximal value of $`\pi _j(H)_{mm}`$ at $`m[p_3/p]j`$ is given by $`e^{pj}`$.
Proof of Theorem 4.5 :
We display the proof for $`H`$, the one for $`H^{}`$ is identical.
We will discuss separately the following two cases :
Case I) $`\left|p_3/p\right|<1`$ :
Employing the explicit formula (4.14) at $`m=n`$ we find, using $`adbc=1`$
$$\pi _j\left(H\right)_{mm}=\left(ad\right)^j\left(\frac{a}{d}\right)^m\underset{l}{}\left(\begin{array}{c}j+m\\ l\end{array}\right)\left(\begin{array}{c}jm\\ l\end{array}\right)\left[1\frac{1}{ad}\right]^l$$
(4.71)
where, as usual, the sum over $`l`$ is over all integers such that no factorials have negative arguments. Since $`a=\mathrm{cosh}\left(p/2\right)+\mathrm{sinh}\left(p/2\right)p_3/p,d=\mathrm{cosh}\left(p/2\right)\mathrm{sinh}\left(p/2\right)p_3/p`$ we have $`ad=\mathrm{cosh}^2\left(p/2\right)\mathrm{sinh}^2\left(p/2\right)\left(p_3/p\right)^2`$ which is large if $`p`$ is large unless $`p_3\pm p`$ which we excluded.
For large $`p`$ we can therefore replace $`11/\left(ad\right)`$ by $`1`$ and can use the addition theorem for binomial coefficients
$$\underset{l}{}\left(\begin{array}{c}\alpha \\ l\end{array}\right)\left(\begin{array}{c}\beta \\ \gamma l\end{array}\right)=\left(\begin{array}{c}\alpha +\beta \\ \gamma \end{array}\right)$$
to arrive at
$$\pi _j\left(H\right)_{mm}\left(ad\right)^j\left(\frac{a}{d}\right)^m\left(\begin{array}{c}2j\\ jm\end{array}\right)$$
(4.72)
For large $`j`$ to which we have focussed attention to, and if also $`j\pm m`$ are large, more precisely, if $`\left|m/j\right|<1`$, we can apply the crudest version of Stirling’s formula $`n!\left(n/e\right)^n`$ to estimate the factorials. Introducing the abbreviations $`s=p_3/p,t=m/j`$ we have $`ad=\mathrm{cosh}^2\left(p/2\right)s^2\mathrm{sinh}^2\left(p/2\right)\mathrm{exp}\left(p\right)\left(1s^2\right)/4,a/d\left(1+s\right)/\left(1s\right)`$, thus
$`\pi _j\left(H\right)_{mm}`$ $``$ $`\left(ad\right)^j\left({\displaystyle \frac{2j}{e}}\right)^{2j}\left(a/d\right)^m{\displaystyle \frac{e^{2j}}{\left(j+m\right)^{j+m}\left(jm\right)^{jm}}}`$ (4.73)
$``$ $`{\displaystyle \frac{e^{pj}}{4^j}}\left(1s^2\right)^j\left(2j\right)^{2j}\left({\displaystyle \frac{1+s}{1s}}\right)^m{\displaystyle \frac{1}{\left(1+t\right)^{j+m}\left(1t\right)^{jm}j^{2j}}}`$
$`=`$ $`e^{pj}\left(1s^2\right)^j\left({\displaystyle \frac{1+s}{1s}}\right)^{jt}{\displaystyle \frac{1}{\left(1+t\right)^{j\left(1+t\right)}\left(1t\right)^{j\left(1t\right)}}}`$
$`=`$ $`e^{pj}\left({\displaystyle \frac{1s^2}{1t^2}}\right)^j\left({\displaystyle \frac{\left(1+s\right)\left(1t\right)}{\left(1s\right)\left(1+t\right)}}\right)^{jt}`$
$`=:`$ $`e^{pj}\left(1s^2\right)^je^{jf\left(t\right)}`$
Let us compute the extrema of the function $`f\left(t\right)`$. We have
$$\dot{f}=\frac{2t}{1t^2}+\mathrm{ln}\left(\frac{\left(1+s\right)\left(1t\right)}{\left(1s\right)\left(1+t\right)}\right)t\left(\frac{1}{1t}+\frac{1}{1+t}\right)=\mathrm{ln}\left(\frac{\left(1+s\right)\left(1t\right)}{\left(1s\right)\left(1+t\right)}\right)$$
(4.74)
which vanishes precisely at $`t=s`$. Moreover, $`\frac{d^2f}{dt^2}=\frac{2}{1t^2}<0`$ for all $`t[0,1]`$, thus $`t=s`$ is the only local and therefore the global maximum. We conclude that $`f\left(t\right)f\left(s\right)`$ and thus $`\pi _j\left(H\right)_{mm}e^{jp}`$ where the maximum is taken at $`m/j=p_3/p`$. Notice that our intermediate assumption that $`\left|m/j\right|<1`$ is justified in retrospect as well. Expanding $`f`$ around $`t=s`$ we get $`f\left(t\right)=f\left(s\right)+f^{\prime \prime }\left(s\right)\left(ts\right)^2+o\left(\left(ts\right)^3\right)=\mathrm{ln}\left(1s^2\right)\frac{\left(ts\right)^2}{1s^2}`$ so that
$$\pi _j\left(H\right)_{mm}e^{2jp}e^{j\frac{\left(m/jp_3/p\right)^2}{1\left(p_3/p\right)^2}}$$
(4.75)
Case II) $`\left|p_3/p\right|1`$ :
In this case $`p_1/p=p_2/p0`$ and the sum over $`l`$ in (4.71) collapses to a single term $`l=0`$ and we find
$$\pi _j\left(H\right)_{mm}\left(ad\right)^j\left(\frac{a}{d}\right)^m$$
(4.76)
Since $`a/d=\left(\mathrm{cosh}\left(p/2\right)+s\mathrm{sinh}\left(p/2\right)\right)/\left(\mathrm{cosh}\left(p/2\right)s\mathrm{sinh}\left(p/2\right)\right)=\mathrm{exp}\left(sp\right)`$ for $`s\pm 1`$ while $`ad=\mathrm{cosh}^2\left(p/2\right)s^2\mathrm{sinh}^2\left(p/2\right)1`$ we get
$$\pi _j\left(H\right)_{mm}e^{smp}$$
(4.77)
which obviously takes its maximum at $`m=sj`$, that is, $`t=m/j=s`$ again. The maximum value is given by $`\pi _j\left(H\right)_{mm}=e^{jp}`$. Thus
$$\pi _j\left(H\right)_{mm}e^{jp}e^{jp\left(sm/j1\right)}e^{jp}e^{jp\left|m/jp_3/p\right|}$$
(4.78)
$`\mathrm{}`$
Notice that at $`m=p_3/pj`$ we have $`\pi _j\left(H\right)_{mm}e^{2pj}`$ while we have shown already that $`\left|\pi _j\left(g\right)_{mn}\right|\sqrt{\chi _j\left(H^2\right)}e^{\left(j+1/2\right)p}`$. This means that for large $`p,j`$ the function $`\left|\pi _j\left(H\right)_{mn}\right|\sqrt{\pi _j\left(H^2\right)_{mm}}e^{pj}`$ is indeed concentrated at $`m=njp_3/p`$. This can be shown explicitly by repeating the above analysis and varying besides $`m`$ also $`n`$.
We summarize the results of this subsection in the following theorem.
###### Theorem 4.6
The probability amplitude $`p_g^t(jmn)`$ is, for large $`p`$, peaked at $`jtp`$ and $`mtp_3^R,ntp_3^L`$. More precisely, there exists a constant $`K_t`$ exponentially vanishing as $`t0`$ and independent of $`g`$ such that
$`p_g^t\left(jmn\right)`$ $`\stackrel{<}{}`$ $`{\displaystyle \frac{1}{p}}{\displaystyle \frac{t^{3/2}}{4\sqrt{\pi }}}{\displaystyle \frac{1}{1K_t}}e^{j/2\frac{(m/j(^Rp_3)/p)^2}{1(^Rp_3/p)^2}}e^{j/2\frac{(n/j(^Lp_3)/p)^2}{1(^Lp_3/p)^2}}e^{\frac{\left[\left(j+1/2\right)tp\right]^2}{t}}`$
$`\text{if }|^{R/L}p_3/p|<1`$
$`p_g^t\left(jmn\right)`$ $`\stackrel{<}{}`$ $`{\displaystyle \frac{1}{p}}{\displaystyle \frac{t^{3/2}}{4\sqrt{\pi }}}{\displaystyle \frac{1}{1K_t}}e^{jp|m/j(^Rp)_3/p|}e^{jp|n/j(^Lp)_3/p|}e^{\frac{\left[\left(j+1/2\right)tp\right]^2}{t}}`$ (4.79)
$`\text{if }|^{R/L}p_3/p|\stackrel{<}{}1`$
The careful reader will notice a seemingly crucial difference between the configuration and momentum representation : While the peak in the configuration representation grows with $`t0`$, in the momentum representation it sinks with $`t0`$. However, this is only apparently so : notice that the configuration Hilbert space is an $`L_2`$ space since the operators $`\widehat{h}_{AB}`$ have continuous spectrum while the momentum Hilbert space is an $`\mathrm{}_2`$ space since the operators $`{}_{}{}^{R}\widehat{p}_{j}^{},^L\widehat{p}_j,\widehat{p}_j^2`$ have discrete spectrum. Now let $`\xi _g^t=\psi _g^t/\psi _g^t`$ then
$$1=\xi _g^t^2=\underset{jmn}{}p_g^t\left(jmn\right)=\underset{jmn}{}\mathrm{\Delta }p_j\mathrm{\Delta }^Rp_m\mathrm{\Delta }^Lp_n\frac{p_g^t\left(jmn\right)}{t^3/2}$$
(4.80)
where $`p_j=jt,^Rp_m=mt,^Lp_n=nt`$ with $`2j\text{ }\mathrm{N},jm,jn\text{ }\mathrm{Z},\left|m\right|,\left|n\right|j`$ and therefore $`2\mathrm{\Delta }p_j=\mathrm{\Delta }^Rp_m=\mathrm{\Delta }^Lp_n=t`$. It follows from (4.6) that $`p_g^t\left(jmn\right)=\stackrel{~}{p}_g^t\left(p_j,^Rp_m,^Lp_n\right)`$ evidently depends only on $`p_j,^Rp_m,^Lp_n`$ and thus (4.80) is a Riemann sum, as $`t0`$, approximating the integral
$$1=_0^{\mathrm{}}𝑑p_j_{p_j}^{p_j}d^Rp_m_{p_j}^{p_j}d^Lp_m\frac{\stackrel{~}{p}_g^t\left(p_j,^Rp_m,^Lp_n\right)}{2t^3}$$
(4.81)
In other words, as $`t0`$ the momentum spectrum approaches a continuum one and the corresponding propability amplitude is up to a constant factor given by (4.6) divided by $`t^3`$ whose peak evidently also grows with $`t0`$ just as in the configuration representation. Thus, the apparent difference of the peak behaviour for the two representations is absent in the limit $`t0`$ if we use a contiuum spectrum approximation.
The trick to use an approximate continuum momentum representation as in (4.80), (4.81) will be used in the proof of Ehrenfest theorems in . In particular, we see from the explicit expression (4.6) that $`\stackrel{~}{p}_g^t\left(p_j,^Rp_m,^Lp_n\right)\delta (p\left(g\right),jt)\delta (^Rp\left(g\right),mt)\delta (^Lp\left(g\right),nt)`$ approaches a $`\delta `$ distribution with respect to the measure (4.81).
### 4.4 Uncertainty Relation and Phase Space Bounds
In this subsection we will compute explictly the Heisenberg uncertainty bound for the operators $`\widehat{g}_{AB}`$, verify that it corresponds to the bound to be expected from the Poisson bracket $`\{g_{AB},\overline{g_{AB}}\}`$ and finally will see explicitly that the overlap function $`i^t(g,g^{})`$ times $`1/t^3`$ can be interpreted as the probability density to find the system at the phase space point $`g^{}`$ in the state $`\widehat{U}_t\psi _g^t`$ with respect to the Liouville measure on phase space.
We will first need the so-called averaged heat kernel measure $`\nu _t`$ on $`G^{\text{ }\mathrm{C}}`$ which one can obtain either by the methods derived in (and advertized in ) which are specific for the heat kernel coherent states or by the more general method derived in for an arbitrary family of coherent states. We will give a direct derivation below for $`SU\left(2\right)`$ as we wish to be as explicit as possible.
###### Lemma 4.11
The measure $`\nu _t`$ underlying the map defined in (4.4) is given for $`G=SU(2)`$ by
$$d\nu _t\left(g\right):=d\mu _H\left(u\right)d\sigma _t\left(H\right):=d\mu _H\left(u\right)\left[\frac{2\sqrt{2}e^{t/4}}{\left(2\pi t\right)^{3/2}}\frac{\mathrm{sinh}\left(p\right)}{p}e^{p^2/t}d^3p\right]=\nu _t\left(g\right)d\mathrm{\Omega }$$
(4.82)
where $`g=Hu`$ is the polar decomposition of $`g`$, $`d^3p`$ is the standard Lebesgue measure on $`\text{ }\mathrm{R}^3`$ and $`d\mathrm{\Omega }=d\mu _Hd^3p`$ is the Liouville measure on $`T^{}G`$.
Proof of Lemma 4.11 :
First of all, to see that $`d\mathrm{\Omega }\left(g\right):=d\mu _H\left(u\right)d^3p`$ with $`g=Hu,H=\mathrm{exp}\left(i\tau _gp_j/2\right)`$ is the Liouville measure on $`G^{\text{ }\mathrm{C}}T^{}G`$ for the case $`G=SU\left(2\right)`$ (up to normalization) it is heplful to think of $`SU\left(2\right)`$ as the sphere $`S^3`$. The phase space $`\widehat{N}=T^{}G`$ can then be thought of as the symplectic reduction of the phase space $`N=T^{}\text{ }\mathrm{R}^4`$ under the co-isotropic constraint $`C:=\left(x^1\right)^2+\left(x^2\right)^2+\left(x^3\right)^2+\left(x^4\right)^21`$. Writing the symplectic structure $`\omega =_{I=1}^4dP_Idx^I`$ on $`N`$ in terms of radial and polar coordinates defined by
$$\stackrel{}{x}=:r\stackrel{}{n}:=r(\mathrm{sin}\left(\theta \right)\mathrm{sin}\left(\varphi \right)\mathrm{cos}\left(\phi \right),\mathrm{sin}\left(\theta \right)\mathrm{sin}\left(\varphi \right)\mathrm{sin}\left(\phi \right),\mathrm{sin}\left(\theta \right)\mathrm{cos}\left(\varphi \right),\mathrm{cos}\left(\theta \right))$$
with $`r[0,\mathrm{});\theta ,\varphi [0,\pi ];\phi [\pi ,\pi ]`$ as well as adapted normal and tangential (to $`S^3`$) momenta defined by $`P_I^{}:=\left(P_Jn^J\right)n^I,P_I^{}:=P_IP_I^{}`$ repectively (the latter of which are Dirac observables) and then pulling it back to the constraint surface $`C=0`$ gives the Liouville measure on $`T^{}S^3\times \text{ }\mathrm{R}^3`$ which is a product measure on $`S^3`$ times the Lebesgue measure on $`\text{ }\mathrm{R}^3`$. The same measure on $`S^3`$ can be obtained as the effective measure induced by $`\delta \left(C\right)d^4x`$ which is obviously proportional to the Haar measure on $`SU\left(2\right)`$ as it is invariant under $`SO\left(4\right)SU\left(2\right)\times SU\left(2\right)`$ (i.e. left and right translations).
We now must verify the isometry relation
$$<\widehat{U}_t\psi ,\widehat{U}_t\psi ^{}>_{\nu _t}=<\psi ,\psi ^{}>_{\mu _H}$$
(4.83)
for any two $`\psi ,\psi ^{}L_2(G,d\mu _H)`$. It will be sufficient to check this on a basis, say the basis $`|jmn>`$ introduced in the previous subsection for which $`\left(\widehat{U}_t|jmn>\right)\left(g\right)=e^{tj\left(j+1\right)/2}\pi _j\left(g\right)_{mn}`$. Using the polar decomposition and writing $`\pi _j\left(g\right)=\pi _j\left(H\right)\pi _j\left(u\right)`$ we see that we can take advantage of the orthogonality relations of the $`\pi _j\left(u\right)_{mn}`$ under Haar measure if we make the ansatz $`d\nu _t\left(g\right)=d\mu _H\left(u\right)d\sigma _t\left(H\right)`$. Thus, choosing $`\psi =|jmn>,\psi ^{}=|j^{}m^{}n^{}>`$, we immediately find the condition
$$\delta _{mm^{}}e^{tj\left(j+1\right)}=𝑑\sigma _t\left(H\right)\pi _j\left(H^2\right)_{m^{}m}$$
(4.84)
for all $`j,m,m^{}`$. We can produce the required Kronecker $`\delta `$ on the right hand side of (4.84) if we choose the measure $`d\sigma _t\left(H\right)`$ to be invariant under $`SO\left(3\right)`$, the homomorphic image of $`SU\left(2\right)`$ under the vector (or spin $`1`$) representation because in that case
$`\pi _j\left(u\right)_{nm^{}}\left[{\displaystyle 𝑑\sigma _t\left(H\right)\pi _j\left(H^2\right)_{m^{}m}}\right]=\left[{\displaystyle 𝑑\sigma _t\left(H\right)\pi _j\left(\left(uHu^1\right)^2\right)_{nm^{}}}\right]\pi _j\left(u\right)_{m^{}m}`$
$`=`$ $`\left[{\displaystyle 𝑑\sigma _t\left(H\right)\pi _j\left(H^2\right)_{nm^{}}}\right]\pi _j\left(u\right)_{m^{}m}`$
that is, the matrix $`A_{m^{}m}=𝑑\sigma _t\left(H\right)\pi _j\left(H^2\right)_{m^{}m}`$ commutes with the irreducible representation $`\pi _j`$ of $`SU\left(2\right)`$ and is therefore proportional to the unit matrix by one of Schur’s lemmata. We therefore are led to the ansatz $`d\sigma _t\left(H\right)=d^3pf_t\left(p\right)`$ where the positive function $`f_t\left(p\right)`$ only depends on $`p=\sqrt{p_jp_j}`$ and $`H=\mathrm{exp}\left(ip_j\tau _j/2\right)`$ as before. Then $`A_{m^{}m}=\text{tr}\left(A\right)\delta _{mm^{}}/d_j`$ and we are left with the condition that
$$\left(2j+1\right)e^{tj\left(j+1\right)}=_{\text{ }\mathrm{R}^3}𝑑\sigma _t\left(H\right)\chi _j\left(H^2\right)=4\pi _0^{\mathrm{}}p^2𝑑pf_t\left(p\right)\frac{\mathrm{sinh}\left(\left(2j+1\right)p\right)}{\mathrm{sinh}\left(p\right)}$$
(4.85)
for all $`j`$. We see that we can produce the $`\left(2j+1\right)`$ factor on the right hand side of (4.85) if we can do an integration by parts. We therefore write $`f_t\left(p\right)=g_t^{}\left(p\right)\mathrm{sinh}\left(p\right)/p^2`$ and find the condition
$$e^{tj\left(j+1\right)}=4\pi _0^{\mathrm{}}𝑑pg_t\left(p\right)\mathrm{cosh}\left(\left(2j+1\right)p\right)$$
(4.86)
provided that $`g_t\left(p\right)`$ is finite at $`p=0`$ where $`\mathrm{sinh}\left(\left(2j+1\right)p\right)`$ vanishes and that $`g_t\left(p\right)`$ decays faster at infinity than $`\mathrm{sinh}\left(\left(2j+1\right)p\right)`$. Finally, assuming that $`g_t\left(p\right)=g_t\left(p\right)`$ is invariant under reflection we find the condition
$$e^{\frac{t}{4}\left(2j+1\right)^2t/4}=2\pi _{\text{ }\mathrm{R}}𝑑pg_t\left(p\right)\mathrm{exp}\left(\left(2j+1\right)p\right)$$
(4.87)
which we recognize as the moment problem for a Gaussian. Thus we define $`g_t\left(p\right)=k_t\mathrm{exp}\left(sp^2/2\right)`$ and find
$$e^{\frac{t}{4}\left(2j+1\right)^2t/4}=2\pi k_t/\sqrt{s}_{\text{ }\mathrm{R}}𝑑xe^{x^2/2}\mathrm{exp}\left(\left(2j+1\right)x/\sqrt{s}\right)=\sqrt{2\pi }^3k_t/\sqrt{s}e^{\frac{\left(2j+1\right)^2}{2s}}$$
(4.88)
from which we read off $`s=2/t,k_t=\frac{e^{t/4}\sqrt{2/t}}{\sqrt{2\pi }^3}`$. Notice that $`g_t`$ is indeed finite at $`p=0`$, decays faster than any exponential of $`p`$ at $`\mathrm{}`$ and is reflection invariant. Therefore
$`d\sigma _t\left(H\right)`$ $`=`$ $`f_t\left(p\right)d^3p=\mathrm{sinh}\left(p\right)/p^2g_t^{}\left(p\right)d^3p=k_t\mathrm{sinh}\left(p\right)/p^2\left(2pe^{p^2/t}/t\right)d^3p`$ (4.89)
$`=`$ $`{\displaystyle \frac{\mathrm{sinh}\left(p\right)}{p}}{\displaystyle \frac{2\sqrt{2}e^{t/4}}{\left(2\pi t\right)^{3/2}}}e^{p^2/t}d^3p`$
$`\mathrm{}`$
The next Lemma is sometimes called the reproducing kernel property and holds completely generally for any system of coherent states defined by a complexifier . We will state and prove it only for the group case for general $`G`$ (see for more details).
###### Lemma 4.12
The coherent state transform of a coherent state at the value $`g^{}`$ is the same as taking inner products in the $`L_2`$ Hilbert space with the coherent state with label $`(g^{})^{}`$ where $`g^{}=(g^1)^{}`$ is the unique involution on $`G^{\text{ }\mathrm{C}}`$ with the property that $`g^{}=g`$ if and only if $`gG`$. That is,
$$\left(\widehat{U}_t\psi _g^t\right)\left(g^{}\right)=\psi _g^{2t}\left(g^{}\right)=<\psi _{\left(g^{}\right)^{}}^t,\psi _g^t>_{\mu _H}$$
(4.90)
Proof of Lemma 4.12 :
The proof is trivial. We have
$$<\psi _g^{}^t,\psi _g^t>=\underset{\pi }{}d_\pi e^{t\lambda _\pi }\chi _\pi \left(g\left(g^{}\right)^{}\right)$$
(4.91)
while
$$\left(\widehat{U}_t\psi _g^t\right)\left(g^{}\right)=\underset{\pi }{}d_\pi e^{t\lambda _\pi }\chi _\pi \left(g\left(g^{}\right)^1\right)=\psi _g^{2t}\left(g^{}\right)$$
(4.92)
$`\mathrm{}`$
Notice that in the polar decomposition of $`g=Hu`$ we have $`g^{}=H^1u`$ which corresponds to $`p_jp_j,uu`$ as it should be.
The following theorem clarifies the meaning of the overlap function of subsection 4.2. We will do this here only for $`SU\left(2\right)`$. The statement for general $`G`$ can be found in .
###### Theorem 4.7
The overlap function $`i^t(g,g^{})`$ approaches exponentially fast with $`t0`$ the function $`p^t(g,(g^{})^{})\frac{\pi t^3}{2}`$ where $`p^t(g,g^{})`$ denotes the probability density to find the system at the phase space point $`g^{}`$ in the state $`\widehat{U}_t\psi _g^t`$ with respect to the Liouville measure on the phase space $`T^{}G`$.
Proof of Theorem 4.7 :
The probability density of the image of the normalized coherent state $`\psi _g^t`$ under the coherent state transform at the phase space point $`g^{}`$ in the Bargmann-Segal Hilbert space with respect to the Liouville measure $`d\mathrm{\Omega }=d\mu _H\left(u\right)d^3p`$ on $`T^{}G`$ is given by
$$p^t(g,g^{})=\nu _t\left(g^{}\right)\frac{\left|\left(\widehat{U}_t\psi _g^t\right)\left(g^{}\right)\right|^2}{\psi _g^t^2}$$
(4.93)
where we have used the isometry property of $`\widehat{U}_t`$, that is, the norm in the denominator of (4.93) can be computed in either Hilbert space. Using Lemma 4.12 and the definition of the overlap function we have
$`p^t(g,g^{})`$ $`=`$ $`\left[\nu _t\left(g^{}\right)\psi _{\left(g^{}\right)^{}}^t^2\right]{\displaystyle \frac{|<\psi _{\left(g^{}\right)^{}}^t,\psi _g^t>|^2}{\psi _g^t^2\psi _{\left(g^{}\right)^{}}^t^2}}`$ (4.94)
$`=`$ $`\left[\nu _t\left(g^{}\right)\psi _g^{}^t^2\right]i^t(g,\left(g^{}\right)^{})`$
where we have used the fact that $`\psi _g^t`$ depends on $`p`$ only. Now using Lemma 4.2 and the explicit expression for $`\nu _t\left(g^{}\right)`$ given in (4.82) we find for the factor multiplying $`i^t(g,\left(g^{}\right)^{})`$ in (4.94) the estimate
$$\frac{2}{\pi t^3}\left(1K_t\right)\frac{p^t(g,g^{})}{i^t(g,\left(g^{}\right)^{})}\frac{2}{\pi t^3}\left(1+\stackrel{~}{K}_t\right)$$
(4.95)
for some constants $`K_t,\stackrel{~}{K}_t`$, independent of $`g,g^{}`$, exponentially vanishing with $`t0`$.
$`\mathrm{}`$
Since $`i^t(g,g^{})`$ is peaked at $`g=g^{}`$ where it equals unity and has a decay width of order $`\sqrt{t}`$ we conclude that like for a particle moving in $`\text{ }\mathrm{R}^3`$ the phase space volume occupied by a coherent state with respect to the measure $`d\mu _H\left(h\right)d^3p`$ is given by $`\left(2\pi t\right)^3\mathrm{}^3`$. In particular, we obtain the interpretation that the normalized coherent states with label $`g`$ in the Bargmann-Segal Hilbert space are concentrated at the phase space point $`g^{}`$ with respect to the Liouville measure. If we would have defined the map $`\widehat{U}_t`$ through heat kernel evolution followed by antianalytical extension, i.e. $`\left(\widehat{U}_t\psi \right)\left(g\right):=\left(\widehat{W}_t\psi \right)\left(h\right)_{hg^{}}`$ (which for $`gG`$ does not make any difference) then the coherent state labelled by $`g`$ in the Bargmann-Segal Hilbert space is concentrated at $`g`$ since the measure $`d\nu _t\left(g\right)=d\nu _t\left(g^{}\right)`$ is involution invariant and so the unitarity and peakedness properties are preserved. With this definition of $`\widehat{U}_t`$ the strange asymmetry $`gg^{}`$ is removed and we assume this to have been done from now on.
Let us now compute the actual uncertainty bound. By the Heisenberg uncertainty relation for the self-adjoint operators
$$\widehat{x}_{AB}:=\frac{1}{2}\left(\widehat{g}_{AB}+\left(\widehat{g}_{AB}\right)^{}\right),\widehat{y}_{AB}:=\frac{1}{2i}\left(\widehat{g}_{AB}\left(\widehat{g}_{AB}\right)^{}\right)$$
(4.96)
where $`(.)^{}`$ means the adjoint with respect to $`L_2(G,d\mu _H)`$ we have for any $`A,B\{1/2,1/2\}`$ and for any state
$$<\left(\mathrm{\Delta }\widehat{x}_{AB}\right)^2>^{1/2}<\left(\mathrm{\Delta }\widehat{y}_{AB}\right)^2>^{1/2}\frac{\left|<[\widehat{g}_{AB},\widehat{g}_{AB}^{}]>\right|}{4}$$
(4.97)
and for coherent states the bound is saturated with equal contributions from $`\widehat{x},\widehat{y}`$. From this we conclude easily that
$$\frac{\left|_{A,B}<[\widehat{g}_{AB},\widehat{g}_{AB}^{}]>\right|}{4}\underset{A,B}{}\frac{\left|<[\widehat{g}_{AB},\widehat{g}_{AB}^{}]>\right|}{4}4\text{max}_{A,B}<\left(\mathrm{\Delta }\widehat{x}_{AB}\right)^2>^{1/2}<\left(\mathrm{\Delta }\widehat{y}_{AB}\right)^2>^{1/2}$$
(4.98)
We will compute the quantity on the left hand side of (4.98) instead of the individual bounds as this is easier and because it gives a uniform bound. It also gives an idea of the individual bounds because they are all of the same order as one can easily check.
We begin with the computation of the Poisson brackets (remember that $`g=Hh`$)
$`\{g_{AB},\overline{g_{AB}}\}`$ $`=`$ $`{\displaystyle \frac{H_{AC}}{p_j}}\{p_j,h_{BD}^1\}h_{CB}H_{DA}{\displaystyle \frac{H_{DA}}{p_j}}\{p_j,h_{CB}\}h_{BD}^1H_{AC}`$ (4.99)
$`+\{p_j,p_k\}{\displaystyle \frac{H_{AC}}{p_j}}h_{CB}{\displaystyle \frac{H_{DA}}{p_k}}h_{BD}^1`$
and using the symplectic structure given in (3.7) we find
$`{\displaystyle \underset{A,B}{}}\{g_{AB},\overline{g_{AB}}\}`$ $`=`$ $`{\displaystyle \frac{\kappa }{a}}{\displaystyle \frac{}{p_j}}\text{tr}\left(H^2\tau _j\right)`$ (4.100)
$`=`$ $`2i{\displaystyle \frac{\kappa }{a}}\left(\mathrm{cosh}\left(p\right)+2{\displaystyle \frac{\mathrm{sinh}\left(p\right)}{p}}\right)`$
where we have made use of $`H^2=\mathrm{cosh}\left(p\right)i\tau _j\frac{p_j}{p}\mathrm{sinh}\left(p\right)`$. Notice that the right hand side of (4.100) depends on the phase space point which is different from the situation with $`T^{}\text{ }\mathrm{R}`$.
We now compare this with the expectation value of the sum of commutators with respect to the normalized coherent state $`\psi _g^t`$. In order to do this we need the following Lemma about the Clebsch-Gordan decomposition.
###### Lemma 4.13
For any $`gSL(2,\text{ }\mathrm{C})`$ we have
$`g_{A_0B_0}\pi _j\left(g\right)_{A_1..A_{2j},B_1..B_{2j}}`$ $`=`$ $`\pi _{j+1/2}\left(g\right)_{A_0..A_{2j},B_0..B_{2j}}`$ (4.101)
$`{\displaystyle \frac{d_{j1/2}}{d_j}}ϵ_{A_0(A_1}\pi _{j1/2}\left(g\right)_{A_2..A_{2j}),(B_2..B_{2j}}ϵ_{B_1)B_0}`$
where all A’s and B’s take the values $`\pm 1/2`$, $`ϵ_{AB}`$ is the totally skew tensor density of weight minus one in two dimensions and $`(.)`$ denotes total symmetrization of indices to be taken as an idempotent operation.
The proof requires elementary linear algebra and is left to the reader. One uses the fact that the space of totally symmetric spinors of rank $`2j`$ provide the representation space of the irreducible representation of spin $`j`$ of $`SU\left(2\right)`$, that is, $`\pi _j\left(g\right)_{A_1..A_{2j},B_1..B_{2j}}=\pi _j\left(g\right)_{A_1+..+A_{2j},B_1+..+B_{2j}}`$ in terms of the former notation with magnetic quantum numbers.
We now use the fact that
$$\widehat{g}_{AB}=e^{t\mathrm{\Delta }/2}\widehat{h}_{AB}e^{t\mathrm{\Delta }/2},\left(\widehat{g}_{AB}\right)^{}=e^{t\mathrm{\Delta }/2}\left(\widehat{h}^1\right)_{BA}e^{t\mathrm{\Delta }/2}$$
(4.102)
and that $`g^1=ϵg^Tϵ^1`$ for any $`gSL(2,\text{ }\mathrm{C})`$. The computations are rather tedious and lengthy. We will not display all the steps but only the main stations of the calculation which require frequent use of Lemma 4.13 and relabelling of indices. One first checks that indeed
$$\left(\widehat{g}_{AB}\psi _g^t\right)\left(h\right)=g_{AB}\psi _g^t\left(h\right)$$
(4.103)
so that we have easily
$$<\psi _g^t,\underset{A,B}{}\left(\widehat{g}_{AB}\right)^{}\widehat{g}_{AB}\psi _g^t>=\text{tr}\left(g^{}g\right)\psi _g^t^2$$
(4.104)
where the dagger in the last line denotes the matrix adjoint. Using the $`SL(2,\text{ }\mathrm{C})`$ Mandelstam identity $`\text{tr}\left(g\right)\chi _j\left(g\right)=\chi _{j+1/2}\left(g\right)+\chi _{j1/2}\left(g\right)`$ which one derives from Lemma 4.13 one can write (4.104) in the equivalent form
$$<\psi _g^t,\underset{A,B}{}\left(\widehat{g}_{AB}\right)^{}\widehat{g}_{AB}\psi _g^t>=\underset{j}{}d_je^{t\lambda _j}\left(\chi _{j+1/2}\left(H^2\right)+\chi _{j1/2}\left(H^2\right)\right)$$
(4.105)
On the other hand, tedious calculations reveal that
$$<\psi _g^t,\underset{A,B}{}\widehat{g}_{AB}\left(\widehat{g}_{AB}\right)^{}\psi _g^t>=\underset{j}{}d_je^{t\lambda _j}\chi _j\left(H^2\right)\left[\frac{d_{j+1/2}}{d_j}e^{t\left(\lambda _j\lambda _{j+1/2}\right)}\frac{d_{j1/2}}{d_j}e^{t\left(\lambda _j\lambda _{j1/2}\right)}\right]$$
(4.106)
Taking the difference of (4.106) and (4.105) and relabelling summation indices one arrives at
$`<\psi _g^t,{\displaystyle \underset{A,B}{}}[\widehat{g}_{AB},\left(\widehat{g}_{AB}\right)^{}]\psi _g^t>`$ (4.107)
$`=`$ $`2{\displaystyle \underset{j}{}}e^{t\lambda _j}\chi _j\left(H^2\right)\left[d_{j1/2}\mathrm{sinh}\left(t\left(j+1/4\right)\right)d_{j+1/2}\mathrm{sinh}\left(t\left(j+3/4\right)\right)\right]`$
Introducing the parameter $`T=\sqrt{t}/2`$ and the function
$`f\left(x\right)=\mathrm{exp}\left(x^2\right)\left(xT\right)\mathrm{sinh}\left(px/T\right)\mathrm{sinh}\left(2TxT^2\right)`$ one can cast (4.107) in a form suitable for an appeal to the Poisson summation formula
$$<\psi _g^t,\underset{A,B}{}[\widehat{g}_{AB},\left(\widehat{g}_{AB}\right)^{}]\psi _g^t>=2\frac{e^{t/4}}{T\mathrm{sinh}\left(p\right)}\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}f\left(nT\right)$$
(4.108)
Computing the Fourier transform of $`f`$ and applying the Poisson summation formula we end up with
$`<\psi _g^t,{\displaystyle \underset{A,B}{}}[\widehat{g}_{AB},\left(\widehat{g}_{AB}\right)^{}]\psi _g^t>={\displaystyle \frac{4\sqrt{\pi }e^{t/4}}{t^{3/2}\mathrm{sinh}\left(p\right)}}{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}\times `$
$`\times `$ $`\{(p/2i\pi n)e^{t/4}e^{4\left(\pi n+i\left(p/2+T^2\right)\right)^2/t}(p/2+i\pi n+2T^2)e^{t/4}e^{4\left(\pi ni\left(p/2+T^2\right)\right)^2/t}`$
$`+(p/2+i\pi n+2T^2)e^{t/4}e^{4\left(\pi n+i\left(p/2T^2\right)\right)^2/t}+(p/2+i\pi n)e^{t/4}e^{4\left(\pi ni\left(p/2T^2\right)\right)^2/t}\}`$
Recalling (4.22) that
$$\psi _g^t^2=\frac{4\sqrt{\pi }e^{t/4}}{t^{3/2}}\frac{1}{\mathrm{sinh}\left(p\right)}\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}\left(p/2i\pi n\right)e^{4\frac{\left(\pi n+ip/2\right)^2}{t}}$$
(4.110)
we see that the prefactors in front of the sums in (4.107) and (4.110) equal each other and an analysis similar to that which has led to Lemma 4.2 reveals that there exist positive constants $`K_t,K_t^{},\stackrel{~}{K}_t,\stackrel{~}{K}_t^{}`$ exponentially vanishing with $`t0`$ such that
$`2{\displaystyle \frac{\mathrm{cosh}\left(p\right)\left(1e^{t/2}\right)t\frac{\mathrm{sinh}\left(p\right)}{p}\stackrel{~}{K}_t^{}}{1+K_t^{}}}`$ (4.111)
$``$ $`{\displaystyle \frac{<\psi _g^t,_{A,B}[\widehat{g}_{AB},\left(\widehat{g}_{AB}\right)^{}]\psi _g^t>}{\psi _g^t^2}}`$
$``$ $`2{\displaystyle \frac{\mathrm{cosh}\left(p\right)\left(1e^{t/2}\right)t\frac{\mathrm{sinh}\left(p\right)}{p}+\stackrel{~}{K}_t}{1K_t}}`$
We conclude that (recall (4.100))
$`{\displaystyle \frac{<\psi _g^t,_{A,B}[\widehat{g}_{AB},\left(\widehat{g}_{AB}\right)^{}]\psi _g^t>}{\psi _g^t^2}}`$ $`=`$ $`t\left(\mathrm{cosh}\left(p\right)+2\mathrm{sinh}\left(p\right)/p\right)\left[1+O\left(t^2\right)\right]`$ (4.112)
$`=`$ $`i\mathrm{}{\displaystyle \underset{A,B}{}}\{g_{AB},\overline{g_{AB}}\}\left[1+O\left(t\right)\right]`$
that is, the uncertainty bound in terms of commutators in the coherent state $`\psi _g^t`$ is given precisely by the value of the associated Poisson bracket at the phase space point $`g`$ up to corrections quadratic in $`t`$.
The fact that the bound depends on the label of the coherent state is due to the fact that we use the operators $`\widehat{x}_{AB},\widehat{y}_{AB}`$ rather than the operators $`\widehat{p}_j,\widehat{q}_{AB}=\left(\widehat{h}_{AB}+\left(\widehat{h}_{AB}\right)^{}\right)/2`$, say for which we get the uncertainty bound
$$\frac{|<\psi _g^t,[\widehat{p}_j,\widehat{q}_{AB}]\psi _g^t>|}{2\psi _g^t^2}=t\frac{|<\psi _g^t,[\left(\tau _j\widehat{h}\right)_{AB}\left(\widehat{h}^1\tau _j\right)_{BA}]\psi _g^t>|}{4\psi _g^t^2}t/2$$
(4.113)
since $`\widehat{h}_{AB}`$ is a bounded operator on $`L_2(G,d\mu _H)`$ with bound $`1`$.
Summarizing, our coherent states saturate the uncertainty bound in precisely the way as they should and occupy a phase space volume (with respect to Liouville measure) of order $`t^3`$ exactly as the harmonic oscillator coherent states.
### 4.5 Extension to Groups of Higher Rank
Looking at the method of proof for all the theorems proved in the present section for $`G=SU\left(2\right)`$ we realize three basic steps :
I) The determination of the exact complexification $`g=g(p,h)`$ of the configuration space of the phase space $`T^{}G`$ induced by the complexifier $`p_j^2/2`$ in order to determine what quantity precisely should be peaked in either representation.
II) The use of the Poisson summation formula which transforms a slowly converging series into a rapidly converging one, allowing us to essentially drop all but one term in estimates.
III) The separate estimate of the series for disjoint ranges of the group angles due to the singular nature of functions that multiply the series, in our case $`0\theta \frac{31}{32}\pi `$ and $`\frac{31}{32}\pi \theta \pi `$ respectively, by rewriting the series in terms of parameters, here $`\delta `$, which cancel the singularities and allow to obtain uniform bounds.
How does this generalize to arbitrary compact gauge groups ?
I) The analysis in section 3 has revealed that for a general compact, semi-simple gauge group the complexification is simply polar decomposition. So this generalizes immediately to any compact, semi-simple gauge group.
II) In a Poisson summation formula is derived for any compact gauge group (see also and references therein). Basically, one uses that the coherent states depend only on the characters in the various representations. The characters in turn can be reduced to the a maximal torus $`T^r`$ (maximal Abelian subgroup) in a rank $`r`$ gauge group (generated by a $`r`$-dimensional Cartan subalgebra) moded by the action of the Weyl group. For instance, in the case $`G=SU\left(2\right)`$ there is the maximal torus $`e^{\theta \tau _3},\theta [0,\pi ]`$ and the action of the Weyl group is given by $`\tau _3\tau _3`$ and we have seen that indeed our characters were invariant under the inversion $`\lambda \lambda ^1`$ with $`\lambda =\mathrm{cosh}\left(\theta \right)+\mathrm{sinh}\left(\theta \right)=e^\theta ,\lambda ^1=\mathrm{cosh}\left(\theta \right)+\mathrm{sinh}\left(\theta \right)=e^\theta =\lambda 1`$. Then one can in fact carry out the Poisson resummation which is again of the form of a series times a product of $`r`$ singular factors of the type $`1/\mathrm{sinh}\left(\theta \right)`$ that we have seen in the case of $`SU\left(2\right)`$ which simply come out of Weyl’s character formula . The series part of that formula again obviously has the typical Gaussian damping factor that we have seen in the case of $`G=SU\left(2\right)`$.
III) For each of these singular prefactors we must make a separate estimate as outlined in this paper which is no problem in principle although the number of cases to be discussed grows as $`2^{N1}`$ !
Concluding, while possibly technically quite difficult, the methodof proof displayed in this paper for $`G=SU\left(2\right)`$ can be taken over, without principal changes, to arbitrary compact, semi-simple $`G`$. We will come back to this in .
## 5 Peakedness Proofs for Gauge-Invariant Coherent States
We first have to compute gauge-invariant coherent states from non-gauge-invariant ones. We will do this by the group averaging procedure ( and references therein). The idea is then to use the peakedness proofs of the previous section exploiting that the gauge group to be averaged over has unit volume. As will become obvious in this section, the peakedness proofs for the gauge-invariant case are under much less control than for the gauge-variant case. Fortunately, as already mentioned in for most of the applications of coherent states we can stick to the gauge-variant ones so that the lack of completeness in this section is not very serious. We leave the improvement of the estimates of the present section for future research.
###### Definition 5.1
Let $`\gamma `$ be a graph and $`E(\gamma )`$ the set of its oriented edges and $`V(\gamma )`$ the set of its vertices. Let
$$\psi _{\gamma ,\stackrel{}{g}}^t\left(\stackrel{}{h}\right)=\underset{eE\left(\gamma \right)}{}\psi _{g_e}^t\left(h_e\right)$$
(5.1)
be the family of gauge-variant coherent states on $`\gamma `$. Then we define a family of gauge-invariant coherent states on $`\gamma `$ by
$$\mathrm{\Psi }_{\gamma ,\stackrel{}{g}}^t\left(\stackrel{}{h}\right):=\eta _\gamma \psi _{\gamma ,\stackrel{}{g}}^t\left(\stackrel{}{h}\right):=\underset{vV\left(\gamma \right)}{}_G𝑑\mu _H\left(h_v\right)\underset{eE\left(\gamma \right)}{}\psi _{g_e}^t\left(h_{e\left(0\right)}h_e\left(h_{e\left(1\right)}\right)^1\right)$$
(5.2)
The operation $`\eta _\gamma `$ can actually be applied to any gauge-variant state, the result is obviously a gauge-invariant state. The following Lemma is elementary.
###### Lemma 5.1
Let, as in section 2, $`T_{\gamma ,\stackrel{}{j},\stackrel{}{J}}`$ be a complete orthonormal basis of spin-network states. Then
$$\mathrm{\Psi }_{\gamma ,\stackrel{}{g}}^t\left(\stackrel{}{h}\right)=\underset{\stackrel{}{j},\stackrel{}{J}}{}e^{\frac{t}{2}_{eE\left(\gamma \right)}\lambda _{j_e}}T_{\gamma ,\stackrel{}{j},\stackrel{}{J}}\left(\stackrel{}{g}\right)\overline{T_{\gamma ,\stackrel{}{j},\stackrel{}{J}}\left(\stackrel{}{h}\right)}$$
(5.3)
Proof of Lemma 5.1 :
Let $`\mathrm{\Delta }_\gamma =_{eE\left(\gamma \right)}\mathrm{\Delta }_e`$. By its very definition we have
$$\psi _{\gamma ,\stackrel{}{g}}^t=\left(e^{t\mathrm{\Delta }_\gamma /2}\delta _{\gamma ,\stackrel{}{h}^{}}^{noninv}\right)_{|\stackrel{}{h^{}}\stackrel{}{g}}$$
(5.4)
where $`\delta _{\gamma ,\stackrel{}{h}^{}}^{noninv}`$ is the distribution defined by $`\delta _{\gamma ,\stackrel{}{h}}^{noninv}\left(f_\gamma \right)=f_\gamma \left(\stackrel{}{h}\right)`$ for any smooth function $`f_\gamma `$ cylindrical with respect to $`\gamma `$. On the other hand
$$\underset{\stackrel{}{j},\stackrel{}{J}}{}e^{\frac{t}{2}_{eE\left(\gamma \right)}\lambda _{j_e}}T_{\gamma ,\stackrel{}{j},\stackrel{}{J}}\left(\stackrel{}{g}\right)\overline{T_{\gamma ,\stackrel{}{j},\stackrel{}{J}}\left(\stackrel{}{h}\right)}=\left(e^{t\mathrm{\Delta }_\gamma /2}\delta _{\gamma ,\stackrel{}{h}^{}}^{inv}\right)_{|\stackrel{}{h}^{}\stackrel{}{g}}$$
(5.5)
since
$$\underset{\stackrel{}{j},\stackrel{}{J}}{}T_{\gamma ,\stackrel{}{j},\stackrel{}{J}}\left(\stackrel{}{h}^{}\right)\overline{T_{\gamma ,\stackrel{}{j},\stackrel{}{J}}\left(\stackrel{}{h}\right)}$$
is indeed a representation of the $`\delta `$distribution on smooth gauge invariant functions cylindrical with respect to $`\gamma `$ as one can check on a spin-network basis.
Since $`\mathrm{\Delta }_\gamma `$ is gauge-invariant it commutes with the operation $`\eta _\gamma `$ and it remains to show that $`\eta _\gamma \delta _{\gamma ,\stackrel{}{h}}^{noninv}`$ is another representation of $`\delta _{\gamma ,\stackrel{}{h}}^{inv}`$. But if $`f_\gamma `$ is gauge invariant, smooth and cylindrical with respect to $`\gamma `$ we have (we use the notation $`\stackrel{}{h}^\stackrel{}{h}^{}:=\left\{h_{e\left(0\right)}^{}h_eh_{e\left(1\right)}^1\right\}_{eE\left(\gamma \right)}`$ for the gauge transformed vector of holonomies)
$`\left[\eta _\gamma \delta _{\gamma ,\stackrel{}{h}}^{noninv}\right]\left(f_\gamma \right)`$ (5.6)
$`=`$ $`{\displaystyle \underset{eE\left(\gamma \right)}{}}{\displaystyle _G}𝑑\mu _H\left(h_e^{\prime \prime }\right)\left[{\displaystyle \underset{vV\left(\gamma \right)}{}}{\displaystyle _G}𝑑\mu _H\left(h_v^{}\right){\displaystyle \underset{eE\left(\gamma \right)}{}}\delta (h_e^{\prime \prime },h_{e\left(0\right)}^{}h_eh_{e\left(1\right)}^1)\right]f_\gamma \left(\stackrel{}{h}^{\prime \prime }\right)`$
$`=`$ $`{\displaystyle \underset{eE\left(\gamma \right)}{}}{\displaystyle _G}𝑑\mu _H\left(h_e^{\prime \prime }\right)\left[{\displaystyle \underset{vV\left(\gamma \right)}{}}{\displaystyle _G}𝑑\mu _H\left(h_v^{}\right){\displaystyle \underset{eE\left(\gamma \right)}{}}\delta (h_{e\left(0\right)}^1h_e^{\prime \prime }h_{e\left(1\right)}^{},h_e)\right]f_\gamma \left(\stackrel{}{h}^{\prime \prime }\right)`$
$`=`$ $`{\displaystyle \underset{eE\left(\gamma \right)}{}}{\displaystyle _G}𝑑\mu _H\left(h_e^{\prime \prime }\right){\displaystyle \underset{eE\left(\gamma \right)}{}}\delta (h_e^{\prime \prime },h_e)\left[{\displaystyle \underset{vV\left(\gamma \right)}{}}{\displaystyle _G}𝑑\mu _H\left(h_v^{}\right)f_\gamma \left(\stackrel{}{h}^{\prime \prime \stackrel{}{h}^{}}\right)\right]`$
$`=`$ $`{\displaystyle \underset{eE\left(\gamma \right)}{}}{\displaystyle _G}d\mu _H\left(h_e^{\prime \prime }\right){\displaystyle \underset{eE\left(\gamma \right)}{}}\delta (h_e^{\prime \prime },h_e)f_\gamma \left(\stackrel{}{h}^{\prime \prime }\right)]`$
$`=`$ $`f_\gamma \left(\stackrel{}{h}\right)`$
by the gauge invariance of $`f_\gamma `$ and the normalization of the Haar measure.
Thus we have shown that the vectors in $`_\gamma `$ on the left hand side and right hand side of (5.3) have equal inner product with a dense set of vectors. Thus they must be the same in the $`L_2`$ sense.
$`\mathrm{}`$
Notice that the properties (i), (ii) and (iii) mentioned at the beginning of section 4 automatically hold also for gauge-invariant coherent states provided that we use only analytically continued entire gauge invariant functions in the connection representation and their complex conjugates as the classical counterparts of operators to be measured by them. The point of working with the integral representation (5.2) rather than with the explicit formula (5.3) is two-fold : First of all, we have established the peakedness proofs for the gauge-variant states (5.1) already and wish to combine those with the integral formula (5.2) while with (5.3) we would need to start from scratch. Secondly, the spin network states are not as explicitly known as one might think, the complication coming from the space of vertex contractions which involves the difficult calculus of the $`3nj`$ symbol for vertices of valence $`n+2`$.
In the next subsections we will need the following Lemma.
###### Lemma 5.2
The relation between gauge-invariant and non-gauge-invariant inner products is given by
$$<\mathrm{\Psi }_{\gamma ,\stackrel{}{g}}^t,\mathrm{\Psi }_{\gamma ,\stackrel{}{g}^{}}^t>=\underset{vV\left(\gamma \right)}{}_G𝑑\mu _H\left(h_v\right)<\psi _{\gamma ,\stackrel{}{g}}^t,\psi _{\gamma ,\stackrel{}{g}^\stackrel{}{h}}^t>$$
(5.7)
Proof of Lemma 5.2 :
The proof follows easily by the invariance properties of the Haar measure. Notice that $`\psi _g^t\left(h_1hh_2^1\right)=\psi _{h_1^1gh_2}^t\left(h\right)`$ for one copy of the group, therefore
$$\psi _{\gamma ,\stackrel{}{g}}^t\left(\stackrel{}{h}^\stackrel{}{h}^{}\right)=\psi _{\gamma ,\stackrel{}{g}^{\left(\stackrel{}{h}^1\right)}}^t\left(\stackrel{}{h}\right)$$
(5.8)
Then we have
$`<\mathrm{\Psi }_{\gamma ,\stackrel{}{g}}^t,\mathrm{\Psi }_{\gamma ,\stackrel{}{g}^{}}^t>`$ (5.9)
$`=`$ $`{\displaystyle \underset{eE\left(\gamma \right)}{}}{\displaystyle _G}𝑑\mu _H\left(\stackrel{~}{h}_e\right){\displaystyle \underset{vV\left(\gamma \right)}{}}{\displaystyle _G}𝑑\mu _H\left(h_v\right){\displaystyle _G}𝑑\mu _H\left(h_v^{}\right)\overline{\psi _{\gamma ,\stackrel{}{g}}^t\left(\stackrel{}{\stackrel{~}{h}}^\stackrel{}{h}\right)}\psi _{\gamma ,\stackrel{}{g}^{\left(\stackrel{}{h}^1\right)}}^t\left(\stackrel{}{\stackrel{~}{h}}\right)`$
$`=`$ $`{\displaystyle \underset{eE\left(\gamma \right)}{}}{\displaystyle _G}𝑑\mu _H\left(\stackrel{~}{h}_e\right){\displaystyle \underset{vV\left(\gamma \right)}{}}{\displaystyle _G}𝑑\mu _H\left(h_v\right){\displaystyle _G}𝑑\mu _H\left(h_v^{}\right)\overline{\psi _{\gamma ,\stackrel{}{g}}^t\left(\stackrel{}{\stackrel{~}{h}}\right)}\psi _{\gamma ,\stackrel{}{g}^{\left(\stackrel{}{h}^1\right)}}^t\left(\stackrel{}{\stackrel{~}{h}}^{\left(\stackrel{}{h}^1\right)}\right)`$
$`=`$ $`{\displaystyle \underset{eE\left(\gamma \right)}{}}{\displaystyle _G}𝑑\mu _H\left(\stackrel{~}{h}_e\right){\displaystyle \underset{vV\left(\gamma \right)}{}}{\displaystyle _G}𝑑\mu _H\left(h_v\right){\displaystyle _G}𝑑\mu _H\left(h_v^{}\right)\overline{\psi _{\gamma ,\stackrel{}{g}}^t\left(\stackrel{}{\stackrel{~}{h}}\right)}\psi _{\gamma ,\stackrel{}{g}^{\left(\stackrel{}{h}\left(\stackrel{}{h}^1\right)\right)}}^t\left(\stackrel{}{\stackrel{~}{h}}\right)`$
$`=`$ $`{\displaystyle \underset{eE\left(\gamma \right)}{}}{\displaystyle _G}𝑑\mu _H\left(\stackrel{~}{h}_e\right){\displaystyle \underset{vV\left(\gamma \right)}{}}{\displaystyle _G}𝑑\mu _H\left(h_v\right)\overline{\psi _{\gamma ,\stackrel{}{g}}^t\left(\stackrel{}{\stackrel{~}{h}}\right)}\psi _{\gamma ,\stackrel{}{g}^\stackrel{}{h}}^t\left(\stackrel{}{\stackrel{~}{h}}\right)`$
$`=`$ $`{\displaystyle \underset{vV\left(\gamma \right)}{}}{\displaystyle _G}d\mu _H\left(h_v\right)<\psi _{\gamma ,\stackrel{}{g}}^t,\psi _{\gamma ,\stackrel{}{g}^\stackrel{}{h}}^t>`$
$`\mathrm{}`$
So far we have defined everything for a general gauge group. We now specialize again to $`SU\left(2\right)`$ since we have proved peakedness theorems only for $`SU\left(2\right)`$ in section 4. Notice, however, that the proofs generalize to any $`G`$ once peakedness is established for the non-gauge invariant states.
### 5.1 Peakedness of the Overlap Function
In section (4.2) we showed that the overlap function for two coherent states with labels $`g,g^{}`$ is strongly peaked at $`g=g^{}`$ for one copy of the gauge group. This immediately implies that the gauge-non-invariant overlap function on $`\gamma `$ defined by
$$i_\gamma ^t(\stackrel{}{g},\stackrel{}{g}^{})=\underset{eE\left(\gamma \right)}{}i^t(g_e,g_e^{})$$
(5.10)
is strongly peaked at $`\stackrel{}{g}=\stackrel{}{g}^{}`$. We define the gauge invariant overlap function on $`\gamma `$ by
$$I_\gamma ^t(\stackrel{}{g},\stackrel{}{g}^{}):=\frac{|<\mathrm{\Psi }_{\gamma ,\stackrel{}{g}}^t,\mathrm{\Psi }_{\gamma ,\stackrel{}{g}^{}}^t>|^2}{\mathrm{\Psi }_{\gamma ,\stackrel{}{g}}^t^2\mathrm{\Psi }_{\gamma ,\stackrel{}{g}^{}}^t^2}$$
(5.11)
Then the following theorem holds.
###### Theorem 5.1
The peakedness of $`I_\gamma ^t(\stackrel{}{g},\stackrel{}{g}^{})`$ at $`[\stackrel{}{g}]=[\stackrel{}{g}^{}]`$ is implied by the peakedness of $`i_\gamma ^t(\stackrel{}{g},\stackrel{}{g}^{})`$ at $`\stackrel{}{g}=\stackrel{}{g}^{}`$. Here $`\stackrel{}{g}=\stackrel{}{H}\stackrel{}{u}`$ denotes the polar decomposition of $`\stackrel{}{g}`$ and $`[\stackrel{}{g}]:=\{\stackrel{}{g}^\stackrel{}{h}^{};\stackrel{}{h}^{}𝒢_\gamma \}`$ the gauge equivalence class of $`\stackrel{}{g}`$.
Proof of Theorem 5.1 :
Defining
$$j_\gamma ^t(\stackrel{}{g},\stackrel{}{g}^{}):=\frac{<\psi _{\gamma ,\stackrel{}{g}}^t,\psi _{\gamma ,\stackrel{}{g}^{}}^t>}{\psi _{\gamma ,\stackrel{}{g}}^t\psi _{\gamma ,\stackrel{}{g}^{}}^t}$$
(5.12)
so that $`\left|j_\gamma ^t\right|^2=i_\gamma ^t`$ we have, using Lemma 5.2
$`I_\gamma ^t(\stackrel{}{g},\stackrel{}{g}^{})`$
$`=`$ $`{\displaystyle \frac{_{vV\left(\gamma \right)}𝑑\mu _H\left(h_v\right)𝑑\mu _H\left(h_v^{}\right)<\psi _{\gamma ,\stackrel{}{g}}^t,\psi _{\gamma ,\stackrel{}{g}^\stackrel{}{h}}^t>\overline{<\psi _{\gamma ,\stackrel{}{g}}^t,\psi _{\gamma ,\stackrel{}{g}^\stackrel{}{h}^{}}^t>}}{_{vV\left(\gamma \right)}d\mu _H\left(h_v\right)d\mu _H\left(h_v^{}\right)<\psi _{\gamma ,\stackrel{}{g}}^t,\psi _{\gamma ,\stackrel{}{g}^\stackrel{}{h}}^t><\psi _{\gamma ,\stackrel{}{g}^{}}^t,\psi _{\gamma ,\stackrel{}{g}^\stackrel{}{h}^{}}^t>}}`$
$`=`$ $`{\displaystyle \frac{_{vV\left(\gamma \right)}𝑑\mu _H\left(h_v\right)𝑑\mu _H\left(h_v^{}\right)j_\gamma ^t(\stackrel{}{g},\stackrel{}{g}^\stackrel{}{h})\psi _{\gamma ,\stackrel{}{g}}^t\psi _{\gamma ,\stackrel{}{g}^\stackrel{}{h}}^t\overline{j_\gamma ^t(\stackrel{}{g},\stackrel{}{g}^\stackrel{}{h}^{})}\psi _{\gamma ,\stackrel{}{g}}^t\psi _{\gamma ,\stackrel{}{g}^\stackrel{}{h}^{}}^t}{_{vV\left(\gamma \right)}𝑑\mu _H\left(h_v\right)𝑑\mu _H\left(h_v^{}\right)j_\gamma ^t(\stackrel{}{g},\stackrel{}{g}^\stackrel{}{h})\psi _{\gamma ,\stackrel{}{g}}^t\psi _{\gamma ,\stackrel{}{g}^\stackrel{}{h}}^tj_\gamma ^t(\stackrel{}{g}^{},\stackrel{}{g}^\stackrel{}{h}^{})\psi _{\gamma ,\stackrel{}{g}^{}}^t\psi _{\gamma ,\stackrel{}{g}^\stackrel{}{h}^{}}^t}}`$
Notice that $`\psi _{\gamma ,\stackrel{}{g}^\stackrel{}{h}}^t=\psi _{\gamma ,\stackrel{}{g}}^t`$. The group integrals are very difficult to perform exactly in the case of a general graph (in fact, even for $`G=U\left(1\right)`$ the problem can be mapped to an Ising model !) and we will confine ourselves to an exact computation in appendix A for a simple graph and only for $`G=U\left(1\right)`$. However, peakedness can still be established heuristically as follows :
We know from section 4.2 that $`j_\gamma ^t(\stackrel{}{g},\stackrel{}{g}^{})`$ is peaked at $`\stackrel{}{g}=\stackrel{}{g}^{}`$ with decay width of order $`t`$. Thus the integral over $`\stackrel{}{h}`$ of $`j_\gamma ^t(\stackrel{}{g},\stackrel{}{g}^\stackrel{}{h})`$ will be large only if there exists $`\stackrel{}{h}`$ such that $`\stackrel{}{g},\stackrel{}{g}^\stackrel{}{h}`$ are lying in the same phase cell in which case we say that $`\left[\stackrel{}{g}\right]\left[\stackrel{}{g}^{}\right]`$. Let $`V_\gamma `$ be the volume with respect to $`_vd\mu _H\left(h_v\right)`$ of the region $`R_\gamma (\stackrel{}{g},\stackrel{}{g}^{})`$ of those $`\stackrel{}{h}`$ such that $`\stackrel{}{g},\stackrel{}{g}^\stackrel{}{h}`$ are lying in the same phase cell. By translation invariance of the Haar measure this volume is independent of $`\stackrel{}{g},\stackrel{}{g}^{}`$ once it is true that $`\left[\stackrel{}{g}\right]\left[\stackrel{}{g}^{}\right]`$. Therefore $`j_\gamma ^t(\stackrel{}{g},\stackrel{}{g}^\stackrel{}{h})1`$ if $`\stackrel{}{h}R_\gamma (\stackrel{}{g},\stackrel{}{g}^{})`$ and $`\left[\stackrel{}{g}\right]\left[\stackrel{}{g}^{}\right]`$ and $`j_\gamma ^t(\stackrel{}{g},\stackrel{}{g}^\stackrel{}{h})0`$ otherwise. In other words,
$$j_\gamma ^t(\stackrel{}{g},\stackrel{}{g}^\stackrel{}{h})\chi _{R_\gamma (\stackrel{}{g},\stackrel{}{g}^{})}\left(\stackrel{}{h}\right)j_\gamma ^t(\left[\stackrel{}{g}\right]_0,\left[\stackrel{}{g}^{}\right]_0)$$
(5.14)
where $`\chi `$ denotes the set-theoretic characteristic function. Here it is understood that we choose from each gauge equivalence class $`\left[\stackrel{}{g}\right]`$ once and for all a representant $`\left[\stackrel{}{g}\right]_0`$. Since $`i_\gamma ^t`$ almost takes only the values $`0`$ or $`1`$ as $`t0`$ we see that the choice of the representant is irrelevant.
Thus, the numerator in (5.1) is given approximately by
$$V_\gamma ^2\left|j_\gamma ^t(\left[\stackrel{}{g}\right]_0,\left[\stackrel{}{g}^{}\right]_0)\right|^2\psi _{\gamma ,\stackrel{}{g}}^t^2\psi _{\gamma ,\stackrel{}{g}^{}}^t^2=V_\gamma ^2i_\gamma ^t(\left[\stackrel{}{g}\right]_0,\left[\stackrel{}{g}^{}\right]_0)\psi _{\gamma ,\stackrel{}{g}}^t^2\psi _{\gamma ,\stackrel{}{g}^{}}^t^2$$
while the denominator is approximately given by
$$V_\gamma ^2\psi _{\gamma ,\stackrel{}{g}}^t^2\psi _{\gamma ,\stackrel{}{g}^{}}^t^2$$
Summarizing, we find
$$I_\gamma ^t(\stackrel{}{g},\stackrel{}{g}^{})i_\gamma ^t(\left[\stackrel{}{g}\right]_0,\left[\stackrel{}{g}^{}\right]_0)$$
(5.15)
meaning that there exists a gauge $`\stackrel{}{h}`$ such that $`I_\gamma ^t(\stackrel{}{g},\stackrel{}{g}^{})i_\gamma ^t(\stackrel{}{g},\stackrel{}{g}^\stackrel{}{h})`$.
$`\mathrm{}`$
Another argument proceeds as follows : it may be difficult to do in practice but it is possible in principle to separate $`\stackrel{}{g}`$ or $`\stackrel{}{p},\stackrel{}{h})`$ into 1) gauge invariant quantities that are non-vanishing on the constraint surface of the phase space on the one hand and 2) pure gauge quantities and those that vanish on constraint surface on the other hand. The gauge-variant overlap function is Gaussian peaked with respect to both sets of quantities and doing the integrals in (5.1) on the constraint surface does not change this behaviour with respect to the first set of quantities, in other words, if not the gauge invariant data of $`\stackrel{}{g},\stackrel{}{g}^{}`$ are close to each other then $`I^t(\stackrel{}{g},\stackrel{}{g}^{})`$ is still small.
Remark :
Given a generic graph $`\gamma `$ with $`\left|E\left(\gamma \right)\right|>2`$ edges and $`\left|V\left(\gamma \right)\right|>1`$ vertices the number of configuration degrees of freedom before taking the Gauss constraint into account is $`\left|E\left(\gamma \right)\right|dim\left(G\right)`$ and after $`\left(\left|E\left(\gamma \right)\right|\left|V\left(\gamma \right)\right|\right)dim\left(G\right)`$ if $`G`$ is non-Abelian and $`\left(\left|E\left(\gamma \right)\right|\left|V\left(\gamma \right)\right|+1\right)dim\left(G\right)`$ if $`G`$ is Abelian since in that case the gauge transformations at one of the vertices can be absorbed into those of another. Therefore, the volume $`V_\gamma `$ of the pure gauge degrees of freedom that contribute to $`I^t(\stackrel{}{g},\stackrel{}{g}^{})`$ in (5.1) should be of the order $`V\left(\gamma \right)=\sqrt{t}^{\left|V\left(\gamma \right)\right|dim\left(G\right)}`$ and $`V\left(\gamma \right)=\sqrt{t}^{\left(\left|V\left(\gamma \right)\right|1\right)dim\left(G\right)}`$ respectively since the decay width of our coherent states is $`\sqrt{t}`$ for all degrees of freedom (unquenched). This is confirmed in our example calculation in appendix A.
### 5.2 Peakedness in the Connection Representation
In section 4.1 we showed that the non-gauge-invariant probability density in the configuration representation is peaked at $`h_e=u_e`$ for all $`eE\left(\gamma \right)`$ in the state $`\psi _{g_e}^t`$ where $`g_e=H_eu_e`$ is the polar decomposition of $`g_e`$. Thus we have also shown that the non-gauge-invariant probability density on the whole graph $`\gamma `$
$$p_{\gamma ,\stackrel{}{g}}^t\left(\stackrel{}{h}\right):=\frac{\left|\psi _{\gamma ,\stackrel{}{g}}^t\left(\stackrel{}{h}\right)\right|^2}{\psi _{\gamma ,\stackrel{}{g}}^t^2}$$
(5.16)
is peaked at $`\stackrel{}{h}=\stackrel{}{u}`$. We define the gauge-invariant probability density by
$$P_{\gamma ,\stackrel{}{g}}^t\left(\stackrel{}{h}\right):=\frac{\left|\mathrm{\Psi }_{\gamma ,\stackrel{}{g}}^t\left(\stackrel{}{h}\right)\right|^2}{\mathrm{\Psi }_{\gamma ,\stackrel{}{g}}^t^2}$$
(5.17)
Then the following theorem is easy to prove.
###### Theorem 5.2
The peakedness of $`P_{\gamma ,\stackrel{}{g}}^t(\stackrel{}{h})`$ at $`[\stackrel{}{h}]=[\stackrel{}{u}]`$ is implied by the peakedness of $`p_{\gamma ,\stackrel{}{g}}^t(\stackrel{}{h})`$ at $`\stackrel{}{h}=\stackrel{}{u}`$. Here $`\stackrel{}{g}=\stackrel{}{H}\stackrel{}{u}`$ denotes the polar decomposition of $`\stackrel{}{g}`$ and $`[\stackrel{}{h}]:=\{\stackrel{}{h}^\stackrel{}{h}^{};\stackrel{}{h}^{}𝒢_\gamma \}`$ the gauge equivalence class of $`\stackrel{}{h}`$.
Proof of Theorem 5.2 :
Let us define the quantity
$$b_{\gamma ,\stackrel{}{g}}^t\left(\stackrel{}{h}\right):=\frac{\psi _{\gamma ,\stackrel{}{g}}^t\left(\stackrel{}{h}\right)}{\psi _{\gamma ,\stackrel{}{g}}^t}$$
(5.18)
so that $`\left|b_{\gamma ,\stackrel{}{g}}^t\left(\stackrel{}{h}\right)\right|^2=p_{\gamma ,\stackrel{}{g}}^t\left(\stackrel{}{h}\right)`$. Then we have by Lemma 5.2 and Definition 5.1
$`P_{\gamma ,\stackrel{}{g}}^t\left(\stackrel{}{h}\right)`$ (5.19)
$`=`$ $`{\displaystyle \frac{_{vV\left(\gamma \right)}_G𝑑\mu _H\left(u_v\right)𝑑\mu _H\left(u_v^{}\right)\psi _{\gamma ,\stackrel{}{g}}^t\left(\stackrel{}{h}^\stackrel{}{u}\right)\overline{\psi _{\gamma ,\stackrel{}{g}}^t\left(\stackrel{}{h}^\stackrel{}{u}^{}\right)}}{\mathrm{\Psi }_{\gamma ,\stackrel{}{g}}^t^2}}`$
$`=`$ $`{\displaystyle \frac{\psi _{\gamma ,\stackrel{}{g}}^t^2_{vV\left(\gamma \right)}_G𝑑\mu _H\left(u_v\right)𝑑\mu _H\left(u_v^{}\right)b_{\gamma ,\stackrel{}{g}}^t\left(\stackrel{}{h}^\stackrel{}{u}\right)\overline{b_{\gamma ,\stackrel{}{g}}^t\left(\stackrel{}{h}^\stackrel{}{u}^{}\right)}}{\mathrm{\Psi }_{\gamma ,\stackrel{}{g}}^t^2}}`$
$``$ $`{\displaystyle \frac{_{vV\left(\gamma \right)}_G𝑑\mu _H\left(u_v\right)𝑑\mu _H\left(u_v^{}\right)b_{\gamma ,\stackrel{}{g}}^t\left(\stackrel{}{h}^\stackrel{}{u}\right)\overline{b_{\gamma ,\stackrel{}{g}}^t\left(\stackrel{}{h}^\stackrel{}{u}^{}\right)}}{V_\gamma }}`$
where in the last line we have used a result established in the course of the proof of Theorem 5.1. Now from section 4.1 we know that $`b_{\gamma ,\stackrel{}{g}}^t\left(\stackrel{}{h}^\stackrel{}{u}\right)`$ is not small only if there exists $`\stackrel{}{u}`$ such that $`\stackrel{}{h}^\stackrel{}{u},\stackrel{}{U}`$ lie in the same configuration cell in which case we say that $`\left[\stackrel{}{U}\right]\left[\stackrel{}{h}\right]`$. Here,$`\stackrel{}{g}=\stackrel{}{H}\stackrel{}{U}`$ is the polar decomposition of $`\stackrel{}{g}`$. The volume of the region $`R_\gamma (\stackrel{}{U},\stackrel{}{h})`$ of $`\stackrel{}{u}`$’s such that this condition is satisfied is $`V_\gamma `$ again. Using the same notation as in Theorem 5.1 and choosing from each class $`\left[\stackrel{}{h}\right]`$ a representant $`\left[\stackrel{}{h}\right]_0`$ such that $`\left[\stackrel{}{h}\right]_0=\left[\stackrel{}{U}\right]_0`$ if $`\left[\stackrel{}{U}\right]=\left[\stackrel{}{h}\right]`$ where $`\left[\stackrel{}{U}\right]_0`$ is determined by $`\left[\stackrel{}{g}\right]_0`$ for $`g=HU`$ we see that
$$b_{\gamma ,\stackrel{}{g}}^t\left(\stackrel{}{h}^\stackrel{}{u}\right)b_{\gamma ,\left[\stackrel{}{g}\right]_0}^t\left(\left[\stackrel{}{h}\right]_0\right)\chi _{R_\gamma (\stackrel{}{g},\stackrel{}{h})}\left(\stackrel{}{u}\right)$$
(5.20)
Therefore (5.19) becomes
$$P_{\gamma ,\stackrel{}{g}}^t\left(\stackrel{}{h}\right)V_\gamma p_{\gamma ,\left[\stackrel{}{g}\right]_0}^t\left(\left[\stackrel{}{h}\right]_0\right)$$
(5.21)
$`\mathrm{}`$
### 5.3 Peakedness in the Electric Field Representation
The gauge-non-invariant coherent states for a graph in the electric field representation are simply given by the product of the ones for each edge
$$\stackrel{~}{\psi }_{\gamma ,\stackrel{}{g}}^t(\stackrel{}{j},\stackrel{}{m},\stackrel{}{n}):=\underset{eE\left(\gamma \right)}{}\stackrel{~}{\psi }_{g_e}^t\left(j_em_en_e\right)$$
(5.22)
Alternatively they can be defined as the inner product of the state $`\psi _{\gamma ,\stackrel{}{g}}^t`$ defined above with the state $`|\stackrel{}{j}\stackrel{}{m}\stackrel{}{n}>`$ given by
$$<\stackrel{}{h},\stackrel{}{j}\stackrel{}{m}\stackrel{}{n}>=\underset{eE\left(\gamma \right)}{}\pi _{j_e}\left(h_e\right)_{m_en_e}$$
(5.23)
Similarly, we define the gauge-invariant coherent states in the electric field representation by
$$\stackrel{~}{\mathrm{\Psi }}_{\gamma ,\stackrel{}{g}}^t(\stackrel{}{j},\stackrel{}{J}):=<\stackrel{}{j}\stackrel{}{J},\mathrm{\Psi }_{\gamma ,\stackrel{}{g}}^t>$$
(5.24)
where
$$<\stackrel{}{h},\stackrel{}{j}\stackrel{}{J}>=T_{\gamma ,\stackrel{}{j},\stackrel{}{J}}\left(\stackrel{}{h}\right)$$
(5.25)
is a spin-network state. Clearly, these gauge invariant Fourier coefficients belong to an $`\mathrm{}_2`$ Hilbert space of sequences equipped with an inner product isometric to the one on the $`L_2`$ space and it is given by $`_{\stackrel{}{j}\stackrel{}{J}}\overline{a_{\stackrel{}{j}\stackrel{}{J}}}b_{\stackrel{}{j}\stackrel{}{J}}`$. Thanks to Lemma (5.1) we can explicitly compute (5.24) to be
$$\stackrel{~}{\mathrm{\Psi }}_{\gamma ,\stackrel{}{g}}^t(\stackrel{}{j},\stackrel{}{J})=e^{\frac{t}{2}_{eE\left(\gamma \right)}j_e\left(j_e+1\right)}T_{\gamma ,\stackrel{}{j},\stackrel{}{J}}\left(\stackrel{}{g}\right)$$
(5.26)
Finally we define the gauge-invariant probability amplitude in the electric field representation by
$$P_{\gamma ,\stackrel{}{g}}^t\left(\stackrel{}{j}\stackrel{}{J}\right):=\frac{\left|\stackrel{~}{\mathrm{\Psi }}_{\gamma ,\stackrel{}{g}}^t\left(\stackrel{}{j}\stackrel{}{J}\right)\right|^2}{\mathrm{\Psi }_{\gamma ,\stackrel{}{g}}^t^2}$$
(5.27)
In order to exploit the peakedness properties established in subsection 4.3 we must know the explicit definition of $`T_{\gamma ,\stackrel{}{j},\stackrel{}{J}}\left(\stackrel{}{g}\right)`$.
###### Lemma 5.3
Denote by $`N(v)`$ the valence of a vertex $`v`$ of a graph $`\gamma `$ and split each edge $`e`$ of $`\gamma `$ into two halves with outgoing orientations from those endpoints that are vertices of $`\gamma `$. For each vertex $`v`$ of $`\gamma `$, choose a labelling of the split edges $`f_k^v,k=1,..,N(v)`$ incident at it. Given an unsplit edge $`e`$, let natural numbers $`k(e),l(e)`$ be defined by $`e=f_{k(e)}^{e(0)}(f_{l(e)}^{e(1)})^1`$ and define $`j_k^v=j_e`$ if $`k=k(e),v=e(0)`$ or $`k=l(e),v=e(1)`$. Also, for each vertex $`v`$ choose a recoupling scheme $`(J_{k1}^vj_{k+1}^v)J_k^v,k=1,..,N(v)1,J_0^v=j_1^v,J_{N(v)2}^v=j_{N(v)}^v,J_{N(v)1}^v=0`$. Finally, let $`\stackrel{}{j}^v=\{j_1^v,..,j_{N(v)}^v\},\stackrel{}{m}^v=\{m_1^v,..,m_{N(v)}^v\},\stackrel{}{J}^v=\{J_1^v,..,J_{N(v)3}^v\}`$.
A spin-network basis is then given by
$$T_{\gamma ,\stackrel{}{j},\stackrel{}{J}}\left(\stackrel{}{h}\right)=\underset{vV\left(\gamma \right)}{}\left[\underset{k=1}{\overset{N\left(v\right)}{}}\pi _{j_k^v}\left(h_{f_k^v}\right)_{m_k^vn_k^v}c_{\stackrel{}{j}^v\stackrel{}{m}^v;\stackrel{}{J}^v}^v\right]\left[\underset{eE\left(\gamma \right)}{}\pi _{j_e}\left(ϵ\right)_{n_{k\left(e\right)}^{e\left(0\right)}n_{l\left(e\right)}^{e\left(1\right)}}\right]$$
(5.28)
where $`ϵ`$ is the totally skew tensor density of weight one in two dimensions and
$$c_{\stackrel{}{j}^v\stackrel{}{m}^v;\stackrel{}{J}^v}^v=<j_1^vm_1^v..j_{N\left(v\right)}^vm_{N\left(v\right)}^v|j_1^v..j_{N\left(v\right)}^v;J_1^v..J_{N\left(v\right)3}^v>$$
(5.29)
is the Clebsch-Gordan-coefficient for recoupling of $`N(v)`$ angular momenta.
Proof of Lemma 5.3 :
We simply have to compute the inner products of two of the states in (5.28) with labels $`\stackrel{}{j},\stackrel{}{J}`$ and $`\stackrel{}{j}^{}\stackrel{}{J}^{}`$ respectively. Clearly we get the non-vanishing result if and only if $`\stackrel{}{j}^v=\stackrel{}{j}^v,\stackrel{}{m}^v=\stackrel{}{m}^v,\stackrel{}{n}^v=\stackrel{}{n}^v`$ for all $`v`$ in which case (5.28) gets mutiplied by $`_{eE\left(\gamma \right)}1/d_{j_e}`$. Thus we find for the inner product
$$=\delta _{\stackrel{}{j}\stackrel{}{j}^{}}\underset{k=1}{\overset{N\left(v\right)}{}}\left[c_{\stackrel{}{j}^v\stackrel{}{m}^v;\stackrel{}{J}^v}^vc_{\stackrel{}{j}^v\stackrel{}{m}^v;\stackrel{}{J}^v}^v\right]\left[\underset{eE\left(\gamma \right)}{}\frac{1}{d_{j_e}}\pi _{j_e}\left(ϵ\right)_{n_{k\left(e\right)}^{e\left(0\right)}n_{l\left(e\right)}^{e\left(1\right)}}\pi _{j_e}\left(ϵ\right)_{n_{k\left(e\right)}^{e\left(0\right)}n_{l\left(e\right)}^{e\left(1\right)}}\right]$$
(5.30)
Performing the sum over $`\stackrel{}{m}^v`$ produces a $`\delta _{\stackrel{}{J}^v,\stackrel{}{J}^v}`$ due to the completeness relations for the CG-coefficients. Performing the sum over $`\stackrel{}{n}`$ produces a $`_e\chi _{j_e}\left(ϵϵ^T\right)=_ed_{j_e}`$. Thus altogether
$$<\stackrel{}{j}\stackrel{}{J},\stackrel{}{j}^{}\stackrel{}{J}^{}>=\delta _{\stackrel{}{j}\stackrel{}{j}^{}}\delta _{\stackrel{}{J}\stackrel{}{J}^{}}$$
(5.31)
$`\mathrm{}`$
We also must compute the multiple CG-coefficients in terms of the elementary $`3j`$-symbols $`<j_1m_1j_2m_2|j_1j_2;jm>=\delta _{m,m_1+m_2}<j_1m_1j_2m_2|j_1j_2;Jm_1+m_2>,\text{max}(\left|m_1+m_2\right|,\left|j_1j_2\right|)Jj_1+j_2`$ for which approximation formulae for large $`j`$’s exist.
###### Lemma 5.4
$`<j_1m_1..j_Nm_N|J_1..J_{N3}J_{N2}=j_NJ_{N1}=0M=0>=\delta _{m_1+..+m_N,0}\times `$
$`\times {\displaystyle \underset{k=1}{\overset{N1}{}}}<J_{k1}n_{k1}j_{k+1}m_{k+1}|J_{k1}j_{k+1};J_kn_k>`$ (5.32)
where $`n_k=_{l=1}^{k+1}m_l,J_0=j_1,J_{N2}=j_N,J_{N_1}=0`$.
Proof of Lemma 5.4 :
This follows from iterating the definition of the CG-coefficients as unitary transformation coefficients between the two orthonormal bases $`|j_1m_1j_2m_2>:=|j_1m_1>|j_2m_2>`$ and $`|j_1j_2;jm>`$. See also the standard literature on angular momentum, e.g. .
$`\mathrm{}`$
The idea is now the following : We have shown in section 4.3 that $`p_{\gamma ,\stackrel{}{g}}^t\left(\stackrel{}{j}\stackrel{}{m}\stackrel{}{n}\right)`$ is peaked at the values $`tj_e=p_e,tm_e=^Rp_e^3,tn_e=^Lp_e^3`$ where the momenta displayed correspond to the polar decompositions $`g_e=(^RH_e)u_e=u_e(^LH_e)`$ and $`{}_{}{}^{R/L}H=\mathrm{exp}(i\tau _j(^{R/L}p_e^j/2))`$. Let us denote these values as $`\stackrel{}{j}\left(\stackrel{}{g}\right),\stackrel{}{m}\left(\stackrel{}{g}\right),\stackrel{}{n}\left(\stackrel{}{g}\right)`$. Using again Theorem 5.1 and the explicit formula (5.28) we find that
$$P_{\gamma ,\stackrel{}{g}}^t(\stackrel{}{j},\stackrel{}{J})V_\gamma \left[\underset{vV\left(\gamma \right)}{}c^v\left(\stackrel{}{j}^v\stackrel{}{J}^v\stackrel{}{m}^v\right)\right]p_{\gamma ,\stackrel{}{g}}^t(\stackrel{}{j},\stackrel{}{m},\stackrel{}{n})$$
(5.33)
where now the $`m_k^v`$ contain both $`m_e`$ and $`n_e`$ according to formula (5.28). That is, the gauge invariant states $`\stackrel{~}{\mathrm{\Psi }}_{\gamma ,\stackrel{}{g}}^t(\stackrel{}{j},\stackrel{}{J})`$ are nothing else than the non-gauge-invariant states $`\stackrel{~}{\psi }_{\gamma ,\stackrel{}{g}}^t(\stackrel{}{j},\stackrel{}{m},\stackrel{}{n})`$ contracted with Clebsch-Gordan coefficients. Since $`p_{\gamma ,\stackrel{}{g}}^t(\stackrel{}{j},\stackrel{}{m},\stackrel{}{n})`$ is small unless $`j_e,m_e,n_e`$ take the values determined by $`\stackrel{}{g}`$ above, a crude estimate of the value of (5.33) is that
$$P_{\gamma ,\stackrel{}{g}}^t(\stackrel{}{j},\stackrel{}{J})V_\gamma \left[\underset{vV\left(\gamma \right)}{}c^v\left(\stackrel{}{j}^v\left(\stackrel{}{g}\right)\stackrel{}{J}^v\stackrel{}{m}^v\left(\stackrel{}{g}\right)\right)\right]p_{\gamma ,\stackrel{}{g}}^t\left(\stackrel{}{j}\stackrel{}{m}\stackrel{}{n}\right)$$
(5.34)
Combining this with Theorem 4.6 we find
###### Theorem 5.3
For large $`p_e`$ there exists a constant $`K_t`$ exponentially vanishing as $`t0`$ and independent of $`\stackrel{}{g}`$ such that
$`P_{\gamma ,\stackrel{}{g}}^t\left(\stackrel{}{j}\stackrel{}{J}\right)`$ $`\stackrel{<}{}`$ $`V_\gamma \left[{\displaystyle \underset{vV\left(\gamma \right)}{}}c^v\left(\stackrel{}{j}^v\left(\stackrel{}{g}\right)\stackrel{}{J}^v\stackrel{}{m}^v\left(\stackrel{}{g}\right)\right)^2\right]\times `$
$`\times [{\displaystyle \underset{eE\left(\gamma \right)}{}}{\displaystyle \frac{1}{2p_e}}{\displaystyle \frac{t^{3/2}}{4\sqrt{\pi }}}{\displaystyle \frac{1}{1K_t}}\times `$
$`\times e^{j/2\frac{(m_e/j_e(^Rp_e^3)/p_e)^2}{1(^Rp_e^3/p_e)^2}}e^{j/2\frac{(n_e/j_e(^Lp_e^3)/p_e)^2}{1(^Lp_e^3/p_e)^2}}e^{\frac{\left(\frac{\left(2j_e+1\right)t}{2}p_e\right)^2}{t}}\text{ if }|^{R/L}p_e^3/p_e|<1`$
$`P_{\gamma ,\stackrel{}{g}}^t(\stackrel{}{j},\stackrel{}{J})`$ $`\stackrel{<}{}`$ $`V_\gamma \left[{\displaystyle \underset{vV\left(\gamma \right)}{}}c^v\left(\stackrel{}{j}^v\left(\stackrel{}{g}\right)\stackrel{}{J}^v\stackrel{}{m}^v\left(\stackrel{}{g}\right)\right)^2\right]\times `$
$`\times [{\displaystyle \underset{eE\left(\gamma \right)}{}}{\displaystyle \frac{1}{2p_e}}{\displaystyle \frac{t^{3/2}}{4\sqrt{\pi }}}{\displaystyle \frac{1}{1K_t}}\times `$
$`\times e^{j_ep_e|m_e/j_e(^Rp)_e^3/p_e|}e^{j_ep_e|n_e/j_e(^Lp)_e^3/p_e|}e^{\frac{\left(\frac{\left(2j_e+1\right)t}{2}p_e\right)^2}{t}}\text{ if }|^{R/L}p_e^3/p_e|\stackrel{<}{}1`$
other mixed cases being treated similarly.
Theorem 5.3 is not entirely satisfactory since one would prefer to know at which values of $`\stackrel{}{J}`$ the probability amplitude is peaked. One might hope that the Clebsh Gordan coefficient itself is peaked at certain values of $`j`$ if the values of $`j_1j_2m_1m_2`$ are given which is the case if we perform the approximation (5.34).
To investigate this question we review pieces of the beautiful paper which rigorizes the classical work of Ponzano and Regge .
Given the values $`j_1,j_2,m_1,m_2,j`$ we can construct the following quantities : Let $`j_3:=j,m_3:=m_1+m_2`$. Then we define
$`\lambda _i`$ $`:=`$ $`\sqrt{j_i^2m_i^2}`$
$`\beta ^2`$ $`:=`$ $`\left[\left(\lambda _1+\lambda _2\right)^2\lambda _3^2\right]\left[\lambda _3^2\left(\lambda _1\lambda _2\right)^2\right]`$ (5.36)
The interpretation of the $`\lambda _i`$ is clear : if we interprete the $`j_i`$ as the length of vectors $`\stackrel{}{p}_i`$ in $`\mathrm{R}`$ satisfying $`\stackrel{}{p}_1+\stackrel{}{p}_2=\stackrel{}{p}_3`$ and $`m_i`$ as their 3-components then the $`\lambda _i`$ are the lengths of the projections of these vectors into the 1-2 plane. Furthermore, it is easy to see by methods of two-dimensional Euclidean geometry that $`\beta ^2`$ is proportional to the square of a triangle with side lengths $`\lambda _1,\lambda _2,\lambda _3`$ provided $`\beta ^20`$ : This defines the (classically) allowed region. Namely, it is easy to see that $`\beta ^20`$ is equivalent to $`\lambda _1+\lambda _2\lambda _3\left|\lambda _1\lambda _2\right|`$. However, there are quantum mechanically allowed ranges of the $`j_i,m_i`$ which satisfy $`\beta ^2<0`$ which defines the (classically) forbidden region. A nice graphical illustration of these regions in parameter space can be found in . The asymptotic behaviour of the CG-coefficients as the $`j_i`$ get large can be obtained by casting the Racah formula for the CG-coefficients into an integral formula and performing a steepest descent contour deformation and a saddle point approximation. These deformations need to be discussed separately for the allowed and the forbidden region.
I) Allowed region :
We define five angles : Consider a triangle in two-dimensional Euclidean space with side lengths $`\lambda _i`$. Let $`0\gamma _1,\gamma _2\pi `$ respectively be the angle between the sides of a triangle of length $`\lambda _1,\lambda _3`$ and $`\lambda _2,\lambda _3`$ respectively. Furthermore, consider a tetrahedron spanned by the vectors $`\stackrel{}{p}_1,\stackrel{}{p}_2`$ and an additional vector $`\stackrel{}{p}`$ which has large and positive 3-component and small 1,2-components. Let $`0\chi _i\pi `$ be the angle between the outward unit vectors of those faces of the tetrahedron intersecting in the edge which corresponds to the vector $`\stackrel{}{p}_i`$. Finally, take the limit of $`R\mathrm{}`$ of $`\stackrel{}{p}R\stackrel{}{e}_3`$ where $`\stackrel{}{e}_3`$ is the standard unit vector of Euclidean space in the 3-direction. Then
$`\left|<j_1m_1j_2m_2|j_3m_3>\right|^2`$ $``$ $`{\displaystyle \frac{4j_3}{\pi \left|\beta \right|}}\mathrm{cos}^2\left(\chi \pi \left[j+{\displaystyle \frac{3}{4}}\right]\right)\text{ where}`$
$`\chi `$ $`=`$ $`m_2\gamma _2m_1\gamma _1+{\displaystyle \underset{i=1}{\overset{3}{}}}\left(j_i+{\displaystyle \frac{1}{2}}\right)\chi _i`$ (5.37)
II) Forbidden region :
The way one obtains the $`\gamma _i,\chi _i`$ in the allowed region is actually by first computing $`\mathrm{cos}\left(\gamma _i\right),\mathrm{cos}\left(\chi _i\right)`$ by analytical formulae. The corresponding expressions take in the allowed region values in $`[1,1]`$. In the forbidden region these values become positive and of modulus greater than one. Thus, the angles become imaginary or the cosines turn into hyperbolic cosines. Furthermore, in the allowed region there are two saddle points which give rise to the cosine in (5.3) upon adding their contribution while in the forbidden region there is only one saddle point so that one ends up only with one exponential function of real argument. Continuing to call the “angles” $`\gamma _i,\chi _i`$ which now take range in the positive real axis one finds that
$`\left|<j_1m_1j_2m_2|j_3m_3>\right|^2`$ $``$ $`{\displaystyle \frac{4j_3}{\pi \left|\beta \right|}}\mathrm{exp}\left(2\chi \right)\text{ where}`$
$`\chi `$ $`=`$ $`m_2\gamma _2+m_1\gamma _1+{\displaystyle \underset{i=1}{\overset{3}{}}}\left(j_i+{\displaystyle \frac{1}{2}}\right)\chi _i`$ (5.38)
Strictly speaking, the forbidden region subdivides into six subregions and the “angles” are a bit differently defined in each subregion but the essential behaviour of (5.3) stays the same.
Let us now analyze (5.3), (5.3) which we interprete as the probability amplitude $`p\left(j\right):=p_{j_1j_2m_1m_2}\left(j\right)`$ for the system of two angular momenta of modulus $`j_1,j_2`$ and 3-components $`m_1,m_2`$ to couple to resulting angular momentum $`j_3=j`$ with 3-component $`m_3=m_1+m_2`$. We are interested in the maximum of that function as $`j`$ varies in its quantum mechanically allowed range max$`(\left|m_1+m_2\right|,\left|j_1j_2\right|)jj_1+j_2`$. The explicit formula for $`p\left(j\right)`$ in terms of $`j_i,m_i`$ is very complicated and the attempt to find the exact critical point leads to unfeasable transcendent equations so that we stick here to a qualitative analysis.
In the allowed region the amplitude of the CG-coefficient is governed by the relatively simple function $`j/\left|\beta \right|`$ while it oscillates rapidly as we change $`j`$ due to the $`j\pi `$ term in the argument of the cosine. The $`\chi _i,\gamma _i`$ on the other hand are slowly varying. Thus, the critical point can be analyzed by studying the function
$$f\left(x\right):=j^2/|\beta |^2=\frac{j^2}{\left[\left(\lambda _1+\lambda _2\right)^2\lambda _3^2\right]\left[\lambda _3^2\left(\lambda _1\lambda _2\right)^2\right]}=:\frac{x+m^2}{\left[\lambda _+^2x^2\right]\left[x\lambda _{}^2\right]}$$
(5.39)
where we have defined $`\lambda _\pm =\lambda _1\pm \lambda _2,x=\lambda ^2`$. One finds that
$$f^{}\left(x\right)=\frac{1}{\beta ^4}\left[j^4\left(m^2+\lambda _+^2\right)\left(m^2+\lambda _{}^2\right)\right]$$
(5.40)
Thus the critical value is at
$$j_0:=\sqrt[4]{\left(m^2+\lambda _+^2\right)\left(m^2+\lambda _{}^2\right)}$$
(5.41)
and has the following interpretation : Suppose $`j^2`$ is the square of the vector $`\stackrel{}{p}_1+\stackrel{}{p}_2`$ then
$$\lambda _{}^2+m^2j^2=j_1^2+j_2^2+2m_1m_2+2\stackrel{}{p}_1^{}\stackrel{}{p}_2^{}\lambda _+^2+m^2$$
where $`\stackrel{}{p}_i^{}`$ are the projections into the 1-2 plane. Thus, $`j_0`$ is the geometric mean of the classical extremal values of $`j^2`$. On the other hand, the expectation value of the operator $`\widehat{j}^2`$ is approximately given by $`j_1^2+j_2^2+2m_1m_2`$ which is the algebraic mean of the two extremal values of $`j^2`$. The geometric mean is never bigger than the algebraic mean. Furthermore we see that $`\sqrt{m^2+\lambda _{}^2}j_0\sqrt{m^2+\lambda _+^2}`$ which means that (5.40) is less/bigger than zero for $`j</>j_0`$ which means that $`j=j_0`$ is the only minimum. Clearly, the formula (5.3) must break down at $`\lambda =\lambda _\pm `$ where it diverges while $`0p\left(j\right)1`$.
In the forbidden region (4.41) is exponentially damped. The function in front of the exponential factor is given by $`f\left(x\right)`$ and so the critical point $`j_0`$ is now a maximum which however has to compete with the exponential dampedness.
The qualitative behaviour of $`p\left(j\right)`$ can therfore be summarized as follows :
If there is an allowed region then $`p\left(j\right)`$ is rapidly oscillating with $`j`$ in that region where the envelope is given by a function which has a minimum at $`j=j_0`$ and is increasing towards the values of $`j`$ corresponding to $`\lambda =\lambda _\pm `$. In the forbidden region $`p\left(j\right)`$ is exponentially damped where the decay width depends on $`j_1j_2m_1m_2`$. In the transition region between allowed and forbidden region we have to join these curves smoothly. If there is no allowed region (e.g. in the case $`m_2=\pm j_2`$) there is only exponential dampedness and the peak is at the transition point $`\lambda =\lambda _\pm `$.
In conclusion, $`p\left(j\right)`$ generically does not display any peakedness properties, the best that one can say is that the expectation value of $`j`$ is $`j_0`$ given above. This agrees qualitatively by fitting the values of $`<\widehat{j}^2>,<\left(\mathrm{\Delta }\widehat{j}^2\right)^2>`$ into a Gaussian distribution. Of course, it is not surprising that the values of the recoupling momenta $`J_k^v,k=1,..,N\left(v\right)3`$ are not so sharply peaked as not even classically any value of $`j`$ in the range allowed by $`j_1,j_2,m_1,m_2`$ is distinguished.
Thus, in order to make progress in that direction one must go back to (5.33) and repeat the analysis by first summing over all $`m_e,n_e`$ and then determine the peakedness properties with respect to $`\stackrel{}{j},\stackrel{}{J}`$. This, however, is beyond the scope of the present paper.
### 5.4 Phase Space Bounds and Heisenberg Uncertainty Relation
These follow essentially from the non-gauge-invariant ones by straightforward but tedious calculations and will be left to the ambitious reader.
Acknowledgements
We are very grateful to Brian Hall for extensive discussions about the coherent states introduced by himself and in particular for pointing out the importance of the Poisson summation formula in estimates. T.T. also thanks The University of California at San Diego for hospitality and financial support. We also thank Laurent Freidel for bringing reference to our attention. O.W. thanks the Studienstiftung des Deutschen Volkes for financial support.
## Appendix A The $`U(1)`$ case
In this appendix we will apply the results of this paper to the case of $`U\left(1\right)`$ as the gauge group. As will become clear, the much simpler structure of $`U\left(1\right)`$ leads to a considerable simplification of the derivation of all the results. The main reason for this is, of course, the fact that $`U\left(1\right)`$ is Abelian and as a consequence of this that all its irreducible representations are one-dimensional. This means that one has to deal with numbers only instead of matrices.
### A.1 Peakedness Proofs for Gauge-Variant Coherent states
We recall from (4) the general form of a coherent state:
$$\psi _g^t\left(h\right)=\underset{\pi }{}d_\pi e^{\frac{t}{2}\lambda _\pi }\chi _\pi \left(gh^1\right)$$
(A.1)
For the case of $`U\left(1\right)`$, $`d_\pi =1`$, the set of all irreducible representations can be parameterized by the set of integers which we denote by $`n`$, and the eigenvalue of the Laplacian is $`n^2`$. Furthermore, we parametrize the $`U\left(1\right)`$ element $`h`$ by the angle $`\theta ,\theta [0,2\pi ]`$ and $`g`$ by $`\varphi [0,2\pi ]`$ and $`p\text{ }\mathrm{R}`$. Summarizing, we have
$`h=e^{i\theta }`$
$`g=e^{i\left(\varphi ip\right)}`$
$`\chi _\pi \left(gh^1\right)=e^{in\left(\varphi \theta \right)}e^{np}\text{ for }\pi =n`$ (A.2)
and thus
$$\psi _g^t\left(h\right)=\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}e^{\frac{t}{2}n^2}e^{in\left(\varphi \theta \right)}e^{np}.$$
(A.3)
#### A.1.1 Peakedness in the Connection Representation
As in the main text we define
$$p_g^t\left(h\right)=\frac{\left|\psi _g^t\left(h\right)\right|^2}{\psi _g^t^2},$$
(A.4)
for which we would like to prove peakedness at $`\theta =\varphi `$, or, equivalently, at $`\varphi =0`$ for $`\psi _g^t\left(1\right)`$. For the norm of $`\psi `$ we immediately get
$$\psi _g^t^2=\psi _{H^2}^{2t}\left(1\right)=\underset{n}{}e^{tn^2}e^{2np}.$$
(A.5)
Now we have to write the formula for $`\psi _g^t\left(1\right)`$ in a form suitable for applying the Poisson formula:
$`\psi _g^t\left(1\right)`$ $`=`$ $`{\displaystyle \underset{n}{}}e^{\frac{t}{2}n^2}e^{in\varphi }e^{np}`$ (A.6)
$`=`$ $`{\displaystyle \underset{n}{}}e^{\left(ns\right)^2/2}e^{i\left(ns\right)\frac{\varphi }{s}}e^{\left(ns\right)\frac{p}{s}}`$
$`=`$ $`{\displaystyle \underset{n}{}}f\left(ns\right),`$
where we introduced $`s=\sqrt{t}`$. Thus we have the following function
$$f\left(x\right)=e^{x^2/2}e^{ix\frac{\varphi }{s}}e^{x\frac{p}{s}}$$
(A.7)
which satisfies all conditions for Poisson summation formula. We obtain for the Fourier transform
$$\stackrel{~}{f}\left(k\right)=\frac{1}{\sqrt{2\pi }}e^{\frac{1}{2}\left(k^2+\frac{\varphi ^2}{s^2}+2ik\frac{p}{s}2k\frac{\varphi }{s}2i\frac{\varphi p}{s^2}\frac{p^2}{s^2}\right)}.$$
(A.8)
Applying Poisson’s formula then leads to
$$\psi _g^t\left(1\right)=\sqrt{\frac{2\pi }{t}}\underset{n}{}e^{\frac{2\pi ^2n^2+\frac{1}{2}\varphi ^2+2i\pi npi\varphi p2\pi n\varphi \frac{1}{2}p^2}{t}}.$$
(A.9)
From this we can immediately read off the Poisson transformed form for the norm of $`\psi _g^t`$ as well, for which we obtain
$$\psi _{H^2}^{2t}\left(1\right)=\sqrt{\frac{\pi }{t}}\underset{n}{}e^{\frac{\pi ^2n^2+2i\pi npp^2}{t}}.$$
(A.10)
Now we are ready to calculate the probability amplitude. Inserting (A.9) and (A.10) into (A.4) we find
$`p_g^t\left(h\right)`$ $`=`$ $`2\sqrt{{\displaystyle \frac{\pi }{t}}}{\displaystyle \frac{\left|_ne^{\frac{2\pi ^2n^2+\frac{1}{2}\varphi ^2+2i\pi npi\varphi \pi 2\pi n\varphi \frac{1}{2}p^2}{t}}\right|^2}{_ne^{\frac{\pi ^2n^2+2i\pi npp^2}{t}}}}`$ (A.11)
$`=`$ $`2\sqrt{{\displaystyle \frac{\pi }{t}}}{\displaystyle \frac{e^{\frac{\varphi ^2}{t}}\left|_ne^{\frac{2\pi ^2n^22\pi n\left(\varphi ip\right)}{t}}\right|^2}{_ne^{\frac{\pi n^2+2i\pi np}{t}}}}.`$
Our next step is to determine bounds for the denominator which we denote by $`D_p^t`$:
$`D_p^t`$ $`=`$ $`{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}e^{\frac{\pi n^2}{t}}\left(\mathrm{cos}\left(2\pi np/t\right)i\mathrm{sin}\left(2\pi np/t\right)\right)`$ (A.12)
$`=`$ $`1+2{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}e^{\frac{\pi n^2}{t}}\mathrm{cos}\left(2\pi np/t\right)`$
So a lower bound is given by
$`\left|D_p^t\right|`$ $``$ $`1+2{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}e^{\frac{\pi n^2}{t}}\underset{p}{\mathrm{min}}\left(\mathrm{cos}\left(2\pi np/t\right)\right)`$ (A.13)
$`=`$ $`12{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}e^{\frac{\pi n^2}{t}}`$
$`=:`$ $`1K_t`$
where $`K_t`$ goes to zero exponentially fast when $`t`$ goes to zero. By an equivalent estimate with signs reversed we get the following upper bound: $`\left|D_p^t\right|1+K_t`$.
For the numerator we obtain the following estimate:
$`\left|N_p^t\right|`$ $``$ $`2\sqrt{{\displaystyle \frac{\pi }{t}}}e^{\frac{\varphi ^2}{t}}{\displaystyle \underset{n}{}}|e^{\frac{2\pi ^2n^22\pi n\left(\varphi ip\right)}{t}|^2}`$ (A.14)
$`=`$ $`2\sqrt{{\displaystyle \frac{\pi }{t}}}e^{\frac{\varphi ^2}{t}}{\displaystyle \underset{n}{}}e^{\frac{4\pi ^2n^2}{t}}e^{\frac{4\pi n\varphi }{t}}`$
$`=`$ $`2\sqrt{{\displaystyle \frac{\pi }{t}}}e^{\frac{\varphi ^2}{t}}\left(1+{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}e^{\frac{4\pi ^2n^2}{t}}\mathrm{cosh}\left({\displaystyle \frac{4\pi n\varphi }{t}}\right)\right)`$
$``$ $`2\sqrt{{\displaystyle \frac{\pi }{t}}}e^{\frac{\varphi ^2}{t}}\left(1+{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}e^{\frac{4\pi ^2n^2}{t}}\mathrm{cosh}\left({\displaystyle \frac{8\pi n}{t}}\right)\right)`$
$`=`$ $`2\sqrt{{\displaystyle \frac{\pi }{t}}}e^{\frac{\varphi ^2}{t}}\left(1+\stackrel{~}{K_t}\right)`$
with $`\stackrel{~}{K_t}0`$ exponentially fast for $`t0`$.
We summarize the results of this subsection in the following theorem :
###### Theorem A.1
There exist positive constants $`K_t,\stackrel{~}{K_t}`$ (independent of $`p`$ and $`\theta `$ ), decaying exponentially fast to $`0`$ as $`t0`$ such that
$$p_g^t\left(1\right)\frac{2\sqrt{\frac{\pi }{t}}e^{\frac{\varphi ^2}{t}}\left(1+\stackrel{~}{K_t}\right)}{1K_t}$$
(A.15)
#### A.1.2 Peakedness of the Overlap Function
We recall from (4.50) the expression for the overlap function:
$$i^t(g,g^{}):=\frac{|<\psi _g^t,\psi _g^{}^t>|^2}{\psi _g^t^2\psi _g^{}^t^2}=\frac{\left|\psi _{HH^{}}^{2t}\left(h\right)\right|^2}{\psi _{H^2}^{2t}\left(1\right)\psi _{\left(H^{}\right)^2}^{2t}\left(1\right)}$$
(A.16)
In our case $`H=e^p,H^{}=e^p^{}`$ and therefore $`HH^{}=e^{p+p^{}}=\stackrel{~}{H}`$, while $`h=u^{}u^1=:\stackrel{~}{h}=e^{i\left(\varphi ^{}\varphi \right)}`$. We would like to show that this overlap function is sharply peaked at $`g=g^{}`$, that is at $`\varphi =\varphi ^{}`$ and $`p=p^{}`$. To make conclusions about the convergence behaviour for $`t0`$, we again need the Poisson transformed expressions. These do not have to be calculated anew again, but can mutatis mutandis simply be taken over from the last subsection. We obtain :
$$\psi _{\stackrel{~}{H}}^{2t}\left(\stackrel{~}{h}\right)=\sqrt{\frac{\pi }{t}}\underset{n}{}e^{\frac{2\pi ^2n^2+\frac{1}{2}\left(\varphi \varphi ^{}\right)^2+2i\pi n\left(p+p^{}\right)i\left(\varphi \varphi ^{}\right)\left(p+p^{}\right)2\pi n\left(\varphi \varphi ^{}\right)\frac{1}{2}\left(p+p^{}\right)^2}{t}}$$
(A.17)
$$\psi _{H^2}^{2t}\left(1\right)=\sqrt{\frac{\pi }{t}}\underset{n}{}e^{\frac{\pi ^2n^2+2i\pi npp^2}{t}}$$
(A.18)
$$\psi _{\left(H^{}\right)^2}^{2t}\left(1\right)=\sqrt{\frac{\pi }{t}}\underset{n}{}e^{\frac{\pi ^2n^2+2i\pi np^{}\left(p^{}\right)^2}{t}}$$
(A.19)
Inserting these results into (A.16) leads to
$$i^t(g,g^{})=\frac{e^{\frac{\frac{1}{2}\left(\varphi \varphi ^{}\right)^2}{t}}e^{\frac{\frac{1}{2}\left(pp^{}\right)^2}{t}}\left|_ne^{\frac{2\pi ^2n^2+2i\pi n\left(p+p^{}\right)i\left(\varphi \varphi ^{}\right)\left(p+p^{}\right)2\pi n\left(\varphi \varphi ^{}\right)}{t}}\right|^2}{D_p^tD_p^{}^t}$$
(A.20)
where $`D_p^t`$ was defined above. Now the argument goes completely analogously to the last subsection, that is, one calculates bounds for the (square of the modulus of the) series and for $`D_p^t`$, for the latter one can actually just take over the previous results.
The final result is :
###### Theorem A.2
There exist constants $`K_t,\stackrel{~}{K_t^{}}`$ (independent of $`g,g^{}`$), decaying exponentially fast to $`0`$ as $`t0`$ such that
$$i^t(g,g^{})\frac{e^{\frac{\frac{1}{2}\left(\varphi \varphi ^{}\right)^2}{t}}e^{\frac{\frac{1}{2}\left(pp^{}\right)^2}{t}}\left(1+\stackrel{~}{K_t^{}}\right)}{\left(1K_t\right)\left(1K_t\right)}$$
(A.21)
#### A.1.3 Peakedness in the Electric Field Representation
In section (4.3) this representation was essentially defined as the “Fourier coefficients” of $`\psi _g^t`$ with respect to the orthonormal system $`|jmn>`$, given by
$$\stackrel{~}{\psi }_g^t\left(jmn\right)=e^{tj\left(j+1\right)/2}\pi _j\left(g\right)_{mn}$$
(A.22)
In the case of $`U\left(1\right)`$, where all irreducible representations are one-dimensional, the corresponding orthonormal system is labelled only by the set of integers $`|n>`$, with $`n`$ corresponding to the $`j`$ above. So for the state $`\psi _g^t`$ in the electric field representation we have
$$\stackrel{~}{\psi }_g^t\left(n\right)=e^{tn^2/2}e^{in\varphi }e^{np}.$$
(A.23)
The aim of this section is to show that
$$p_g^t\left(n\right):=\frac{\left|\stackrel{~}{\psi }_g^t\left(n\right)\right|^2}{\psi _g^t^2}=\frac{e^{tn^2}e^{2np}}{\psi _g^t^2}$$
(A.24)
is peaked at $`tn=p`$. This would complement the peakedness property in the connection representation, leading to the conclusion that the coherent states used here have the desirable property to be ”localised” at the phase space point given by $`g`$. The proof is straightforward given the results of the previous sections. Recalling that
$$\psi _g^t^2=\psi _{H^2}^{2t}\left(1\right)=\sqrt{\frac{\pi }{t}}e^{\frac{p^2}{t}}D_p^t$$
(A.25)
and the estimate for $`D_p^t`$ we can conclude that
$`p_g^t\left(n\right)`$ $``$ $`{\displaystyle \frac{\sqrt{\frac{t}{\pi }}e^{tn^2}e^{2np}e^{p^2/t}}{\left(1K_t\right)}}`$ (A.26)
$`=`$ $`{\displaystyle \frac{\sqrt{\frac{t}{\pi }}e^{\frac{\left(tnp\right)^2}{t}}}{\left(1K_t\right)}}`$
From this we immediately see that $`p_g^t`$ is bounded as $`n\mathrm{}`$ and that it is peaked sharply at $`tn=p`$ as desired. Other than in the $`SU\left(2\right)`$ case there is no qualitatively different behaviour according to whether $`p`$ is large or not. The reason is simply that $`p`$ shows up in the exponent only in this case. Notice again that in order to get an approximately continuum momentum distribution we should introduce $`p_n=nt`$ as a summation variable and have to divide (A.26) by $`\mathrm{\Delta }p_n=t`$ which then approaches indeed a $`\delta `$-distribution as $`t0`$.
#### A.1.4 Uncertainty Relation and Phase Space Bounds
There are two things we would like to show in this section. First we want to verify for the case of $`U\left(1\right)`$ that the overlap function $`i_(^tg,g^{})`$ can be interpreted as the probability density $`p^t(g,g^{})`$ to find the system at the phase space point $`g^{}`$ in the state $`\widehat{U}_t\psi _g^t`$ (see section 4 for the definition of $`\widehat{U}_t`$) times the volume of a phase space cell. Second we will calculate the commutator between the operators $`\widehat{g}`$ and $`\widehat{g}^{}`$ to verify that it has the correct semiclasssical limt, that is, the Poisson bracket between $`g`$ and $`\overline{g}`$, thus ensuring the validity of the Heisenberg uncertainty bound.
Our first task is to determine the measure $`d\nu _t`$ on the target space of the coherent state transform. We recall its definition from (4.4):
$$\widehat{U}_t:L_2(G,d\mu _H)L_2(G^{\text{ }\mathrm{C}},d\nu _t);f\left(\widehat{U}_tf\right)\left(g\right):=<\overline{\psi _g^t},f>$$
(A.27)
where in our case $`G`$ is $`U\left(1\right)`$, $`d\mu _H=\frac{d\varphi }{2\pi }`$ and $`G^{\text{ }\mathrm{C}}=\text{ }\mathrm{C}\left\{0\right\}`$. The measure $`d\nu _t`$ is to be determined from the unitarity requirement of the transform. It is easiest to check that requirement given by
$$<\widehat{U}_tf,\widehat{U}_tf^{}>_{\nu _t}=<f,f^{}>_{\mu _H}$$
(A.28)
for any $`f,f^{}`$ $`L_2(U\left(1\right),d\theta )`$ on the basis of the electric field representation. There we have $`\left(\widehat{U}_t|n>\right)\left(g\right)=e^{n^2t/2}e^{in\varphi }e^{np}`$ so the unitarity condition reads
$$\frac{d\varphi }{2\pi }𝑑\sigma \left(H\right)e^{n^2t/2}e^{in\varphi }e^{np}e^{\left(n^{}\right)^2t/2}e^{in^{}\varphi }e^{n^{}p}=\frac{d\varphi }{2\pi }e^{i\varphi \left(nn^{}\right)}$$
(A.29)
where we made the product ansatz $`d\nu _t=d\mu _Hd\sigma _t`$. This simplifies to
$$𝑑\sigma _t\left(H\right)e^{tn^2}e^{2np}=1$$
(A.30)
which by inspection is solved by a Gaussian measure in $`p`$. The precise result is given in the following Lemma:
###### Lemma A.1
The measure $`\nu _t`$ on the target space of the coherent state transform is given by
$$d\nu _t\left(g\right)=d\mu _H\left(u\right)d\sigma _t\left(H\right)=d\mu _H\left(u\right)\sqrt{\frac{1}{\pi t}}e^{\frac{p^2}{t}}dp=:\nu _t\left(g\right)d\mathrm{\Omega }$$
(A.31)
where $`g=Hu`$ is the polar decomposition of $`g`$ and $`d\mathrm{\Omega }=d\mu _H(u)dp`$ is the Liouville measure on $`T^{}U(1)`$.
Now we are ready to address our first problem. We take over from section (4.4) the general expression for $`p^t(g,g^{})`$:
$$p^t(g,g^{})=\nu _t\left(g^{}\right)\frac{\left|\left(\widehat{U}_t\psi _g^t\right)\left(g^{}\right)\right|^2}{\psi _g^t^2}=\nu _t\left(g^{}\right)\psi _g^{}^t^2i^t(g,\left(g^{}\right)^{})$$
(A.32)
where $``$ in the $`U\left(1\right)`$ case is just complex conjugation. Then using the previous results (A.10), (A.12), (A.13) that
$$\sqrt{\frac{\pi }{t}}e^{\frac{p^2}{t}}\left(1K_t\right)\psi _g^{}^t^2\sqrt{\frac{\pi }{t}}e^{\frac{p^2}{t}}\left(1+K_t\right)$$
(A.33)
and the expression for $`\nu _t`$ we find:
$$\frac{2\pi }{\left(2\pi t\right)}\left(1K_t\right)\frac{p^t(g,g^{})}{i^t(g,\left(g^{}\right)^{})}\frac{2\pi }{\left(2\pi t\right)}\left(1+K_t\right)$$
(A.34)
for some constant $`K_t`$ exponentially vanishing for $`t0`$. We summarize :
###### Theorem A.3
The overlap function $`i^t(g,g^{})`$ approaches exponentially fast with $`t0`$ the function $`p^t(g,(g^{})^{})t`$ where $`p^t(g,g^{})`$ denotes the probability density to find the system at the phase space point $`g^{}`$ in the state $`\widehat{U}_t\psi _g^t`$ with respect to the measure $`d\mathrm{\Omega }`$ in the phase space $`T^{}U(1)`$.
Thus the phase space volume occupied by a coherent state with respect to the measure $`d\mu _Hdp`$ is given by $`t\mathrm{}`$.
We now come to the commutator calculation. The classical variables are $`g`$ and $`\overline{g}`$ where $`g=Hu=e^pe^{i\varphi }`$. For their Poisson bracket we get
$`\{g,\overline{g}\}`$ $`=`$ $`\{e^pe^{i\varphi },e^pe^{i\varphi }\}`$ (A.35)
$`=`$ $`\{e^p,e^{i\varphi }\}e^{i\varphi }e^p+\{e^{i\varphi },e^p\}e^pe^{i\varphi }`$
$`=`$ $`ie^pe^{i\varphi }e^{i\varphi }e^pie^{i\varphi }e^pe^pe^{i\varphi }`$
$`=`$ $`2ie^p`$
This Poisson bracket should be proportional to the first order term (in t) of the expression
$$\frac{<\psi _g^t,[\widehat{g},\widehat{g}^{}]\psi _g^t>}{\psi _g^t^2}$$
(A.36)
One of the characteristic properties of coherent states is that $`\widehat{g}\psi _g^t\left(h\right)=g\psi _g^t\left(h\right)`$, so one term of the commutator is easy :
$$<\psi _g^t,\widehat{g}^{}\widehat{g}\psi _g^t>=\left(\overline{g}g\right)\psi _g^t^2=e^{2p}\psi _g^t^2$$
(A.37)
The second term requires a bit more work. We recall the following expressions
$`\psi _g^t\left(h\right)={\displaystyle \underset{n}{}}e^{n^2t/2}\left(gh^1\right)^n`$
$`\widehat{g}^{}=e^{t\mathrm{\Delta }/2}\left(\widehat{h}^1\right)e^{t\mathrm{\Delta }/2}`$ (A.38)
from which we calculate
$`\left(\widehat{g}^{}\psi _g^t\right)\left(h\right)`$ $`=`$ $`e^{t\mathrm{\Delta }/2}\left(\widehat{h}^1\right){\displaystyle \underset{n}{}}e^{n^2t/2}e^{n^2t/2}\left(gh^1\right)^n`$ (A.39)
$`=`$ $`e^{t\mathrm{\Delta }/2}{\displaystyle \underset{n}{}}e^{n^2t/2}g^nh^{n1}`$
$`=`$ $`{\displaystyle \underset{n}{}}e^{tn^2}e^{\frac{t}{2}\left(n+1\right)^2}g^nh^{n1}`$
Applying $`\widehat{g}`$ in a similar way we obtain
$$\left(\widehat{g}\widehat{g}^{}\psi _g^t\right)\left(h\right)=\underset{n}{}e^{tn^2}e^{\frac{t}{2}\left(n+1\right)^2}e^{\frac{t}{2}\left(n+1\right)^2}e^{\frac{t}{2}n^2}g^nh^n$$
(A.40)
and thus after some simple algebra
$`<\psi _g^t,\widehat{g}\widehat{g}^{}\psi _g^t>`$ $`=`$ $`{\displaystyle \underset{n}{}}e^{2nt+t}e^{tn^2}e^{2np}`$ (A.41)
$`=`$ $`{\displaystyle \underset{n}{}}e^{t\left(n1\right)^2}e^{2t}e^{2np}`$
$`=`$ $`e^{2t}e^{2p}\psi _g^t^2`$
where we relabelled the summation index in the last line. Combining all results we find
$`{\displaystyle \frac{<\psi _g^t,[\widehat{g},\widehat{g}^{}]\psi _g^t>}{\psi _g^t^2}}`$ $`=`$ $`\left(e^{2t}1\right)e^{2p}`$ (A.42)
$`=`$ $`2te^{2p}+O\left(t^2\right)`$
$`=`$ $`it\{g,\overline{g}\}+O\left(t^2\right)`$
which is the desired result.
### A.2 Peakedness Proofs for Gauge-Invariant Coherent States
In section 5 of the main text first the notion of gauge-invariant coherent states was introduced and several properties were proved. This was done for a general compact Lie group so there is no need to repeat that part here for $`U\left(1\right)`$. Then the discussion was specialized to the case of $`SU\left(2\right)`$ where the peakedness in the connection representation, the electric field representation and for the overlap function was proved. We will not repeat those proofs here for the case of $`U\left(1\right)`$, rather we will illustrate them by means of a concrete example. To avoid tedious book-keeping problems we take a very simple graph $`\gamma _0`$ for the coherent state to be considered. It consists of two vertices $`v_1,v_2`$ which are connected by three edges $`e_1,e_2,e_3`$. Without loss of generality we can assume that the edges are outgoing from the same vertex. This example will be underlying all the discussion in the following subsections.
#### A.2.1 Peakedness of the Overlap function
In this subsection we want to illustrate theorem 5.1 by explicitly calculating the last line of (5.1) for our example. First we determine the form of $`j_{\gamma _0}^t(\stackrel{}{g},\stackrel{}{g}^{})`$ and then perform the integrations over the gauge group of which there are two in our case. One remark about the notation: ”$``$” will stand for equality in the limit $`t0`$, that is, terms which vanish to first order in this limit are omitted. A general $`J_\gamma ^t(\stackrel{}{g},\stackrel{}{g}^{})`$ is given in terms of
$$j_{\gamma _0}^t(\stackrel{}{g},\stackrel{}{g}^\stackrel{}{h})=\frac{<\psi _{\gamma _0,\stackrel{}{g}}^t,\psi _{\gamma _0,\stackrel{}{g}^\stackrel{}{h}}^t>}{\psi _{\gamma _0,\stackrel{}{g}}^t\psi _{\gamma _0,\stackrel{}{g}^\stackrel{}{h}}^t}$$
(A.43)
The expression for the norms for our special graph $`\gamma _0`$ can be taken over from the previous sections:
$$\psi _{\gamma _0,\stackrel{}{g}}^t=\left(\frac{\pi }{t}\right)^{3/4}e^{\frac{_{i=1}^3p_i^2}{2t}}\left(1K_t\right)$$
(A.44)
$$\psi _{\gamma _0,\stackrel{}{g}^\stackrel{}{h}}^t=\psi _{\gamma _0,\stackrel{}{g}^{}}^t=\left(\frac{\pi }{t}\right)^{3/4}e^{\frac{_{i=1}^3\left(p_i^{}\right)^2}{2t}}\left(1K_t^{}\right)$$
(A.45)
To calculate the numerator we introduce the two gauge angles $`\theta _1,\theta _2`$, associated with the gauge transformations at $`v_1,v_2`$ and their difference $`\mathrm{\Delta }\theta =\theta _2\theta _1`$. From (A.17) we find by inserting $`\mathrm{\Delta }\theta `$ in the right place
$`<\psi _{\gamma _0,\stackrel{}{g}}^t,\psi _{\gamma _0,\stackrel{}{g}^\stackrel{}{h}}^t>`$ $`=`$ $`\left({\displaystyle \frac{\pi }{t}}\right)^{3/2}e^{\frac{3\left(\mathrm{\Delta }\theta \right)^2}{4t}}e^{\frac{\mathrm{\Delta }\theta _{i=1}^3\left(\varphi _i\varphi _i^{}\right)}{2t}}e^{\frac{i\mathrm{\Delta }\theta _{i=1}3\left(p_i+p_i^{}\right)}{2t}}\times `$ (A.46)
$`\times e^{\frac{\frac{1}{2}_{i=1}^3\left(\varphi _i\varphi _i^{}\right)^2}{2t}}e^{\frac{i_{i=1}^3\left(\varphi _i\varphi _i^{}\right)\left(p_i+p_i^{}\right)}{2t}}e^{\frac{\frac{1}{2}_{i=1}^3\left(p_i+p_i^{}\right)^2}{2t}}\left(1+\stackrel{~}{K}_t\right)`$
Only the first three exponents are relevant for the integration so we will separate them out. Also, due to the special dependence of the expression on the gauge angles we perform a change of integration variables. Instead of integrating over $`\theta _1`$ an d $`\theta _2`$ we will integrate over $`\theta _{}=\left(\theta _2\theta _1\right)/2=\mathrm{\Delta }\theta /2`$ and $`\theta _+=\left(\theta _2+\theta _1\right)/2`$ where the latter integration is trivial. We obtain for the integral
$`Int`$ $`:=`$ $`{\displaystyle \frac{1}{4\pi ^2}}{\displaystyle _\pi ^\pi }𝑑\theta _{}{\displaystyle _{\left|\theta _{}\right|}^{2\pi \left|\theta _{}\right|}}𝑑\theta _+2e^{\frac{3\theta _{}^2}{t}}e^{\frac{\theta _{}\left(_{i=1}^3\left(\varphi _i\varphi _i^{}\right)i\left(p_i+p_i^{}\right)\right)}{t}}`$ (A.47)
$`=`$ $`{\displaystyle \frac{1}{4\pi ^2}}4{\displaystyle _\pi ^\pi }𝑑\theta _{}\left(2\pi \left|\theta _{}\right|\right)e^{\frac{\left(3\theta _{}^2+\theta _{}\left(_{i=1}^3\left(\varphi _i\varphi _i^{}\right)i\left(p_i+p_i^{}\right)\right)\right)}{t}}`$
$`=`$ $`{\displaystyle \frac{1}{4\pi ^2}}4\sqrt{t}{\displaystyle _{\pi /\sqrt{t}}^{\pi /\sqrt{t}}}𝑑x\left(\pi \sqrt{t}\left|x\right|\right)e^{3\left(x^2+\frac{x\left(_{i=1}^3\left(\varphi _i\varphi _i^{}\right)+i\left(r_i+r_i^{}\right)\right)}{3\sqrt{t}}\right)}`$
$``$ $`{\displaystyle \frac{1}{4\pi ^2}}4\sqrt{t}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑x\left(\pi \sqrt{t}\left|x\right|\right)e^{3\left(x^2+\frac{x\left(_{i=1}^3\left(\varphi _i\varphi _i^{}\right)i\left(p_i+p_i^{}\right)\right)}{3\sqrt{t}}\right)}`$
The second term is a derivative which when evaluated gives a finite result independent of $`t`$ times $`t`$ and therefore vanishes in the limit $`t0`$. Thus we can continue
$`Int`$ $``$ $`{\displaystyle \frac{1}{4\pi ^2}}4\pi \sqrt{t}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑xe^{3\left(x^2+\frac{x\left(_{i=1}^3\left(\varphi _i\varphi _i^{}\right)i\left(p_i+p_i^{}\right)\right)}{3\sqrt{t}}\right)}`$ (A.48)
$`=`$ $`{\displaystyle \frac{1}{4\pi ^2}}4\pi \sqrt{t}e^{\frac{\left(_{i=1}^3\left(\varphi _i\varphi _i^{}\right)i\left(p_i+p_i^{}\right)\right)^2}{12t}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑xe^{3x^2}`$
$`=`$ $`{\displaystyle \frac{1}{\sqrt{3\pi }}}\sqrt{t}e^{\frac{\left(_{i=1}^3\left(\varphi _i\varphi _i^{}\right)i\left(p_i+p_i^{}\right)\right)^2}{12t}}`$
Putting everything together we get the following expression
$$j_{\gamma _0}^t(\stackrel{}{g},\stackrel{}{g}^\stackrel{}{h})=\frac{\left(\frac{\pi }{t}\right)^{3/2}\sqrt{\frac{t}{3\pi }}e^{\frac{\left(_{i=1}^3\left(\varphi _i\varphi _i^{}\right)i\left(p_i+p_i^{}\right)\right)^2}{12t}}e^{\frac{\frac{1}{2}_{i=1}^3\left(\varphi _i\varphi _i^{}\right)^2}{2t}}e^{\frac{i_{i=1}^3\left(\varphi _i\varphi _i^{}\right)\left(p_i+p_i^{}\right)}{2t}}e^{\frac{\frac{1}{2}_{i=1}^3\left(p_i+p_i^{}\right)^2}{2t}}\left(1+\stackrel{~}{K}_t\right)}{\left(\frac{\pi }{t}\right)^{3/2}e^{\frac{_{i=1}^3\left(p_i^2+\left(p_i^{}\right)^2\right)}{2t}}\left(1K_t\right)\left(1K_t^{}\right)}$$
(A.49)
for
$$J_{\gamma _0}^t(\stackrel{}{g},\stackrel{}{g}^{}):=𝑑\mu _H\left(h_1\right)𝑑\mu _H\left(h_2\right)j_{\gamma _0}^t(\stackrel{}{g},\stackrel{}{g}^\stackrel{}{h})$$
This result is of course the same as for the integral over $`j_{\gamma _0}^t(\stackrel{}{g},\stackrel{}{g}^\stackrel{}{h}^{})`$, so the numerator of $`I_{\gamma _0}^t(\stackrel{}{g},\stackrel{}{g}^{})`$ is just given by the norm square of (A.49) (ignoring the norms which cancel with those in the denominator.) The $`j_{\gamma _0}^t`$ terms in the denominator can be obtained from (A.49) by setting $`\varphi _i=\varphi _i^{}`$ and $`p_i=p_i^{}`$. Putting everything together and after some lengthy algebra we obtain as the final result :
$$I_{\gamma _0}^t(\stackrel{}{g},\stackrel{}{g}^{})=e^{\frac{\frac{1}{3}_{i<j}\left(\varphi _i\varphi _i^{}+\varphi _j^{}\varphi _j\right)^2}{2t}}e^{\frac{_{i=1}^3\left(p_ip_i^{}\right)^2}{2t}}\left(1+\overline{K}_t\right)$$
(A.50)
where $`\overline{K}_t`$ is a constant that goes to $`0`$ for $`t0`$ exponentially fast. This expression obviously has the same peakedness property as in the gauge-variant case, but the $`\varphi _i`$ appear in a gauge-invariant way as should be expected. The $`p_i`$ are of course gauge-invariant by themselves. We still want to make the relation between the gauge-variant and the gauge-invariant overlap functions a bit clearer. To do so we choose as a set of independent variables the following combinations of the $`\varphi _i`$:
$`x_1=\varphi _2\varphi _1`$
$`x_2=\varphi _3\varphi _1`$
$`\varphi _1`$ (A.51)
and similar for the primed angles. It follows that $`\mathrm{\Delta }\varphi _2=\varphi _2\varphi _2^{}=\mathrm{\Delta }x_1+\mathrm{\Delta }\varphi _1`$ and $`\mathrm{\Delta }\varphi _3=\mathrm{\Delta }x_2+\mathrm{\Delta }\varphi _1`$. The exponent for $`i_{\gamma _0}^t`$ in the new variables reads then
$`\left(\mathrm{\Delta }\varphi _1\right)^2+\left(\mathrm{\Delta }\varphi _2\right)^2+\left(\mathrm{\Delta }\varphi _3\right)^2`$ $`=`$ $`\left(\mathrm{\Delta }\varphi _1\right)^2+\left(\mathrm{\Delta }x_1+\mathrm{\Delta }\varphi _1\right)^2+\left(\mathrm{\Delta }x_2+\mathrm{\Delta }\varphi _1\right)^2`$
$`=`$ $`3\left(\mathrm{\Delta }\varphi _1+{\displaystyle \frac{\mathrm{\Delta }x_1+\mathrm{\Delta }x_2}{3}}\right)^2+{\displaystyle \frac{2}{3}}\left(\left(\mathrm{\Delta }x_1\right)^2+\left(\mathrm{\Delta }x_2\right)^2\mathrm{\Delta }x_1\mathrm{\Delta }x_2\right)`$
where we left out some simple manipulations. To see the peakedness of the expression we have to render the last term into a quadratic form. The ansatz
$$\left(\mathrm{\Delta }x_1\right)^2+\left(\mathrm{\Delta }x_2\right)^2\mathrm{\Delta }x_1\mathrm{\Delta }x_2=\beta \left[\left(\mathrm{\Delta }x_1\alpha \mathrm{\Delta }x_2\right)^2+\left(\mathrm{\Delta }x_2\alpha \mathrm{\Delta }x_1\right)^2\right]$$
(A.53)
leads to $`\beta =\frac{1}{4\alpha }`$ and $`\alpha =2\pm \sqrt{3}`$. Thus we find
$$\left(\mathrm{\Delta }\varphi _1\right)^2+\left(\mathrm{\Delta }\varphi _2\right)^2+\left(\mathrm{\Delta }\varphi _3\right)^2=3\left(\mathrm{\Delta }\varphi _1+\frac{\mathrm{\Delta }x_1+\mathrm{\Delta }x_2}{3}\right)^2+\frac{2}{3}\frac{1}{4\alpha }\left[\left(\mathrm{\Delta }x_1\alpha \mathrm{\Delta }x_2\right)^2+\left(\mathrm{\Delta }x_2\alpha \mathrm{\Delta }x_1\right)^2\right]$$
(A.54)
which is peaked at $`\mathrm{\Delta }x_1=\mathrm{\Delta }x_2=\mathrm{\Delta }\varphi _1=0`$ as $`\alpha 1`$. Performing similar manipulations on the exponent of $`I_{\gamma _0}^t`$ leads to
$$\underset{i<j}{}\left(\varphi _i\varphi _i^{}+\varphi _j\varphi _j^{}\right)^2=2\frac{1}{4\alpha }\left[\left(\mathrm{\Delta }x_1\alpha \mathrm{\Delta }x_2\right)^2+\left(\mathrm{\Delta }x_2\alpha \mathrm{\Delta }x_1\right)^2\right]$$
(A.55)
In conclusion,
$`i_{\gamma _0}^t`$ $`=`$ $`\mathrm{exp}({\displaystyle \frac{1}{3t}}\{{\displaystyle \frac{1}{4\alpha }}[(\mathrm{\Delta }x_1\alpha \mathrm{\Delta }x_2)^2+(\mathrm{\Delta }x_2\alpha \mathrm{\Delta }x_1)^2]+{\displaystyle \frac{1}{2}}(3\mathrm{\Delta }\varphi _1+\mathrm{\Delta }x_1+\mathrm{\Delta }x_2)^2\})\times `$
$`\times e^{\frac{_{i=1}^3\left(p_ip_i^{}\right)^2}{2t}}\left(1+K_t\right)`$
$`I_{\gamma _0}^t`$ $`=`$ $`\mathrm{exp}({\displaystyle \frac{1}{3t}}{\displaystyle \frac{1}{4\alpha }}[(\mathrm{\Delta }x_1\alpha \mathrm{\Delta }x_2)^2+(\mathrm{\Delta }x_2\alpha \mathrm{\Delta }x_1)^2])\times `$ (A.56)
$`\times e^{\frac{_{i=1}^3\left(p_ip_i^{}\right)^2}{2t}}\left(1+\overline{K}_t\right)`$
So we see that in the gauge $`\mathrm{\Delta }\varphi _1+\left(\mathrm{\Delta }x_1+\mathrm{\Delta }x_2\right)/3=0`$ both expressions become equal (as usual understood in the limit $`t0`$).
Remark :
Notice that $`p_i,\varphi _i`$ and $`p_i^{},\varphi _i^{}`$ respectively are defined by $`g_i=g_{e_i}\left((A,E)\right)`$ and $`g_i^{}=g_{e_i}\left((A^{},E^{})\right)`$ respectively, that is, they come from two different points in the phase space. Nevertheless, they transform under the same gauge transformation function and so it seems surprising that we have two independent sets of gauge functions $`\theta _I,\theta _I^{},I=1,2`$. The reason is simply that there are two independent gauge group integrals appearing in (5.1).
#### A.2.2 Peakedness in the Connection Representation
In this subsection we want to illustrate the peakedness of the gauge-invariant probability density given by
$`P_{\gamma _0,\stackrel{}{g}}^t\left(\stackrel{}{h}\right)`$ $`=`$ $`{\displaystyle \frac{\left|\mathrm{\Psi }_{\gamma _0,\stackrel{}{g}}^t\left(\stackrel{}{h}\right)\right|^2}{\mathrm{\Psi }_{\gamma _0,\stackrel{}{g}}^t^2}}`$ (A.57)
$`=`$ $`{\displaystyle \frac{\psi _{\gamma _0,\stackrel{}{g}}^t^2_{vV\left(\gamma _0\right)}_G𝑑\mu _H\left(u_v\right)𝑑\mu _H\left(u_v^{}\right)b_{\gamma _0,\stackrel{}{g}}^t\left(\stackrel{}{h}^\stackrel{}{u}\right)\overline{b_{\gamma _0,\stackrel{}{g}}^t\left(\stackrel{}{h}^\stackrel{}{u}^{}\right)}}{\mathrm{\Psi }_{\gamma _0,\stackrel{}{g}}^t^2}}`$
$`\mathrm{\Psi }_{\gamma _0,\stackrel{}{g}}^t^2`$ is nothing but $`\psi _{\gamma _0,\stackrel{}{g}}^t^2_{vV\left(\gamma _0\right)}_G𝑑\mu _H\left(u_v\right)j_{\gamma _0}^t(\stackrel{}{g},\stackrel{}{g}^\stackrel{}{u})`$ which we calculated already in the last section. It remains to determine $`b_{\gamma _0,\stackrel{}{g}}^t\left(\stackrel{}{h}^\stackrel{}{u}\right)`$. As seen earlier the norm of $`\psi _{\gamma _0,\stackrel{}{g}}^t`$ is given by
$$\psi _{\gamma _0,\stackrel{}{g}}^t=\left(\frac{\pi }{t}\right)^{3/4}e^{\frac{1}{4t}_{i=1}^3p_i^2}\left(1+K_t\right)$$
(A.58)
while
$$\psi _{\gamma _0,\stackrel{}{g}}^t\left(\stackrel{}{h}^\stackrel{}{u}\right)=\underset{i=1}{\overset{3}{}}\underset{n_i}{}e^{\frac{t}{2}n_i^2}e^{in_i\varphi _i+n_ip_i}e^{in_i\alpha _i+in_i\mathrm{\Delta }\theta }$$
(A.59)
where the $`\alpha _i`$ parameterize the $`h_e`$, $`\mathrm{\Delta }\theta `$ is again the difference of the gauge angles and $`K_t`$ etc. denote constants esponentially decaying to zero as $`t0`$. After applying the Poisson summation formula and keeping only the relevant terms in an explicit form we obtain
$`\psi _{\gamma _0,\stackrel{}{g}}^t\left(\stackrel{}{h}^\stackrel{}{u}\right)`$ $`=`$ $`\left({\displaystyle \frac{2\pi }{t}}\right)^{3/2}e^{\frac{\frac{3}{2}\left(\mathrm{\Delta }\theta \right)^2}{t}}e^{+\frac{i\mathrm{\Delta }\theta _ip_i}{t}}e^{\frac{\mathrm{\Delta }\theta _i\left(\varphi _i\alpha _i\right)}{t}}\times `$ (A.60)
$`e^{\frac{\frac{1}{2}_i\left(\varphi _i\alpha _i\right)^2}{t}}e^{\frac{i_i\left(\varphi _i\alpha _i\right)p_i}{t}}e^{\frac{\frac{1}{2}_ip_i^2}{t}}\left(1+\stackrel{~}{K}_t\right)`$
We see that the $`\theta `$dependent terms are nearly the same as in the last section, the only difference being that $`\varphi _i^{}`$ is substituted by $`\alpha _i`$ and $`2t`$ in the exponent by $`t`$. Thus we can take over the final result with minor changes only. We obtain as the final result:
$$P_{\gamma _0,\stackrel{}{g}}^t\left(\stackrel{}{h}\right)=\left(2\sqrt{\frac{\pi }{t}}\right)^3\left[\sqrt{\frac{t}{3\pi }}\right]e^{\frac{_{i<j}\left(\varphi _i\alpha _i+\alpha _j\varphi _j\right)^2}{2t}}\left(1+\stackrel{~}{K}_t\right)$$
(A.61)
where the volume factor $`V_{\gamma _0}=\left[\sqrt{\frac{t}{3\pi }}\right]`$ has popped out. As expected in section 5, it is proportional to $`\sqrt{t}`$. The result is obviously gauge invariant. To compare it with the gauge-variant case one should set the $`\left\{\alpha _i\right\}`$ to $`0`$, as we calculated peakness for $`\psi _g^t\left(1\right)`$ there. A further analysis of their relation can be done analogously to the last section. We will not repeat this here.
#### A.2.3 Peakedness in the Electric Field Representation
We recall from section 5.3 the form of a gauge-invariant coherent state in the electric field representation
$$\stackrel{~}{\mathrm{\Psi }}_{\gamma _0\stackrel{}{g}}^t(\stackrel{}{j},\stackrel{}{J})=e^{\frac{t}{2}_{eE\left(\gamma _0\right)}j_e\left(j_e+1\right)}T_{\gamma _0\stackrel{}{j}\stackrel{}{J}}\left(\stackrel{}{g}\right)$$
(A.62)
which in our case becomes
$$\stackrel{~}{\mathrm{\Psi }}_{\gamma _0\stackrel{}{g}}^t\left(\stackrel{}{n}\right)=e^{\frac{t}{2}_in_i^2}T_{\gamma _0\stackrel{}{n}}\left(\stackrel{}{g}\right)=e^{\frac{t}{2}_in_i^2}e^{i_in_i\varphi _i}e^{_in_ip_i}$$
(A.63)
with the additional condition that only those $`\left\{n_i\right\}`$ are allowed for which $`_in_i=0`$. Analogously, the Gauss constraint requires that $`_ip_i=0`$. The aim is now to prove peakedness for the gauge-invariant probability amplitude given by
$$P_{\gamma _0,\stackrel{}{g}}^t\left(\stackrel{}{n}\right)=\frac{\left|\stackrel{~}{\mathrm{\Psi }}_{\gamma _0,\stackrel{}{g}}^t\left(\stackrel{}{n}\right)\right|^2}{\stackrel{~}{\mathrm{\Psi }}_{\gamma _0,\stackrel{}{g}}^t^2}$$
(A.64)
The numerator is immediately obvious from above:
$$\left|\right|\stackrel{~}{\mathrm{\Psi }}_{\gamma _0\stackrel{}{g}}^t\left(\stackrel{}{n}\right)|^2=e^{t_in_i^2}e^{2_in_ip_i}$$
(A.65)
On the other hand we verify as in the main text that
$$\left|\stackrel{~}{\mathrm{\Psi }}_{\gamma _0\stackrel{}{g}}^t\right||^2=\left|\right|\stackrel{~}{\psi }_{\gamma _0\stackrel{}{g}}^t||^2V_{\gamma _0}$$
(A.66)
which finally leads to
$$P_{\gamma _0,\stackrel{}{g}}^t\left(\stackrel{}{n}\right)=\frac{\left(\frac{t}{\pi }\right)^{3/2}e^{\frac{_i\left(tn_ip_i\right)^2}{t}}}{V_{\gamma _0}}$$
(A.67)
where still the additional condition $`_in_i=0`$ holds. This condition can be interpreted as the remains of the Clebsch-Gordan coefficients which appear in the $`SU\left(2\right)`$ case.
## Appendix B Graphical Illustration of Peakedness
This appendix contains extensive graphics, illustrating the peakedness properties of the heat kernel coherent states for $`G=U\left(1\right),SU\left(2\right)`$. Except for the first four figures, all graphics were calculated by Mathematica on the basis of a numerical approximation of the respective Poisson transformed formulas.
### B.1 The $`U(1)`$ case
#### B.1.1 Peakedness in the Connection Representation
We numerically compute the function $`\frac{\left|\psi _g^t\left(h\right)\right|}{\psi _g^t}`$, with the parameterization $`g=e^pe^{i\varphi },h=e^{i\varphi _0}`$. Without restriction, one can choose $`\varphi _0=0`$ and plot only $`p`$ and $`\varphi `$. For $`p`$ we choose the range $`[5,5]`$ while $`\varphi `$ is varied over its full range $`[\pi ,\pi ]`$. The first four plots are obtained without Poisson transformation. We contrast those with figure 5 obtained after Poisson summation, demonstrating the drastic difference in the convergence behaviour.
Figures 1 and 2 reveal the bad convergence behaviour of the non-transformed series resulting in grossly misleading plots : taking the same number $`N`$ of summation terms in the non-transformed series, the one with the higher value of $`t`$ (figure 2) is the better approximation to the actual situation (figure 5). To improve on the non-transformed series, one has to considerably increase $`N`$, as shown in figures 3 and 4.
#### B.1.2 Peakedness of the Overlap Function
We compute the function $`|<\psi _g^t,\psi _g^{}^t>|/\left(\right|\left|\psi _g^t\right|\left|\right|\left|\psi _g^{}^t\right||`$ with the parameterizations $`g=e^{p_1}e^{i\varphi },g^{}=e^{p_2}e^{i\varphi ^{}}`$. W.l.g. we set $`\varphi ^{}=0`$, choose the range of $`p_1`$ to be $`[p_25,p_2+5]`$ and take $`\varphi [\pi ,\pi ]`$. We compute plots for the values $`p_2=0,1,2,3,4`$ in figures 6 through 10.
Finally, figure 11 resolves the exact nature of the peak, exhibiting its Gauss-like structure. Its decay width is compatible with the expected value $`\sqrt{t}0.03`$.
### B.2 The $`SU(2)`$ case
#### B.2.1 Peakedness in the Connection Representation
We numerically compute the function $`\frac{\left|\psi _g^t\left(h\right)\right|}{\psi _g^t}`$, with the parametrization $`g=e^{ip^j\tau _j/2}e^{\theta ^j\tau _j}`$. Without restriction, one can chosse $`h=1`$. For the parameterization vectors we take $`p^j=(0,0,p)`$ and $`\theta ^j=\theta (\mathrm{sin}\left(\chi \right),\mathrm{cos}\left(\varphi \right),\mathrm{sin}\left(\chi \right)\mathrm{sin}\left(\varphi \right),\mathrm{cos}\left(\chi \right))`$. The variables for each graphic are $`\theta [0,\pi ]`$ and $`\varphi [\pi ,\pi ]`$. Graphics were calculated for all combinations of the values $`p=\pm 1,\pm 2`$ and $`\chi =0,\pi /4,\pi /2,3\pi /4,\pi `$. However, since they look completely identical in our resolution we display only the plots with $`p=\pm 2`$ in figures 12 through 21.
#### B.2.2 Peakedness of the Overlap Function
We numerically compute the function $`\frac{\left|\psi _g^t,\psi _g^{}^t\right|}{\psi _g^t\psi _g^{}^t}`$, with the parametrizations $`g=e^{ip^j\tau _j/2}e^{\theta ^j\tau _j}`$ $`g=e^{i\left(p^{}\right)^j\tau _j/2}e^{\left(\theta ^{}\right)^j\tau _j}`$. Without restriction, one can choose $`\theta ^{}=0`$. For the parameterization vector we take $`\theta ^j=\theta _0(\mathrm{sin}\left(\chi \right)\mathrm{cos}\left(\varphi \right),\mathrm{sin}\left(\chi \right)\mathrm{sin}\left(\varphi \right),\mathrm{cos}\left(\chi \right))`$. The variables for each graphic are $`\theta _0[0,\pi ]`$ and $`p[p^{}5,p^{}+5]`$. Graphics were calculated for all combinations of the values $`p^{}=\pm 3`$, $`\chi =0,\pi /2,\pi `$ and $`\varphi =0,\pm \pi /2,\pm \pi `$. However, since again all of them look completely identical we display only the plots with $`\varphi =0`$ in figures 22 through 27 for parallel vectors and in figure 28 for orthogonal ones.
Case A: Parallel vectors We choose $`p^j=(0,0,p)`$ and $`\left(p^{}\right)^j=(0,0,p^{})`$, that is, the momentum vectors are parallel and peakedness is therefore to be expected for $`p=p^{}`$, independent of the values for $`\chi `$ and $`\varphi `$.
Case B: Orthogonal vectors We now choose $`p^j=(0,0,p)`$ and $`\left(p^{}\right)^j=(0,p^{},0)`$, that is, the vectors are orthogonal, and therefore the overlap function should approximately vanish for all values of $`\chi `$ and $`\varphi `$. We have calculated all 30 figures as for case A but will display only one here (figure 28), as they look all the same and expectedly rather boring.
Finally, figure 29 resolves the exact nature of the peak, exhibiting its Gauss-like structure. Notice that the decay width of the peak is indeed of the order of the expected value $`\sqrt{t}0.03`$. This gives an idea of how drastically semi-classical these states will be in applications to quantum gravity in $`D=3`$ where $`t=10^{64}`$ for $`a=1`$cm !
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# 1 Introduction
## 1 Introduction
The Lagrangian appropriate for a generalized Skyrme model in the leading classical field approximation yields chiral soliton solutions. This solitons are associated with the nucleon and will be used for obtaining spatial structure information on baryons. Effectively, then, the baryon or nucleon form factors can be extracted from the soliton model equations of motion. Experimentally measured form factors can be compared with these soliton predicted spatial densities to test the model.
The present paper includes discussion of our recent work, extensions, and calculations of the nucleon electromagnetic form factors in the generalized Skyrme model. This model involves the theoretical description of the dilaton-quarkonium scalar field and shows its importance in the description of soliton dynamics. We use ”dilaton-quarkonium” scalar field to indicate the way we subdivide the gluon condensate in the calculation. This generalized model reproduces the experimental value of the nucleon mass, the input being the experimental value of the pion decay constant and the theoretically derived value of the Skyrme constant, $`e=2\pi `$.. The naive, straightforward calculation of the electromagnetic form factors has shortcomings: the values of $`F_\pi `$ and $`e`$ give too small a nucleon size and the calculated curves do not give the approximate dipole form factor values.
The generalized Skyrme model under consideration follows the formulation of Andrianov et al. This approach uses the framework of the joint chiral and conformal bosonization of the QCD Lagrangian, including chiral and scalar dilaton-quarkonium fields. In such a model the properties of the topological solitons are dramatically changed in numerical value from those in the original Skyrme model. Several authors have introduced an additional scalar field to the Skyrme model for different motivational reasons. For example, Riska and Schwesinger appear to be the first to investigate the isospin independent part of the nucleon-nucleon spin-orbit interaction when a scalar field is added. A number of papers studied the effects a scalar $`\sigma `$ meson would have by introducing it as a gluon condensate.,,, and also related, and . The purely theoretical and convincing reason is that with the introduction of a scalar field, the conformal anomaly, one of the distinctive features of the QCD Lagrangian, is reproduced. In the SU(2) sector of this Lagrangian, one can construct an effective theory which reproduces the conformal anomaly in the framework of the effective Lagrangian method, introducing a field corresponding to scale invariance. As shown in , ,and , it leads to the necessary strong attraction at intermediate internucleon distances. In such an approach the starting point is the fermion integral over quark fields, in the low energy regime of QCD. The integral is specified by the finite mode regularization scheme with a cut-off that also plays the role of a low energy boundary. Performing the joint chiral and conformal bosonization on this integral leads to an effective action for chiral $`U(x)`$ and dilaton $`\sigma (x)`$ fields. This Lagrangian favors the linear sigma model in terms of the composite field $`U(x)exp(\sigma (x))`$. The resulting effective Lagrangian, generalizing the original Skyrme Lagrangian is
$`L_{eff}(U,\sigma )`$ $`=`$ $`{\displaystyle \frac{F_\pi ^2}{4}}exp(2\sigma )Tr[_\mu U^\mu U^+]+`$ (4)
$`{\displaystyle \frac{N_fF_\pi ^2}{4}}(_\mu \sigma )^2\mathrm{exp}(2\sigma )+`$
$`+{\displaystyle \frac{1}{128\pi ^2}}Tr[_\mu UU^+,_\nu UU^+]^2`$
$`{\displaystyle \frac{C_gN_f}{48}}(e^{4\sigma }1+{\displaystyle \frac{4}{\epsilon }}(1e^{\epsilon \sigma }))`$
where the pion decay constant is taken as the experimental value, $`F_\pi =93MeV`$ and $`N_f`$ is the number of flavors. The gluon condensate, according to QCD sum rules, is $`C_g=(300400MeV)^4`$. The first two terms are the kinetic terms for the chiral and scalar fields and the third term, the well-known Skyrme term. The effective potential for the scalar field is the result of an extrapolation of the low energy potential to high energies by use of a one-loop-approximation to the Gell-Mann Low QCD $`\beta `$ \- function. The parameter $`\epsilon `$ is determined by the number of flavors $`N_f`$ as $`\epsilon =8N_f/(332N_f)`$.
## 2 The Nucleon
In the baryon sector we choose the chiral field as the spherically symmetric ansatz of Skyrme and Witten, $`U(\stackrel{}{x})=exp[i\stackrel{}{\tau }\stackrel{}{n}F(r)]`$, where $`\stackrel{}{n}=\stackrel{}{r}/|\stackrel{}{r}|`$. It is convenient to introduce a new field, $`\rho (x)=exp(\sigma (x))`$. Then, the mass functional in dimensionless variables, $`x=eF_\pi r`$, has the form $`M=M_2+M_4+V`$, where
$`M_2`$ $`=`$ $`4\pi {\displaystyle \frac{F_\pi }{e}}{\displaystyle _0^+\mathrm{}}dx[{\displaystyle \frac{N_f}{4}}x^2(\rho )^2+\rho ^2({\displaystyle \frac{x^2(F^{})^2}{2}}+sin^2F)],`$ (5)
$`M_4`$ $`=`$ $`4\pi {\displaystyle \frac{F_\pi }{e}}{\displaystyle _0^+\mathrm{}}dx({\displaystyle \frac{sin^2F}{2x^2}}+(F)^2)sin^2F,`$ (7)
$`V`$ $`=`$ $`4\pi {\displaystyle \frac{F_\pi }{e}}D_{eff}{\displaystyle _0^+\mathrm{}}𝑑xx^2[\rho ^41+{\displaystyle \frac{4}{\epsilon }}(1\rho ^\epsilon )].`$ (9)
In the last equations , the same Skyrme parameter value, $`e=2\pi `$, is used. The contribution of the potential to the mass is determined by the factor $`D_{eff}=C_gN_f/48e^2F_\pi ^4`$. The mass functional leads to a system of equations for the profile functions $`F(x)`$ and $`\rho (x)`$, where a prime is used to denote the derivative with respect to $`x`$,
$`F^{\prime \prime }[\rho ^2x^2+2sin^2F]+2F^{}x[x\rho \rho ^{}+\rho ^2]+(F^{})^2sin(2F)`$
(10)
$`\rho ^2sin(2F)sin(2F)sin^2F/x^2=0,`$ (11)
$`{\displaystyle \frac{N_f}{2}}x[x\rho ^{\prime \prime }+2\rho ^{}]2\rho [{\displaystyle \frac{x^2(F^{})^2}{2}}+sin^2F]`$
(12)
$`4D_{eff}[\rho ^3\rho ^{\epsilon 1}]x^2=0,`$ (13)
At small distances, $`F=\pi N\alpha x`$ and $`\rho =\rho (0)+\beta x^2`$, with $`\rho (0)0`$. For large $`x`$, these functions behave as $`F(x)a/x^2`$, and $`\rho (x)1b/x^6+\mathrm{}`$.
According to the virial theorem, the contributions of the individual terms of the mass functional to the energy of the system must obey the condition,
$`M_4M_23V=0,`$ (14)
which can be used to control the accuracy of the numerical solution of the system. There are nontrivial equations between the numbers $`\alpha `$ and $`\beta `$, $`a`$ and $`b`$,
$`b`$ $`=`$ $`{\displaystyle \frac{1}{2}}a^2/D_{eff},`$ (15)
$`\beta `$ $`=`$ $`\left[\rho (0)\alpha ^2+{\displaystyle \frac{4}{3}}\left(\rho ^3(0)1\right)D_{eff}\right]/N_f.`$ (17)
The choice of boundary conditions ensures a finiteness of the mass functional for a given value of the topological charge $`B=N`$. Performing canonical quantization of the rotational degrees of freedom with the collective variable method, one obtains for the nucleon mass,
$`M_B=M+S(S+1)/(2I),`$ (18)
where the moment of inertia is
$`I={\displaystyle \frac{8\pi }{3}}(F_\pi e)^3{\displaystyle _0^{\mathrm{}}}𝑑xsin^2[\rho ^2x^2+(F^{})^2x^2+sin^2F].`$ (19)
Some numerical results are presented in Table 1,
where the soliton mean square radius of the corresponding baryon $`<r_{IS}^2>`$ density distribution is given
$`<r_B^2>^{1/2}={\displaystyle \frac{1}{F_\pi e}}\left\{{\displaystyle \frac{2}{\pi }}{\displaystyle _0^{\mathrm{}}}𝑑xx^2F^{}sin^2F\right\}^{1/2}.`$ (20)
A discussion of partial restoration of chiral symmetry in this model is given in Ref.(. The restoration appears as a large deviation of $`\rho (0)`$ from its asymptotic value of $`\rho (0)=1`$. The dependence of the mass spectra on the gluon condensate in the generalized Skyrme model was also discussed in .
## 3 Form factors of charge distributions
The nucleon electric and magnetic form factors, $`G_E(q^2)`$and $`G_M(q^2)`$ can be calculated from the electromagnetic currents, in the Breit frame where the photon does not transfer energy.
$`<N_f({\displaystyle \frac{\stackrel{}{q}}{2}})|J_0(0)|N_i({\displaystyle \frac{\stackrel{}{q}}{2}})>`$ $`=`$ $`G_E(\stackrel{}{q}^2)\xi _f^+\xi _i,`$
$`<N_f({\displaystyle \frac{\stackrel{}{q}}{2}})|\stackrel{}{J}(0)|N_i({\displaystyle \frac{\stackrel{}{q}}{2}})>`$ $`=`$ $`{\displaystyle \frac{G_M(\stackrel{}{q}^2)}{2M_N}}\xi _f^+i\stackrel{}{\sigma }\stackrel{}{q}\xi _i.`$ (21)
Here, $`|N(\stackrel{}{p})>`$ is the nucleon state with momentum $`\stackrel{}{p}`$, $`\xi _i`$, $`\xi _f`$ and two component Pauli spinors, and $`\stackrel{}{q}`$ momentum transfer.
The isoscalar (S) and isovector (V) nucleon form factors are related to those for the proton and neutron by
$`G_{E,M}^{\mathrm{p},\mathrm{n}}=G_{E,M}^S\pm G_{E,M}^V`$ (22)
These form factors are normalized to the respective charge and magnetic moments by
$`G_E^\mathrm{p}(0)=1G_E^\mathrm{n}(0)=0`$ (23)
$`G_M^\mathrm{p}(0)\mu _p=2.79G_M^\mathrm{n}=\mu _n=1.91.`$ (24)
We remaked above on the smallness of the nucleon size as determined by the baryon charge density distribution in the model with a dilaton-quarkonium field.
Vector meson dominance means that the isoscalar photon sees $`\omega `$ meson $`\mathrm{𝑠𝑡𝑟𝑢𝑐𝑡𝑢𝑟𝑒}`$, but not the isoscalar baryon density $`B_0(r)`$.
According to vector meson dominance, the isoscalar current is proportional to the $`\omega _\mu `$-field,
$`J_{I=0}^\mu ={\displaystyle \frac{m_\omega ^2}{3g}}\omega _\mu (r)`$ (25)
and the corresponding charge form factor,
$`G_E^S(\stackrel{}{q}^2)={\displaystyle \frac{m_\omega ^2}{3g}}{\displaystyle d^3r\mathrm{exp}i\stackrel{}{q}\stackrel{}{r}\omega (r)}.`$ (26)
The static $`\omega (r)`$ obeys the equation,
$`(^2m_\omega ^2)\omega (r)={\displaystyle \frac{3g}{2}}B(r)={\displaystyle \frac{3gF^{}(r)}{4\pi r^2}}sin^2F(r).`$ (27)
From this equation, we obtain,
$`G_E^S(\stackrel{}{q}^2)={\displaystyle \frac{1}{2}}{\displaystyle \frac{m_\omega ^2}{m_\omega ^2+\stackrel{}{q}^2}}4\pi {\displaystyle 𝑑rr^2B_0(r)j_0(qr)}.`$ (28)
Therefore, the effecive isoscalar nucleon density is equal baryon charge density $`B_0(r)`$ times the $`\omega `$-meson propagator.
The isovector electromagnetic formfactor has analogous structure,
$`G_E^V(\stackrel{}{q}^2)={\displaystyle \frac{1}{2}}{\displaystyle \frac{m_\rho ^2}{m_\rho ^2+\stackrel{}{q}^2}}F_E^V(q^2),`$ (29)
In writing the propagators separately, as a factor multiplied into $`F_E^V`$, $`\omega `$ and $`\rho `$, themselves have no substructure or internal dynamics; the corresponding Skyrmion densities are considered as the sources of these $`\omega `$ and $`\rho `$ fields. Explicit considerations of the role of vector mesons in the electromagnetic form factors in the $`\sigma `$ model has been given by Holzwarth and quantum corrections to the relevant baryon properties in the chiral soliton models has been calculated.
The results of the present calculations are given in Figures 1 to 4.
The isoscalar part of the Skyrmion electric charge coincides with the baryon density distribution, and for the isovector density from the Skyrmion model one obtains,
$`\rho ^V(x)=sin^2F(x)\left[x^2\rho ^2(x)+\left(F^{}(x)\right)^2x^2+sin^2F(x)\right].`$ (30)
To take chiral symmetry breaking into account, we must add the pion mass term,
$`_\pi ={\displaystyle \frac{1}{4}}m_\pi ^2F_\pi ^2e^{3\sigma }Tr\left[U+U^+{\displaystyle \frac{3}{2}}e^\sigma \right],`$ (31)
to our Skyrme model Lagrangian. The theoretical predictions for the proton electric, neutron electric, proton magnetic and neutron magnetic form factors, compared with data are shown in Figs. 1, 2, 3, and 4, respectively. Corresponding values of the proton and neutron mean square radius of the electric charge distribution are 0.78 $`Fm^2`$ and - 0.19 $`Fm^2`$.
## 4 Conclusions
We have presented our calculations on the nucleon electromagnetic form factors in the framework of the generalized Skyrme model with dilaton quarkonium field. The first calculation in such a model yielded large deviations of the calculated form factors from the dipole approximation formula. In the present work, we use the empirical value of the pion decay constant and the theoretical value for the Skyrme term constant in the vector meson dominance approach to obtain a good description of the form factor data in the finite range of momentum transfer in the measurements. The vector mesons are included only as elements of the hadron substructure of the photon and are not considered as components of the structure in the soliton self-dynamics. Implicit in the approach, though not explicitly proposed, is the possibility of having the role of vector mesons given by higher derivative terms in the effective Lagrangian for soliton dynamics . For example, keeping terms to four orders in the expansion of the effective Lagrangian would lead to a $`\rho `$ meson-like term and the sixth order terms would give $`\omega `$-like terms which are important in the calculations of the form factors at the larger momentum transfers.
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# Rapidly Rotating Fermi Gases
## Abstract
We show that the density profile of a Fermi gas in rapidly rotating potential will develop prominent features reflecting the underlying Landau level like energy spectrum. Depending on the aspect ratio of the trap, these features can be a sequence of ellipsoidal volumes or a sequence of quantized steps.
Currently, there is an intense effort to cool trapped Fermi gases down to the degenerate limit. Recent experiments at JILA on <sup>40</sup>K has reached one half of its Fermi temperature. One of the motivations of cooling Fermions is to reveal their possible superfluid ground states, which can be quite novel in the case of multi-component Fermion systems. Current theoretical estimates, however, indicate that the interactions between different spin states of <sup>40</sup>K are all positive, implying a normal instead of superfluid ground state. The absence of superfluid ground states, however, does not mean that the system can not have novel macroscopic quantum phenomena. Quantum Hall effect is an excellent example. In a strong magnetic field, the energy levels of a two dimensional electron system organize into highly degenerate Landau levels, leading to a whole host of dramatic effects. In this paper, we show that similar organizations will take place in a fast rotating three dimensional trapped Fermi gas, leading to many macroscopic quantum phenomena.
For neutral atoms in rotating harmonic traps, we shall see that Landau level like energy spectrum will appear when the rotation frequency $`\mathrm{\Omega }`$ approaches the transverse confining frequency $`\omega _{}`$. This “fast rotating” limit might appear hard to achieve as the system is at the verge of flying apart due to centrifugal instability. Such instability, however, can be prevented by imposing an additional repulsive potential which dominates over the centrifugal force beyond certain radius. With centrifugal instability eliminated, the Landau-like levels will show up in many ways. We shall see that for cylindrical harmonic traps with $`\omega _z<<\omega _{}`$, where $`\omega _{}`$ and $`\omega _z`$ are the transverse and longitudinal trapping frequencies, the density is a sum of one-dimensional like density distributions residing in different “Landau volumes”. For traps with $`\omega _z`$ comparable or smaller than $`\omega _{}`$, the density consists of a set discs along $`z`$, each of which is made up of a sequence of density steps quantized in units of $`M\omega _{}/(\pi \mathrm{})`$.
2D case: We first consider 2D rotating Fermi gases in harmonic potentials since they illustrate the basic physics. The Hamiltonian in the rotating frame is
$$H\mathrm{\Omega }L_z=\frac{1}{2M}𝐩_{}^2+\frac{1}{2}M\omega _{}^2r^2\mathrm{\Omega }\widehat{𝐳}𝐫\times 𝐩_{},$$
(1)
where $`𝐩_{}=(p_x,p_y)`$ and $`𝐫=(x,y)`$. The eigenfunctions and eigenvalues of eq.(1) are
$$u_{n,m}(r,\theta )=\frac{e^{|w|^2/2}_+^m_{}^ne^{|w|^2}}{\sqrt{\pi a_{}^2n!m!}}$$
(2)
$$ϵ_{n,m}=\mathrm{}(\omega _{}+\mathrm{\Omega })n+\mathrm{}(\omega _{}\mathrm{\Omega })m+\omega _{},$$
(3)
where $`n,m`$ are non-negative integers $`0,1,2,..`$; $`w(x+iy)/a_{}`$; $`a_{}=\sqrt{\mathrm{}/M\omega _{}}`$ ; and $`_\pm (a_{}/2)(_x\pm i_y)`$. To derive eqs.(2) and (3), we note that eq.(1) can be written as $`𝚷^2/2M+(\omega _{}\mathrm{\Omega })L_z`$ with $`𝚷=𝐩_{}M\omega _{}\widehat{𝐳}\times 𝐫`$, which is precisely the canonical momentum $`𝚷=𝐩_{}\frac{eB}{2c}\widehat{𝐳}\times 𝐫`$ of an electron in a magnetic field $`𝐁`$ in the symmetric gauge, with $`eB/Mc=2\omega _{}`$. The eigenfunctions of $`𝚷^2/2M`$ are those in eq.(2), with eigenvalues $`ϵ_{n,m}^L=\mathrm{}\omega _{}(2n+1)`$, where $`n`$ is the Landau level index, and $`m`$ is an “angular momentum” index labelling the degeneracy in each level. Since $`L_zu_{n,m}=\mathrm{}(mn)u_{n,m}`$, eq.(2) is also an eigenstate of eq.(1) with eigenvalues eq.(3). Note that the function $`u_{n,m}`$ in eq.(2) peaks at
$$r_{n,m}r^2_{n,m}=a_{}^2(n+m+1),$$
(4)
and decays away as a Gaussian over a distance $`a_{}`$.
Eq.(3) shows that the system is unbounded when $`\mathrm{\Omega }>\omega _{}`$ unless an additional repulsive potential $`V_{\mathrm{wall}}(r)`$ (say, introduced by an additional optical trap) is present. We shall in particular consider potentials $`V_{\mathrm{wall}}(r)`$ which are zero for $`r<R`$ but become strongly repulsive for $`r>R`$, with $`R>>a_{}`$. The specific form of $`V_{\mathrm{wall}}`$ is not important for the key features discussed below, as long as it is smooth over length scale $`a_{}`$. The condition $`R>>a_{}`$, however, allows us to fit many $`m`$ states inside $`r<R`$ and is a necessary feature for many effects discussed below. Since $`V_{\mathrm{wall}}(r)`$ is cylindrically symmetric, the eigenstates are still labeled by quantum numbers $`(n,m)`$. For states originally with $`r_{n,m}<R`$, eqs.(2) and (3) remain valid because $`V_{\mathrm{wall}}=0`$ for $`r<R`$. For states $`(n,m)`$ originally peaked beyond $`R`$, their energies increase rapidly because $`V_{\mathrm{wall}}`$ is strongly repulsive. (For <sup>40</sup>K in a tight trap $`\omega _{}=4000`$Hz and $`10^5`$Hz, we have $`a_{}2.5\times 10^5`$cm and $`5\times 10^6`$cm resp. The condition $`R>>a_{}`$ is satisfied for $`R>5\times 10^4`$cm.)
Let us first consider the case $`\mathrm{\Omega }<\omega _{}`$ with $`V_{\mathrm{wall}}=0`$. The density in the ground state is $`\rho (r)=_{n,m}|u_{n,m}(𝐫)|^2\mathrm{\Theta }(\mu ϵ_{n,m})`$, where $`\mathrm{\Theta }(x)=1`$ or $`0`$ if $`x>0`$ or $`<0`$, $`\mu `$ is the chemical potential related to the particle number $`N`$ as $`N=_{n,m}\mathrm{\Theta }(\mu ϵ_{n,m})`$. We can write $`\rho (r)=_{n=0}^n^{}\rho _n(r;m_n^{})`$, where $`\rho _n(r;L)`$ is density contribution of the $`n`$-th Landau level with angular momentum states filled up to $`m=L`$;
$$\rho _n(r;L)=\underset{m=0}{\overset{L}{}}|u_{n,m}(𝐫)|^2\mathrm{\Theta }(\mu ϵ_{n,m}),$$
(5)
$`m_n^{}`$ is the highest angular momentum state in the $`n`$-th Landau level with energy less than $`\mu `$, and $`n^{}`$ is the highest Landau level below $`\mu `$,
$$m_n^{}=\mathrm{Int}\left[\frac{\mu /\mathrm{}\omega _{}(\omega _{}+\mathrm{\Omega })n}{(\omega _{}\mathrm{\Omega })}\right],$$
(6)
$`n^{}`$$`=\mathrm{Int}`$ $`\left[\frac{\mu /\mathrm{}\mathrm{\Omega }}{\omega _{}+\mathrm{\Omega }}\right]`$, where Int$`[x]`$ denotes the integer part of $`x`$, and $`x`$ is understood to be positive. Since $`(n,m_n^{})`$ is the state in $`\rho _n`$ farthest from the origin, its peak location ($`r_n=r_{n,m_n^{}}=a_{}\sqrt{n+m_n^{}+1}`$) gives the size of $`\rho _n`$. When $`\mathrm{\Omega }`$ is very close to $`\omega _{}`$, we have $`m_n^{}>>1`$ and
$$r_n^2=\frac{\mu \mathrm{}\mathrm{\Omega }(2n+1)}{M\omega _{}(\omega _{}\mathrm{\Omega })}.$$
(7)
Note that the difference in area between successive Landau discs is a constant
$$\pi (r_{n1}^2r_n^2)=(\pi a_{}^2)\left(\frac{2\mathrm{\Omega }}{\omega _{}\mathrm{\Omega }}\right).$$
(8)
Using eq.(2), it is straightforward to show that
$$\rho _0(r;m_0^{})=\frac{1}{2\pi a_{}^2}[1\mathrm{erf}\left(\frac{s}{\sqrt{2m_0^{}}}\right)(1+[..])].$$
(9)
where $`s=(r/a_{})^2m_0^{}`$, $`\mathrm{erf}(x)=(2/\sqrt{\pi })_0^xe^{z^2}dz`$, and the term $`[..]`$ in eq.(9) is of order $`(12m_0^{})^1`$ and smaller. The densities $`\rho _n(r)`$ of the higher Landau levels $`(n>0)`$ can be generated from $`\rho _0`$ as $`\rho _n(r;m_n^{})`$$`=\frac{1}{n!}`$ $`(\frac{1}{4}_\lambda ^2)^n`$ $`[\rho _0(\stackrel{}{r}a_{}\stackrel{}{\lambda };m_n^{})`$ $`e^{\stackrel{}{\lambda }^2}]_{\stackrel{}{\lambda }=0}`$. ¿From the properties of $`\mathrm{erf}(x)`$, it is clear that $`\rho _0`$ is a constant $`1/(\pi a_{}^2)`$ within a disc of radius $`r_0=\sqrt{m_0^{}+1}a_{}`$ and has an edge of width $`\mathrm{\Delta }_0\frac{3}{\sqrt{2}}a_{}`$. If $`m_0^{}>>1`$, $`\mathrm{\Delta }_0<<r_0`$, and $`\rho _0`$ can be approximated as a step function on scales larger than $`\mathrm{\Delta }_0`$. Likewise, $`\rho _n`$ can be approximated as a step function with somewhat larger width within the same approximation. Thus, when $`m_n^{}>>1`$, we have
$$\rho _n(r)=(\pi a_{}^2)^1\mathrm{\Theta }(r_n^2r^2).$$
(10)
Eq.(10) and (8) then imply $`\rho (r)=\frac{M\omega _{}}{\pi \mathrm{}}I(r)`$, with $`I(r)=_n\mathrm{\Theta }\left(\mu M\omega _{}(\omega _{}\mathrm{\Omega })r^2\mathrm{}\mathrm{\Omega }(2n+1)\right)`$. Using the identity $`\mathrm{Int}[x+1]=_n\mathrm{\Theta }(\alpha (xn))`$, for all $`x>0`$ and $`\alpha >0`$, $`\rho (r)`$ can be simplified to
$$\rho (r)=\frac{M\omega _{}}{\pi \mathrm{}}\mathrm{Int}\left[\frac{\mu M\omega _{}(\omega _{}\mathrm{\Omega })r^2+\mathrm{}\mathrm{\Omega }}{2\mathrm{}\mathrm{\Omega }}\right].$$
(11)
It is, however, instructive to re-derive eq.(11) in a way generalizable to arbitrary potentials. We rewrite eq.(1) as
$$H\mathrm{\Omega }L_z=\frac{(𝐩_{}M\mathrm{\Omega }\widehat{𝐳}\times 𝐫)^2}{2M}+\frac{1}{2}M(\omega _{}^2\mathrm{\Omega }^2)r^2.$$
(12)
The first term gives a set of Landau levels with spacing $`2\mathrm{}\mathrm{\Omega }`$, each of which contributes $`(\pi A^2)^1`$ to the density, where $`A^2=\mathrm{}/(M\mathrm{\Omega })`$. If the second term in eq.(12) were absent, the density is given by $`\rho =I/(\pi A^2)`$, where $`I`$ is the total number of Landau levels below the chemical potential, $`I=\mathrm{Int}[\frac{\mu +\mathrm{}\mathrm{\Omega }}{2\mathrm{}\mathrm{\Omega }}]`$. When $`\mathrm{\Omega }`$ is close to $`\omega _{}`$, the second term in eq.(12) is slowly varying over the scale of $`Aa_{}`$ and can be absorbed in the chemical potential. The density profile within local density approximation (LDA) is then
$$\rho (r)=\frac{M\mathrm{\Omega }}{\pi \mathrm{}}I(r),I(r)=\mathrm{Int}\left[\frac{\mu (r)+\mathrm{}\mathrm{\Omega }}{2\mathrm{}\mathrm{\Omega }}\right]$$
(13)
where $`\mu (r)=\mu \frac{1}{2}M(\omega _{}^2\mathrm{\Omega }^2)r^2`$. Clearly, eq.(13) is equivalent to (11) up to correction $`(1\mathrm{\Omega }/\omega _{})`$$`<<1`$ as $`\mathrm{\Omega }\omega _{}`$. Eq.(13), once established, is easily generalized to other potentials. In the presence of $`V_{\mathrm{wall}}`$, one simply replaces $`\mu (r)`$ in eq.(13) by
$$\mu (r)=\mu \frac{1}{2}M(\omega _{}^2\mathrm{\Omega }^2)r^2V_{\mathrm{wall}}(r).(2D)$$
(14)
Eq.(13) and eq.(14) constitute the LDA solution for the 2D rotating Fermi gas for both $`\mathrm{\Omega }<\omega _{}`$ and $`\mathrm{\Omega }>\omega _{}`$. The schematics of LDA is shown in fig.1(a) and 1(b).
To understand the validity of LDA (eq.(13)), we have calculated the density numerically using eq.(2). The result for a system of $`2000`$ Fermions at $`\mathrm{\Omega }/\omega _{}=0.996`$ is shown in figure 2a. The system exhibits a sequence of quantized steps at locations well described by LDA. The evolution of the density within the range $`0.99<\mathrm{\Omega }/\omega _{}<1`$ of this Fermion system is shown in fig.2b. As $`\mathrm{\Omega }`$ decreases, more Landau levels are populated while the steps near the surfaces are closer together, (as expected from eq.(6) and (8)), yet the step structures remain discernable and correctly described by LDA, (the LDA construction is not displayed so as to keep the fig.2b readable). The behaviors of the densities at lower frequency $`0.98<\mathrm{\Omega }/\omega _{}<1`$ are shown in fig.3a and 3b for a system of 1000 Fermions. At the lowest frequency displayed, (fig.3a), the step structure near the surface is completely smeared out by the spread of the edges. Nevertheless, the density of the innermost plateau remained quantized, with a size correctly described by LDA. This core of quantized density (or “quantized core” for short) is a clear evidence for Landau levels. Our studies show that for about 2000 particles, Landau levels will show up as a sequence of discernable steps only when $`0.99<\mathrm{\Omega }/\omega _{}<1`$, which is achievable with the current capability to control frequencies, especially for large $`\omega _{}`$. For lower frequencies, the existence of Landau levels can only be revealed through the presence of a “quantized core”, which shrinks in size as $`\mathrm{\Omega }`$ decreases. On the other hand, the LDA in fig.1b shows that by introducing an additional potential $`V_{\mathrm{wall}}`$, Landau levels (in the form of a sequence of steps near the center or a “quantized core”) can still be revealed at frequencies farther beyond $`\omega _{}`$, even though the steps near the surface are smeared out.
3D case : For a 3D harmonic trap, eq.(1) acquires terms $`p_z^2/2M+\frac{1}{2}M\omega _zz^2`$, which give rise to harmonic oscillator eigenfunctions $`f_{n_z}(z)`$ with eigenvalue $`ϵ_{n_z}=\mathrm{}\omega _z(n_z+\frac{1}{2})`$. The density is $`\rho (r,z)`$$`=_{n_z,n,m}`$$`|f_{n_z}(z)|^2`$$`|u_{n,m}(𝐫)|^2`$ $`\mathrm{\Theta }(\mu ϵ_{n_z,n,m})`$, where $`ϵ_{n_z,n,m}=ϵ_{n_z}+ϵ_{n,m}`$. When $`\mathrm{\Omega }/\omega _{}1`$, we first perform the $`m`$-sum because it produces the smoothest change in the energy. Repeating the steps leading to eq.(10), we have
$$\rho (r,z)=\underset{n_z,n}{}\frac{|f_{n_z}(z)|^2}{\pi a_{}^2}\mathrm{\Theta }(\mu (r)\mathrm{}\mathrm{\Omega }(2n+1)ϵ_{n_z})$$
(15)
where $`\mu (r)=\mu M\omega _{}(\omega _{}\mathrm{\Omega })r^2\mu \frac{M}{2}(\omega _{}^2\mathrm{\Omega }^2)r^2`$.
The order of summation of the remaining integers $`n_z`$ and $`n`$ depends on the relative strength between $`\omega _z`$ and $`\omega _{}`$. When $`\omega _z<<\omega _{}`$, we first sum $`n_z`$. To do that, we note that the density of a 1D Fermi gas is $`\rho _{1D}=\sqrt{2M\mu }/(\pi \mathrm{})`$, where $`\mu `$ is the chemical potential. In a harmonic trap, the density within LDA is obtained by changing $`\mu `$ to $`\mu (z)=\mu \frac{1}{2}M\omega _z^2z^2`$. We then have $`\rho _{_{1D}}(z)=_{n_z}|f_{n_z}(z)|^2\mathrm{\Theta }(\mu ϵ_{n_z})=\frac{\sqrt{2M\mu (z)}}{\pi \mathrm{}}`$, which turns eq.(15) into a useful LDA
$$\rho (r)=\frac{1}{\pi ^2}\frac{1}{a_{}^2a_z}\underset{n}{}\sqrt{\left[\mu (r,z)\mathrm{}\mathrm{\Omega }(2n+1)\right]/(\frac{1}{2}\mathrm{}\omega _z)},$$
(16)
where $`a_z=\sqrt{\mathrm{}/M\omega _z}`$, and $`\mu (r,z)=\mu \frac{M}{2}(\omega _{}^2\mathrm{\Omega }^2)r^2\frac{M}{2}\omega _z^2z^2V_{\mathrm{wall}}`$. We have included $`V_{\mathrm{wall}}`$ in $`\mu `$ for the more general situation as in the 2D case. Eq.(16) shows that $`\rho (r)`$ is a sum of 1D densities (labelled by “$`n`$”) each of which distributed over in a “Landau volume” bounded by the “Landau surface” $`\mu (r,z)=\mathrm{}\mathrm{\Omega }(2n+1)`$. When $`V_{\mathrm{wall}}=0`$, $`\mathrm{\Omega }<\omega _{}`$ the Landau surface are ellipsoidal surfaces. It is easy to verify that the surface areas $`A_n`$ for successive ellipsoids differ by a constant $`A_{n1}A_n=\frac{\mathrm{}}{M\mathrm{\Omega }}\left(\frac{16\pi \mathrm{\Omega }^2}{\omega _{}^2\mathrm{\Omega }^2}\right)`$. When $`\mathrm{\Omega }\omega _{}`$, centrifugal instability against harmonic confinement sets in and stability can only be established by $`V_{\mathrm{wall}}`$.
To demonstrate the validity of the LDA eq.(16), we have evaluated the density numerically for a system of 2000 Fermions for $`\omega _z/\omega _{}=0.2`$ at $`\mathrm{\Omega }/\omega _{}=0.99`$. The results are shown in figure 4. It shows that the LDA (dotted line) is a good approximation. The Landau volumes can be clearly identified by the change of slope in the density. The appearance of a plateau the center is because $`\omega _z/\omega _{}`$ is only 0.2, revealing the 2D feature of the $`n_z`$ levels. For smaller ratios of $`\omega _z/\omega _{}`$, the plateau disappears and the LDA expression (eq.(16)) is achieved.
When $`\omega _z>\omega _{}`$, summation of $`n`$ in eq.(15) gives
$$\rho (r)=\frac{M\omega _{}}{\pi \mathrm{}}\underset{n_z}{}|f_{n_z}(z)|^2\mathrm{Int}\left[\frac{\mu (r)\mathrm{}\omega _z(n_z+\frac{1}{2})+\mathrm{}\mathrm{\Omega }}{2\mathrm{}\mathrm{\Omega }}\right]$$
(17)
where $`\mu (r)=\mu \frac{1}{2}M(\omega _{}^2\mathrm{\Omega }^2)r^2V_{\mathrm{wall}}`$. In this limit, the density consists of a sequence of discs (labeled by $`n_z`$) in the $`z`$-direction. Each disc $`|f_{n_z}(z)|^2`$ consists of a sequence of density steps in the $`xy`$-plane reflecting the number of filled Landau levels. The behavior of the density within each disc in the $`xy`$-plane is identical to the 2D case discussed before. Finally, we note that as temperature increases, the Landau level structure near the surface will first melt away, and the melting will proceed toward the center. The temperature below which the Landau level effect begin to appear is $`T=2\mathrm{}\omega _{}/k_B`$, which is $`3.8\times 10^7`$K and $`9.6\times 10^6`$K for for $`\omega _{}=4000`$Hz and $`10^5`$Hz resp., a temperature range achievable in current experiments.
So far, we have only discussed the effect of Landau levels on the density profiles of fast rotating Fermi gases. If the development of quantum Hall effect in the last decade is a guide, one expects many more novel phenomena in Fermi gases in the fast rotating regime. This work is completed during a workshop at the Lorentz Center of University of Leiden. We thank Professor Henk Stoof and the Lorentz Center for generous support. This work is supported by a Grant from NASA NAG8-1441, and by the NSF Grants DMR-9705295 and DMR-9807284.
Caption
Fig.1a and 1b: Fig.1a and 1(b) corresponds to $`\mathrm{\Omega }<\omega _{}`$ and $`\mathrm{\Omega }>\omega _{}`$ resp. The rapid drop at large $`r`$ is due to the strongly repulsive potential $`V_{\mathrm{wall}}`$. The LDA densities are indicated by the steps in thick lines. The integer value of $`(\mu (r)+\mathrm{}\mathrm{\Omega })/(2\mathrm{}\mathrm{\Omega })`$, (i.e. $`I`$), is related to the index $`n`$ of the intersected Landau level as $`I=n+1`$. And the relation $`(\mu (r)+\mathrm{}\mathrm{\Omega })/(2\mathrm{}\mathrm{\Omega })=I`$ is equivalent to $`\mu (r)=(2n+1)\mathrm{\Omega }`$.
Figure.2a: LDA (dotted line) and numerical calculations (solid line) of the density for $`N=2000`$ Fermions at $`\mathrm{\Omega }/\omega _{}`$$`=0.996`$.
Figure 2b. Density profiles of $`N=2000`$ Fermions at $`\mathrm{\Omega }/\omega _{}`$$`=0.993`$ (crosses), $`0.996`$ (dots), and 0.998 (solid line).
Figure 3a: LDA (dotted line) and numerical calculations (solid line) of the density for $`N=1000`$ Fermions at $`\mathrm{\Omega }/\omega _{}`$$`=0.982`$.
Figure 3b. Density profiles of $`N=2000`$ Fermions at $`\mathrm{\Omega }/\omega _{}`$$`=0.982`$ (crosses), $`0.989`$ (dots), and 0.995 (solid line).
Figure 4. LDA (dotted line) and numerical calculations (solid line) of the density for $`N=2000`$ Fermions at $`\mathrm{\Omega }/\omega _{}`$$`=0.99`$ and $`\omega _z/\omega _{}`$$`=0.2`$.
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# Problems on the geometry of finitely generated solvable groups
## 1. Introduction
Our story begins with a theorem of Gromov, proved in 1980.
###### Theorem 1 (Gromov’s Polynomial Growth Theorem \[Gr1\]).
Let $`G`$ be any finitely generated group. If $`G`$ has polynomial growth then $`G`$ is virtually nilpotent, i.e. $`G`$ has a finite index nilpotent subgroup.
Gromov’s theorem inspired the more general problem (see, e.g. \[GH1, BW, Gr1, Gr2\]) of understanding to what extent the asymptotic geometry of a finitely generated solvable group determines its algebraic structure. One way in which to pose this question precisely is via the notion of quasi-isometry.
A (coarse) quasi-isometry between metric spaces is a map $`f:XY`$ such that, for some constants $`K,C,C^{}>0`$:
1. $`\frac{1}{K}d_X(x_1,x_2)Cd_Y(f(x_1),f(x_2))Kd_X(x_1,x_2)+C`$ for all $`x_1,x_2X`$.
2. The $`C^{}`$-neighborhood of $`f(X)`$ is all of $`Y`$.
$`X`$ and $`Y`$ are *quasi-isometric* if there exists a quasi-isometry $`XY`$. Note that the quasi-isometry type of a metric space $`X`$ is unchanged upon removal of any bounded subset of $`X`$; hence the term “asymptotic”.
Quasi-isometries are the natural maps to study when one is interested in the geometry of a group. In particular:
* The word metric on any f.g. group is unique up to quasi-isometry.
* Any injective homomorphism with finite index image is a quasi-isometry, as is any surjective homomorphism with finite kernel. The equivalence relation generated by these two types of maps can be described more compactly: two groups $`G,H`$ are equivalent in this manner if and only if they are *weakly commensurable*, which means that there exists a group $`Q`$ and homomorphisms $`QG`$, $`QH`$ each having finite kernel and finite index image (proof: show that “weakly commensurable” is in fact an equivalence relation). This weakens the usual notion of commensurability, i.e. when $`G`$ and $`H`$ have isomorphic finite index subgroups. Weakly commensurable groups are clearly quasi-isometric.
* Any two cocompact, discrete subgroups of a Lie group are quasi-isometric. There are cocompact discrete subgroups of the same Lie group which are not weakly commensurable, for example arithmetic lattices are not weakly commensurable to non-arithmetic ones.
The Polynomial Growth Theorem was an important motivation for Gromov when he initiated in \[Gr2, Gr3\] the problem of classifying finitely-generated groups up to quasi-isometry.
Theorem 1, together with the fact that nilpotent groups have polynomial growth (see §3 below), implies that the property of being nilpotent is actually an asymptotic property of groups. More precisely, the class of nilpotent groups is quasi-isometrically rigid: any finitely-generated group quasi-isometric to a nilpotent group is weakly commensurable to some nilpotent group <sup>1</sup><sup>1</sup>1In fact such a group must have a finite-index nilpotent subgroup. In general, weak commensurability is the most one can hope for in quasi-isometric rigidity problems.. Sometimes this is expressed by saying that the property of being nilpotent is a geometric property, i.e. it is a quasi-isometry invariant (up to weak commensurability). The natural question then becomes:
###### Question 2 (Rigidity question).
Which subclasses of f.g. solvable groups are quasi-isometrically rigid? For example, are polycyclic groups quasi-isometrically rigid? metabelian groups? nilpotent-by-cyclic groups?
In other words, which of these algebraic properties of a group are actually geometric, and are determined by apparently cruder asymptotic information? A. Dioubina \[Di\] has recently found examples which show that the class of finitely generated solvable groups is not quasi-isometrically rigid (see §2). On the other hand, at least some subclasses of solvable groups are indeed rigid (see §4).
Along with Question 2 comes the finer classification problem:
###### Problem 3 (Classification problem).
Classify f.g. solvable (resp. nilpotent, polycyclic, metabelian, nilpotent-by-cyclic, etc.) groups up to quasi-isometry.
As we shall see, the classification problem is usually much more delicate than the rigidity problem; indeed the quasi-isometry classification of finitely-generated nilpotent groups remains one of the major open problems in the field. We discuss this in greater detail in §3.
The corresponding rigidity and classification problems for irreducible lattices in semisimple Lie groups have been completely solved. This is a vast result due to many people and, among other things, it generalizes and strengthens the Mostow Rigidity Theorem. We refer the reader to \[Fa\] for a survey of this work.
In contrast, results for finitely generated solvable groups have remained more elusive. There are several reasons for this:
* Finitely generated solvable groups are defined algebraically, and so they do not always come equipped with an obvious or well-studied geometric model (see, e.g., item 3 below).
* Dioubina’s examples show not only that the class of finitely-generated solvable groups is not quasi-isometrically rigid; they also show (see §4 below) that the answer to Question 2 for certain subclasses of solvable groups (e.g. abelian-by-cyclic) differs in the finitely presented and finitely generated cases.
* There exists a finitely presented solvable group $`\mathrm{\Gamma }`$ of derived length 3 with the property that $`\mathrm{\Gamma }`$ has unsolvable word problem (see \[Kh\]). Solving the word problem for a group is equivalent to giving an algorithm to build the Cayley graph of that group. In this sense there are finitely presented solvable groups whose geometry cannot be understood, at least by a Turing machine.
* Solvable groups are much less rigid than irreducible lattices in semisimple Lie groups. This phenomenon is exhibited concretely by the fact that many finitely generated solvable groups have infinite-dimensional groups of self quasi-isometries<sup>2</sup><sup>2</sup>2The set of self quasi-isometries of a space $`X`$, with the operation of composition, becomes a group $`\mathrm{QI}(X)`$ once one mods out by the relation $`fg`$ if $`d(f,g)<\mathrm{}`$ in the $`sup`$ norm. (see below).
###### Problem 4 (Flexibility of solvable groups).
For which infinite, finitely generated solvable groups $`\mathrm{\Gamma }`$ is $`\mathrm{QI}(\mathrm{\Gamma })`$ infinite dimensional?
In contrast, all irreducible lattices in semisimple Lie groups $`G\mathrm{SO}(n,1),\mathrm{SU}(n,1)`$ have countable or finite-dimensional quasi-isometry groups.
At this point in time, our understanding of the geometry of finitely-generated solvable groups is quite limited. In §3 we discuss what is known about the quasi-isometry classification of nilpotent groups (the rigidity being given by Gromov’s Polynomial Growth Theorem). Beyond nilpotent groups, the only detailed knowledge we have is for the finitely-presented, nonpolycyclic abelian-by-cyclic groups. We discuss this in depth in §4, and give a conjectural picture of the polycyclic case in §5. One of the interesting discoveries described in these sections is a connection between finitely presented solvable groups and the theory of dynamical systems. This connection is pursued very briefly in a more general context in §6, together with some questions about issues beyond the limits of current knowledge.
This article is meant only as a brief survey of problems, conjectures, and theorems. It therefore contains neither an exhaustive history nor detailed proofs; for these the reader may consult the references. It is a pleasure to thank David Fisher, Pierre de la Harpe, Ashley Reiter, Jennifer Taback, and the referee for their comments and corrections.
## 2. Dioubina’s examples
Recall that the wreath product of groups $`A`$ and $`B`$, denoted $`A\mathrm{}B`$, is the semidirect product $`(_AB)A`$, where $`_AB`$ is the direct sum of copies of B indexed by elements of A, and A acts via the “shift”, i.e. the left action of $`A`$ on the index set $`A`$ via left multiplication. Note that if $`A`$ and $`B`$ are finitely-generated then so is $`A\mathrm{}B`$.
The main result of Dioubina \[Di\] is that, if there is a bijective quasi-isometry between finitely-generated groups $`A`$ and $`B`$, then for any finitely-generated group $`C`$ the groups $`C\mathrm{}A`$ and $`C\mathrm{}B`$ are quasi-isometric. Dioubina then applies this theorem to the groups $`A=C=𝐙,B=𝐙D`$ where $`D`$ is a finite nonsolvable group. It is easy to construct a one-to-one quasi-isometry between $`A`$ and $`B`$. Hence $`G=C\mathrm{}A`$ and $`H=C\mathrm{}B`$ are quasi-isometric.
Now $`G`$ is torsion-free solvable, in fact $`G=𝐙\mathrm{}𝐙`$ is an abelian-by-cyclic group of the form $`𝐙[𝐙]`$-by-$`𝐙`$. On the other hand $`H`$ contains $`_𝐙D`$, and so is not virtually solvable, nor even weakly commensurable with a solvable group. Hence the class of finitely-generated solvable groups is not quasi-isometrically rigid.
Dioubina’s examples never have any finite presentation. In fact if $`A\mathrm{}B`$ is finitely presented then either $`A`$ or $`B`$ is finite (see \[Bau\]). This leads to the following question.
###### Question 5.
Is the class of finitely presented solvable groups quasi-isometrically rigid?
Note that the property of being finitely presented is a quasi-isometry invariant (see \[GH1\]).
## 3. Nilpotent groups and Pansu’s Theorem
While the Polynomial Growth Theorem shows that the class of finitely generated nilpotent groups is quasi-isometrically rigid, the following remains an important open problem.
###### Problem 6.
Classify finitely generated nilpotent groups up to quasi-isometry.
The basic quasi-isometry invariants for a finitely-generated nilpotent group $`G`$ are most easily computed in terms of the set $`\{c_i(G)\}`$ of ranks (over $`𝐐`$) of the quotients $`G_i/G_{i+1}`$ of the lower central series for $`G`$, where $`G_i`$ is defined inductively by $`G_1=G`$ and $`G_{i+1}=[G,G_i]`$.
One of the first quasi-isometry invariants to be studied was the growth of a group, studied by Dixmier, Guivar’ch, Milnor, Wolf, and others (see \[H\], Chapters VI-VII for a nice discussion of this, and a careful account of the history). The growth of $`G`$ is the function of $`r`$ that counts the number of elements in a ball of radius $`r`$ in $`G`$. There is an important dichotomy for solvable groups:
###### Theorem 7 (\[Wo, Mi\]).
Let $`G`$ be a finitely generated solvable group. Then either $`G`$ has polynomial growth and is virtually nilpotent, or $`G`$ has exponential growth and is not virtually nilpotent.
When $`G`$ has polynomial growth, the degree $`\mathrm{deg}(G)`$ of this polynomial is easily seen to be a quasi-isometry invariant. It is given by the following formula, discovered around the same time by Guivar’ch \[Gu\] and by Bass \[Ba\]:
$$\mathrm{deg}(G)=\underset{i=1}{\overset{n}{}}ic_i(G)$$
where $`n`$ is the degree of nilpotency of $`G`$. Another basic invariant is that of virtual cohomological dimension $`\mathrm{vcd}(G)`$. For groups $`G`$ with finite classifying space (which is not difficult to check for torsion-free nilpotent groups), this number was shown by Gersten \[Ge\] and Block-Weinberger \[BW\] to be a quasi-isometry invariant. On the other hand it is easy to check that
$$\mathrm{vcd}(G)=\underset{i=1}{\overset{n}{}}c_i(G)$$
where $`n`$ is the degree of nilpotency, also known as the Hirsch length, of $`G`$. As Bridson and Gersten have shown (see \[BG, Ge\]), the above two formulas imply that any finitely generated group $`\mathrm{\Gamma }`$ which is quasi-isometric to $`𝐙^n`$ must have a finite index $`𝐙^n`$ subgroup: by the Polynomial Growth Theorem such a $`\mathrm{\Gamma }`$ has a finite index nilpotent subgroup $`N`$; but
$$\mathrm{deg}(N)=\mathrm{deg}(𝐙^n)=n=\mathrm{vcd}(𝐙^n)=\mathrm{vcd}(N)$$
and so
$$\underset{i}{}ic_i(N)=\underset{i}{}c_i(N)$$
which can only happen if $`c_i(N)=0`$ for $`i>1`$, in which case $`N`$ is abelian.
###### Problem 8 (\[GH2\]).
Give an elementary proof (i.e. without using Gromov’s Polynomial Growth Theorem) that any finitely generated group quasi-isometric to $`𝐙^n`$ has a finite index $`𝐙^n`$ subgroup.
As an exercise, the reader is invited to find nilpotent groups $`N_1,N_2`$ which are not quasi-isometric but which have the same degree of growth and the same $`\mathrm{vcd}`$.
There are many other quasi-isometry invariants for finitely-generated nilpotent groups $`\mathrm{\Gamma }`$. All known invariants are special cases of the following theorem of Pansu \[Pa1\]. To every nilpotent group $`\mathrm{\Gamma }`$ one can associate a nilpotent Lie group $`\mathrm{\Gamma }𝐑`$, called the Malcev completion of $`\mathrm{\Gamma }`$ (see \[Ma\]), as well as the associated graded Lie group $`\mathrm{gr}(\mathrm{\Gamma }𝐑)`$.
###### Theorem 9 (Pansu’s Theorem \[Pa1\]).
Let $`\mathrm{\Gamma }_1,\mathrm{\Gamma }_2`$ be two finitely-generated nilpotent groups. If $`\mathrm{\Gamma }_1`$ is quasi-isometric to $`\mathrm{\Gamma }_2`$ then $`\mathrm{gr}(\mathrm{\Gamma }_1𝐑)`$ is isomorphic to $`\mathrm{gr}(\mathrm{\Gamma }_2𝐑)`$.
We remark that there are nilpotent groups with non-isomorphic Malcev completions where the associated gradeds are isomorphic; the examples are 7-dimensional and somewhat involved (see \[Go\], p.24, Example 2). It is not known whether or not the Malcev completion is a quasi-isometry invariant.
Theorem 9 immediately implies:
###### Corollary 10.
The numbers $`c_i(\mathrm{\Gamma })`$ are quasi-isometry invariants.
In particular we recover (as special cases) that growth and cohomological dimension are quasi-isometry invariants of $`\mathrm{\Gamma }`$.
To understand Pansu’s proof one must consider Carnot groups. These are graded nilpotent Lie groups $`N`$ whose Lie algebra $`𝐍`$ is generated (via bracket) by elements of degree one. Chow’s Theorem \[Ch\] states that such Lie groups $`N`$ have the property that the left-invariant distribution obtained from the degree one subspace $`𝐍^1`$ of $`𝐍`$ is a totally nonintegrable distribution: any two points $`x,yN`$ can be connected by a piecewise smooth path $`\gamma `$ in $`N`$ for which the vector $`d\gamma /dt(t)`$ lies in the distribution. Infimizing the length of such paths between two given points gives a metric on $`N`$, called the Carnot Carethéodory metric $`d_{car}`$. This metric is non-Riemannian if $`N𝐑^n`$. For example, when $`N`$ is the $`3`$-dimensional Heisenberg group then the metric space $`(N,d_{car})`$ has Hausdorff dimension $`4`$.
One important property of Carnot groups is that they come equipped with a $`1`$-parameter family of dilations $`\{\delta _t\}`$, which gives a notion of (Carnot) differentiability (see \[Pa1\]). Further, the differential $`Df(x)`$ of a map $`f:N_1N_2`$ between Carnot groups $`N_1,N_2`$ which is (Carnot) differentiable at the point $`xN_1`$ is actually a Lie group homomorphism $`N_1N_2`$.
Sketch of Pansu’s proof of Theorem 9. If $`(\mathrm{\Gamma },d)`$ is a nilpotent group endowed with a word metric $`d`$, the sequence of scaled metric spaces $`\{(\mathrm{\Gamma },\frac{1}{n}d)\}_{n𝐍}`$ has a limit in the sense of Gromov-Hausdorff convergence:
$$(\mathrm{\Gamma }_{\mathrm{}},d_{\mathrm{}})=\underset{n\mathrm{}}{lim}(\mathrm{\Gamma },\frac{1}{n}d)$$
(See \[Pa2\], and \[BrS\] for an introduction to Gromov-Hausdorff convergence). It was already known, using ultralimits, that some subsequence converges \[Gr1, DW\]. Pansu’s proof not only gives convergence on the nose, but it yields some additional important features of the limit metric space $`(\mathrm{\Gamma }_{\mathrm{}},d_{\mathrm{}})`$:
* (Identifying limit) It is isometric to the Carnot group $`\mathrm{gr}(\mathrm{\Gamma }𝐑)`$ endowed with the Carnot metric $`d_{car}`$.
* (Functoriality) Any quasi-isometry $`f:\mathrm{\Gamma }_1\mathrm{\Gamma }_2`$ bewteen finitely-generated nilpotent groups induces a bilipschitz homeomorphism
$$\widehat{f}:(\mathrm{gr}(\mathrm{\Gamma }_1𝐑),d_{car})(\mathrm{gr}(\mathrm{\Gamma }_2𝐑),d_{car})$$
Note that functoriality follows immediately once we know the limit exists: the point is that if $`f:(\mathrm{\Gamma }_1,d_1)(\mathrm{\Gamma }_2,d_2)`$ is a $`(K,C)`$ quasi-isometry of word metrics, then for each $`n`$ the map
$$f:(\mathrm{\Gamma }_1,\frac{1}{n}d_1)(\mathrm{\Gamma }_2,\frac{1}{n}d_2)$$
is a $`(K,C/n)`$ quasi-isometry, hence the induced map
$$\widehat{f}:((\mathrm{\Gamma }_1)_{\mathrm{}},d_{car})((\mathrm{\Gamma }_2)_{\mathrm{}},d_{car})$$
is a $`(K,0)`$ quasi-isometry, i.e. is a bilipschitz homeomorphism.
Given a quasi-isometry $`f:\mathrm{\Gamma }_1\mathrm{\Gamma }_2`$, we thus have an induced bilipschitz homeomorphism $`\widehat{f}:\mathrm{gr}(\mathrm{\Gamma }_1𝐑)\mathrm{gr}(\mathrm{\Gamma }_2𝐑)`$ between Carnot groups endowed with Carnot-Carethéodory metrics. Pansu then proves a regularity theorem, generalizing the Rademacher-Stepanov Theorem for $`𝐑^n`$. This general regularity theorem states that a bilipschitz homeomorphism of Carnot groups (endowed with Carnot-Carethéodory metrics) is differentiable almost everywhere. Since the differential $`D_x\widehat{f}`$ is actually a group homomorphism, we know that for almost every point $`x`$ the differential $`D_x\widehat{f}:\mathrm{gr}(\mathrm{\Gamma }_1𝐑)\mathrm{gr}(\mathrm{\Gamma }_2𝐑)`$ is an isomorphism.
## 4. Abelian-by-cyclic groups: nonpolycyclic case
The first progress on Question 2 and Problem 3 in the non-(virtually)-nilpotent case was made in \[FM1\] and \[FM2\]. These papers proved classification and rigidity for the simplest class of non-nilpotent solvable groups: the solvable Baumslag-Solitar groups
$$\mathrm{BS}(1,n)=<a,b:aba^1=b^n>$$
These groups are part of the much broader class of abelian-by-cyclic groups. A group $`\mathrm{\Gamma }`$ is abelian-by-cyclic if there is an exact sequence
$$1A\mathrm{\Gamma }Z1$$
where $`A`$ is an abelian group and $`Z`$ is an infinite cyclic group. If $`\mathrm{\Gamma }`$ is finitely generated, then $`A`$ is a finitely generated module over the group ring $`𝐙[Z]`$, although $`A`$ need not be finitely generated as a group.
By a result of Bieri and Strebel \[BS1\], the class of finitely presented, torsion-free, abelian-by-cyclic groups may be described in another way. Consider an $`n\times n`$ matrix $`M`$ with integral entries and $`detM0`$. Let $`\mathrm{\Gamma }_M`$ be the ascending HNN extension of $`𝐙^n`$ given by the monomorphism $`\varphi _M`$ with matrix $`M`$. Then $`\mathrm{\Gamma }_M`$ has a finite presentation
$$\mathrm{\Gamma }_M=t,a_1,\mathrm{},a_n|[a_i,a_j]=1,ta_it^1=\varphi _M(a_i),i,j=1,\mathrm{},n$$
where $`\varphi _M(a_i)`$ is the word $`a_1^{m_1}\mathrm{}a_n^{m_n}`$ and the vector $`(m_1,\mathrm{},m_n)`$ is the $`i^{\text{th}}`$ column of the matrix $`M`$. Such groups $`\mathrm{\Gamma }_M`$ are precisely the class of finitely presented, torsion-free, abelian-by-cyclic groups (see \[BS1\] for a proof involving a precursor of the Bieri-Neumann-Strebel invariant, or \[FM2\] for a proof using trees). The group $`\mathrm{\Gamma }_M`$ is polycyclic if and only if $`\left|detM\right|=1`$ (see \[BS2\]).
The results of \[FM1\] and \[FM2\] are generalized in \[FM3\], which gives the complete classification of the finitely presented, nonpolycyclic abelian-by-cyclic groups among all f.g. groups, as given by the following two theorems.
The first theorem in \[FM3\] gives a classification of all finitely-presented, nonpolycyclic, abelian-by-cyclic groups up to quasi-isometry. It is easy to see that any such group has a torsion-free subgroup of finite index, so is commensurable (hence quasi-isometric) to some $`\mathrm{\Gamma }_M`$. The classification of these groups is actually quite delicate—the standard quasi-isometry invariants (ends, growth, isoperimetric inequalities, etc.) do not distinguish any of these groups from each other, except that the size of the matrix $`M`$ can be detected by large scale cohomological invariants of $`\mathrm{\Gamma }_M`$.
Given $`M\mathrm{GL}(n,𝐑)`$, the *absolute Jordan form* of $`M`$ is the matrix obtained from the Jordan form for $`M`$ over $`𝐂`$ by replacing each diagonal entry with its absolute value, and rearranging the Jordan blocks in some canonical order.
###### Theorem 11 (Nonpolycyclic, abelian-by-cyclic groups: Classification).
Let $`M_1`$ and $`M_2`$ be integral matrices with $`\left|detM_i\right|>1`$ for $`i=1,2`$. Then $`\mathrm{\Gamma }_{M_1}`$ is quasi-isometric to $`\mathrm{\Gamma }_{M_2}`$ if and only if there are positive integers $`r_1,r_2`$ such that $`M_1^{r_1}`$ and $`M_2^{r_2}`$ have the same absolute Jordan form.
###### Remark.
Theorem 11 generalizes the main result of \[FM1\], which is the case when $`M_1,M_2`$ are positive $`1\times 1`$ matrices; in that case the result of \[FM1\] says even more, namely that $`\mathrm{\Gamma }_{M_1}`$ and $`\mathrm{\Gamma }_{M_2}`$ are quasi-isometric if and only if they are commensurable. When $`n2`$, however, it’s not hard to find $`n\times n`$ matrices $`M_1,M_2`$ such that $`\mathrm{\Gamma }_{M_1},\mathrm{\Gamma }_{M_2}`$ are quasi-isometric but not commensurable. Polycyclic examples are given in \[BG\]; similar ideas can be used to produce nonpolycyclic examples.
The following theorem shows that the algebraic property of being a finitely presented, nonpolycyclic, abelian-by-cyclic group is in fact a geometric property.
###### Theorem 12 (Nonpolycyclic, abelian-by-cyclic groups: Rigidity).
Let $`\mathrm{\Gamma }=\mathrm{\Gamma }_M`$ be a finitely presented abelian-by-cyclic group, determined by an $`n\times n`$ integer matrix $`M`$ with $`\left|detM\right|>1`$. Let $`G`$ be any finitely generated group quasi-isometric to $`\mathrm{\Gamma }`$. Then there is a finite normal subgroup $`NG`$ such that $`G/N`$ is commensurable to $`\mathrm{\Gamma }_N`$, for some $`n\times n`$ integer matrix $`N`$ with $`\left|detN\right|>1`$.
###### Remark.
Theorem 12 generalizes the main result of \[FM2\], which covers the case when $`M`$ is a positive $`1\times 1`$ matrix. The $`1\times 1`$ case is given a new proof in \[MSW\], which is adapted in \[FM3\] to prove Theorem 12.
###### Remark.
The “finitely presented” hypothesis in Theorem 12 cannot be weakened to “finitely generated”, since Diuobina’s example (discussed in §2) is abelian-by-cyclic, namely $`𝐙[𝐙]`$-by-$`𝐙`$.
One new discovery in \[FM3\] is that there is a strong connection between the geometry of solvable groups and the theory of dynamical systems. Assuming here for simplicity that the matrix $`M`$ lies on a $`1`$-parameter subgroup $`M^t`$ in $`\mathrm{GL}(n,𝐑)`$, let $`G_M`$ be the semi-direct product $`𝐑^n_M𝐑`$, where $`𝐑`$ acts on $`𝐑^n`$ by the $`1`$-parameter subgroup $`M^t`$. We endow the solvable Lie group $`G_M`$ with a left-invariant metric. The group $`G_M`$ admits a vertical flow:
$$\mathrm{\Psi }_s(x,t)=(x,t+s)$$
There is a natural horizontal foliation of $`G_M`$ whose leaves are the level sets $`P_t=𝐑^n\times \{t\}`$ of time. A quasi-isometry $`f:G_MG_N`$ is horizontal respecting if it coarsely permutes the leaves of this foliation; that is, if there is a constant $`C0`$ so that
$$d_{}(f(P_t),P_{h(t)})C$$
where $`d_{}`$ denotes Hausdorff distance and $`h:𝐑𝐑`$ is some function, which we think of as a time change between the flows.
A key technical result of \[FM3\] is the phenomenon of time rigidity: the time change $`h`$ must actually be affine, so taking a real power of $`M`$ allows one to assume $`h(t)=t`$.
It is then shown that “quasi-isometries remember the dynamics”. That is, $`f`$ coarsely respects several foliations arising from the partially hyperbolic dynamics of the flow $`\mathrm{\Psi }`$, starting with the weak stable, weak unstable, and center-leaf foliations. By keeping track of different exponential and polynomial divergence properties of the action of $`\mathrm{\Psi }`$ on tangent vectors, the weak stable and weak unstable foliations are decomposed into flags of foliations. Using time rigidity and an inductive argument it is shown that these flags are coarsely respected by $`f`$ as well. Relating the flags of foliations to the Jordan Decomposition then completes the proof of:
###### Theorem 13 (Horizontal respecting quasi-isometries).
If there is a horizontal-respecting quasi-isometry $`f:G_MG_N`$ then there exist nonzero $`a,b𝐑`$ so that $`M^a`$ and $`M^b`$ have the same absolute Jordan form.
The “nonpolycyclic” hypothesis (i.e. $`\left|detM\right|>1`$) in Theorem 11 is used in two ways. First, the group $`\mathrm{\Gamma }_M`$ has a model space which is topologically a product of $`𝐑^n`$ and a regular tree of valence $`\left|detM\right|+1`$, and when this valence is greater than $`2`$ we can use coarse algebraic topology (as developed in \[FS\], \[EF\], and \[FM1\]) to show that any quasi-isometry $`\mathrm{\Gamma }_M\mathrm{\Gamma }_N`$ induces a quasi-isometry $`G_MG_N`$ satisfying the hypothesis of Theorem 13. Second, we are able to pick off integral $`a,b`$ by developing a “boundary theory” for $`\mathrm{\Gamma }_M`$; in case $`\left|detM\right|>1`$ this boundary is a self-similar Cantor set whose bilipschitz geometry detects the primitive integral power of $`\left|detM\right|`$ by Cooper’s Theorem \[FM1\], finishing the proof of Theorem 11.
###### Problem 14 (Nilpotent-by-cyclic groups).
Extend Theorem 11 and Theorem 12 to the class of finitely-presented nilpotent-by-cyclic groups.
Of course, as the classification of finitely-generated nilpotent groups is still open, Problem 14 is meant in the sense of reducing the nilpotent-by-cyclic case to the nilpotent case, together with another invariant. This second invariant for a nilpotent-by-cyclic group $`G`$ will perhaps be the absolute Jordan form of the matrix which is given by the action of the generator of the cyclic quotient of $`G`$ on the nilpotent kernel of $`G`$.
## 5. Abelian-by-cyclic groups: polycyclic case
The polycyclic, abelian-by-cyclic groups are those $`\mathrm{\Gamma }_M`$ for which $`\left|detM\right|=1`$, so that $`\mathrm{\Gamma }_M`$ is cocompact and discrete in $`G_M`$, hence quasi-isometric to $`G_M`$. In this case the proof of Theorem 11 outlined above breaks down, but this is so in part because the answer is quite different: the quasi-isometry classes of polycyclic $`\mathrm{\Gamma }_M`$ are much coarser than in the nonpolycyclic case, as the former are (conjecturally) determined by the absolute Jordan form up to real, as opposed to integral, powers. The key conjecture is:
###### Conjecture 15 (Horizontal preserving).
Suppose that $`\left|detM\right|,\left|detN\right|=1`$, and that $`M`$ and $`N`$ have no eigenvalues on the unit circle. Then every quasi-isometry of $`G_MG_N`$ is horizontal-respecting.
The general (with arbitrary eigenvalues) case of Conjecture 15, which is slightly more complicated to state, together with Theorem 13 easily implies:
###### Conjecture 16 (Classification of polycyclic, abelian-by-cyclic groups).
Suppose that $`\left|detM\right|,\left|detN\right|=1`$. Then $`\mathrm{\Gamma }_M`$ is quasi-isometric to $`\mathrm{\Gamma }_N`$ if and only if there exist nonzero $`a,b𝐑`$ so that $`M^a`$ and $`N^b`$ have the same absolute Jordan form.
Here by $`M^a`$ we mean $`\varphi (a)`$, where $`\varphi :𝐑\mathrm{GL}(n,𝐑)`$ is a $`1`$-parameter subgroup with $`\varphi (1)=M`$ (we are assuming that $`M`$ lies on such a subgroup, which can be assumed after squaring $`M`$).
Now let us concentrate on the simplest non-nilpotent example, which is also one of the central open problems in the field. The 3-dimensional geometry solv is the Lie group $`G_M`$ where $`M\mathrm{SL}(2,𝐙)`$ is any matrix with 2 distinct real eigenvalues (up to scaling, it doesn’t matter which such $`M`$ is chosen).
###### Conjecture 17 (Rigidity of solv).
The $`3`$-dimensional Lie group solv is quasi-isometrically rigid: any f.g. group $`G`$ quasi-isometric to solv is weakly commensurable with a cocompact, discrete subgroup of solv.
There is a natural boundary for solv which decomposes into two pieces $`^s\text{solv}`$ and $`^u\text{solv}`$; these are the leaf spaces of the weak stable and weak unstable foliations, respectively, of the vertical flow on solv, and are both homeomorphic to $`𝐑`$.
The isometry group $`\mathrm{Isom}(\text{solv})`$ acts on the pair $`(^s\text{solv},^u\text{solv})`$ affinely and induces a faithful representation $`\text{solv}\mathrm{Aff}(𝐑)\times \mathrm{Aff}(𝐑)`$ whose image consists of the pairs
$$(ax+b,a^1x+c),a𝐑^+,b,c𝐑$$
Just as quasi-isometries of hyperbolic space $`𝐇^n`$ are characterized by their quasiconformal action on $`𝐇^n`$ (a fact proved by Mostow), giving the formula $`\mathrm{QI}(H^n)=\mathrm{QC}(𝐇^n)`$, we conjecture:
###### Conjecture 18 (QI group of solv).
$$\mathrm{QI}(\text{solv})=(\mathrm{Bilip}(𝐑)\times \mathrm{Bilip}(𝐑))𝐙/2$$
where $`\mathrm{Bilip}(𝐑)`$ denotes the group of bilipschitz homeomorphisms of $`𝐑`$, and $`𝐙/2`$ acts by switching factors.
There is evidence for Conjecture 18: the direction $``$ is not hard to check (see \[FS\]), and the analogous theorem $`\mathrm{QI}(\mathrm{BS}(1,n))=\mathrm{Bilip}(𝐑)\times \mathrm{Bilip}(𝐐_n)`$ was proved in \[FM1\]. By using convergence groups techniques and a theorem of Hinkkanen on uniformly quasisymmetric groups (see \[FM2\]), we have been able to show:
Conjecture 15 (in the $`2\times 2`$ case) $``$ Conjecture 18 $``$ Conjecture 17
Here is a restatement of Conjecture 15 in the $`2\times 2`$ case:
###### Conjecture 19.
Every quasi-isometry $`f:\text{solv}\text{solv}`$ is horizontal respecting.
Here is one way *not* to prove Conjecture 19.
One of the major steps of \[FM1\] in studying $`\mathrm{BS}(1,n)`$ was to construct a model space $`X_n`$ for the group $`\mathrm{BS}(1,n)`$, study the collection of isometrically embedded hyperbolic planes in $`X_n`$, and prove that for any quasi-isometric embedding of the hyperbolic plane into $`X_n`$, the image has finite Hausdorff distance from some isometrically embedded hyperbolic plane.
However, solv has quasi-isometrically embedded hyperbolic planes which are *not* Hausdorff close to isometrically embedded ones. The natural left invariant metric on solv has the form
$$e^{2t}dx^2+e^{2t}dy^2+dt^2$$
from which it follows that the $`xt`$-planes and $`yt`$-planes are the isometrically embedded hyperbolic planes. But none of these planes is Hausdorff close to the set
$$\{(x,y,t)\text{solv}:x0\text{and}y=0\}\{(x,y,t)\text{solv}:y0\text{and}x=0\}$$
which is a quasi-isometrically embedded hyperbolic plane. An even stranger example is shown in Figure 1.
These strange quasi-isometric embeddings from $`𝐇^2`$ to solv do share an interesting property with the standard isometric embeddings, which may point the way to understanding quasi-isometric rigidity of solv. We say that a quasi-isometric embedding $`\varphi :𝐇^2\text{solv}`$ is *$`A`$-quasivertical* if for each $`x𝐇^2`$ there exists a vertical line $`\mathrm{}\text{solv}`$ such that $`\varphi (x)`$ is contained in the $`A`$-neighborhood of $`\mathrm{}`$, and $`\mathrm{}`$ is contained in the $`A`$-neighborhood of $`\varphi (𝐇^2)`$.
In order to study solv, it therefore becomes important to understand whether every quasi-isometrically embedded hyperbolic plane is quasi-vertical. Specifically:
###### Problem 20.
Show that for all $`K,C`$ there exists $`A`$ such that each $`K,C`$-quasi-isometrically embedded hyperbolic plane in solv is $`A`$-quasivertical.
Arguing by contradiction, if Problem 20 were impossible, fixing $`K,C`$ and taking a sequence of examples whose quasi-vertical constant $`A`$ goes to infinity, one can pass to a subsequence and take a renormalized limit to produce a quasi-isometric embedding $`𝐇^2\text{solv}`$ whose image is entirely contained in the upper half $`t0`$ of solv. But we conjecture that this is impossible:
###### Conjecture 21.
There does not exist a quasi-isometric embedding $`𝐇^2\text{solv}`$ whose image is entirely contained in the upper half space $`\text{solv}^+=\{(x,y,t)\text{solv}:t0\}`$.
## 6. Next steps
While we have already seen that there is a somewhat fine classification of finitely presented, nonpolycyclic abelian-by-cyclic groups up to quasi-isometry, this class of groups is but a very special class of finitely generated solvable groups. We have only exposed the tip of a huge iceberg. An important next layer is:
###### Problem 22 (Metabelian groups).
Classify the finitely presented (nonpolycyclic) metabelian groups up to quasi-isometry.
The first step in attacking this problem is to find a workable method of describing the geometry of the natural geometric model of such groups $`G`$. Such a model should fiber over $`𝐑^n`$, where $`n`$ is the rank of the maximal abelian quotient of $`G`$; inverse images under this projection of (translates of) the coordinate axes should be copies of the geometric models of abelian-by-cyclic groups.
Polycyclic versus nonpolycyclic. We’ve seen the difference, at least in the abelian-by-cyclic case, between polycyclic and nonpolycyclic groups. Geometrically these two classes can be distinguished by the trees on which they act: such trees are lines in the former case and infinite-ended in the latter. It is this branching behavior which should combine with coarse topology to make the nonpolycyclic groups more amenable to attack.
Note that a (virtually) polycyclic group is never quasi-isometric to a (virtually) nonpolycyclic solvable group. This follows from the theorem of Bieri that polycylic groups are precisely those solvable groups satisfying Poincare duality, together with the quasi-isometric invariance of the latter property (proved by Gersten \[Ge\] and Block-Weinberger \[BW\]).
Solvable groups as dynamical systems. The connection of nilpotent groups with dynamical systems was made evident in \[Gr1\], where Gromov’s Polynomial Growth Theorem was the final ingredient, combining with earlier work of Franks and Shub \[Sh\], in the positive solution of the Expanding Maps Conjecture: every locally distance expanding map on a closed manifold $`M`$ is topologically conjugate to an expanding algebraic endomorphism on an infranil manifold (see \[Gr1\]).
In §4 and §5 we saw in another way how invariants from dynamics give quasi-isometry invariants for abelian-by-cyclic groups. This should be no big surprise: after all, a finitely presented abelian-by-cyclic group is describable up to commensurability as an ascending HNN extension $`\mathrm{\Gamma }_M`$ over a finitely-generated abelian group $`𝐙^n`$. The matrix $`M`$ defines an endomorphism of the $`n`$-dimensional torus. The mapping torus of this endomorphism has fundamental group $`\mathrm{\Gamma }_M`$, and is the phase space of the suspension semiflow of the endomorphism, a semiflow with partially hyperbolic dynamics (when $`M`$ is an automorphism, and so $`\mathrm{\Gamma }_M`$ is polycyclic, the suspension semiflow is actually a flow). Here we see an example of how the geometric model of a solvable group is actually the phase space of a dynamical system.
But Bieri-Strebel \[BS1\] have shown that every finitely presented solvable group is, up to commensurability, an ascending HNN extension with base group a finitely generated solvable group. In this way every finitely presented solvable group is the phase space of a dynamical system, probably realizable geometrically as in the abelian-by-cyclic case.
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# Low temperature dynamics and laser-cooling of two-species Coulomb chains for quantum logic
## I Introduction
The rapid development of trapping techniques for neutral and charged particles has constituted a breakthrough in the investigation of quantum mechanical systems . Among the many interesting experiments, ordered structures of charged ions have been achieved in Paul and Penning traps . Such structures are composed from few up to thousands of particles, and they originate at low temperature from the combined effect of the Coulomb repulsive interaction among the ions and the trapping potential . Therefore, their geometry depends intrinsically on the trap set-up.
The field of interest of these so-called ”Coulomb Crystals” is rather broad, and in quantum optics they find an application in the ion-trap quantum computer . Here, a string of ions is proposed as a system for processing information, using two stable or metastable internal states of the ions for storing the quantum information and coherent interaction of the internal degrees of freedom of the single ions with the laser light for generating the unitary operations which process the information, while the coupling among ions is provided by the collective vibrational excitations of the chain. Present schemes for quantum information processing are based on the harmonicity of the ionic motion . This regime can be achieved by laser cooling the string of ions, and to a good extent the ions can be considered to vibrate harmonically around their equilibrium positions.
The coupling to the environment gives rise to phenomena which destroy the quantum coherence required for information processing: Decoherence affects both the internal and the motional quantum states. For motional states, decoherence can be inhibited by applying laser-cooling to the ion-chain on a regular rate or even continuously. Laser-cooling does not destroy the quantum information stored in the internal states, provided only some ions are addressed by the cooling laser (cooling ions), while the chain is sympathetically cooled via the Coulomb interaction. Furthermore it allows for simultaneous information processing on the other ions (qu-bits), provided the quantum gates do not require quantum coherence among the vibrational levels, as for the gate proposed in and realized in . In this respect, it is rather difficult to find a candidate ion which is, at the same time, a good qu-bit and cooling ion. In addition, some realizations of ion strings for quantum information do not rely on the spatial resolution of the ions with the laser . Hence, one of the most recent issues in the ion-trap quantum computer is to use two different ionic species which compose the ion chain, one for quantum information processing, the other for laser-cooling . This type of crystal, which we call here the two-species Coulomb crystal, has been already used for imaging the mechanical effect of radiation pressure on the crystallized ions and for sympathetic cooling of big ordered ionic structures .
In this work we study the mechanical motion in a two-species linear crystal from the point of view of quantum logic, illustrating the general features and the differences from a linear crystal composed of ions of equal masses, and we discuss in some detail sympathetic cooling of the chain. The issue of decoherence is also discussed in connection to the mechanical properties of the system. In particular, in this paper we consider the decoherence due to the coupling of the ionic motion to the fluctuations of the electric field (which we assume to have an instantaneous value with zero spatial gradient) along the crystal, and thus couples with the center of mass motion. Finally, we discuss the harmonic approximation and the effect of the anharmonic corrections on laser cooling and quantum information processing. An analogous analysis has been presented in , where they studied one specific scheme for sympathetic cooling. Here, we study the more general system and its properties, and seek schemes which are more suitable for realizing quantum logic. In particular, we provide some examples calculated with <sup>115</sup>In<sup>+</sup> and <sup>25</sup>Mg<sup>+</sup> ions, which have been currently trapped and cooled in Garching .
The paper is organized as follows. In Section II the small oscillations formalism is introduced, and the mechanical properties of a chain of ions are discussed in detail. In particular, sympathetic cooling of the two-species linear chain is studied, and the rates of cooling for different crystal configurations are derived. In Section III the effect of the anharmonic corrections on cooling and quantum logic is discussed. Summary and conclusions are presented in Section IV.
## II Small oscillations
In this section we investigate the collective vibrations around the equilibrium position of a chain of ions confined in a linear Paul trap, following the lines of the literature of small oscillations . We assume the motion to be one-dimensional, i.e. confined along the trap axis of the linear Paul trap. This corresponds to assuming a very steep radial potential .
We consider a one–dimensional string of $`N`$ ions with charge $`e`$ and mass equal either to $`M`$ or $`m`$, which are aligned along the $`\widehat{x}`$ axis, which corresponds to the axis of the linear Paul trap. Indicating with $`i`$ the position of the ion in the chain ($`i=1,\mathrm{},N`$), the sequence of ionic masses is described by the array $`𝐦=(m_1,\mathrm{},m_N)`$ with $`m_i=M,m`$. The ions are confined by the electrostatic potential $`V_S`$ and interact via the Coulomb repulsion . Sufficiently far away from the electrodes, $`V_S`$ can be considered harmonic and the total potential $`V`$ has the form:
$$V=\underset{i=1}{\overset{N}{}}\frac{1}{2}u_0x_i^2+\frac{1}{2}\underset{i=1}{\overset{N}{}}\underset{j=1,ji}{\overset{N}{}}\frac{e^2/4\pi ϵ_0}{|x_ix_j|}.$$
(1)
where $`x_i`$ is the coordinate of the ion $`i`$ and $`u_0`$ is a constant with the dimensions of an energy over a distance squared. If the ions are sufficiently cold , they crystallize around the classical equilibrium positions $`x_i^{(0)}`$, which are the solutions of the set of equations $`V/x_i|_{x_i^{(0)}}=0`$. Those solutions are independent of the mass, as the potential of the electrodes interacts only with the ionic charges. A characteristic quantity is the equilibrium distance between two ions:
$$x_0=x_2^{(0)}x_1^{(0)}=\left(\frac{2e^2/4\pi ϵ_0}{u_0}\right)^{1/3},$$
(2)
where the ions are displaced symmetrically with respect to the center of the trap. This quantity scales the inter-ionic distance for a $`N`$ ion chain .
Assuming that the collective motion of the ions around the equilibrium position is harmonic, we approximate $`V`$ with its Taylor expansion around $`x_i^{(0)}`$ truncated to second order. The dynamics of the system are described by the Lagrangian
$$L=\frac{1}{2}\left[\underset{i=1}{\overset{N}{}}m_i\dot{q}_i^2\underset{i=1}{\overset{N}{}}V_{ij}q_iq_j\right]$$
(3)
where $`q_i=x_ix_i^{(0)}`$ are the displacements of the ions from the equilibrium positions, and $`V_{ij}`$ are real coefficients which have the form:
$`V_{ij}`$ $`=`$ $`{\displaystyle \frac{^2}{x_ix_j}}V(x_1,\mathrm{},x_N)|_{\{x_i^{(0)}\}}`$ (4)
$`=`$ $`u_0+2{\displaystyle \underset{k=1,ki}{}}{\displaystyle \frac{e^2/4\pi ϵ_0}{|x_i^{(0)}x_k^{(0)}|^3}}\text{ if }i=j`$ (5)
$`=`$ $`2{\displaystyle \frac{e^2/4\pi ϵ_0}{|x_i^{(0)}x_j^{(0)}|^3}}\text{ if }ij.`$ (6)
From (3) the equations for the normal modes of the motion are
$$\underset{j=1}{\overset{N}{}}V_{ij}\beta _j^\alpha =\lambda _\alpha m_i\beta _i^\alpha \text{ with }\alpha =1,\mathrm{},N$$
(7)
where the eigenvalues $`\lambda _\alpha `$ are real given the hermiticity of $`V_{ij}`$, and where $`\beta ^\alpha `$ is eigenvector at $`\lambda _\alpha `$. Stable and harmonic oscillations exist if the condition $`\lambda _\alpha >0`$ is fulfilled for any $`\alpha `$, as it occurs in this case. Under this condition the frequency of the normal mode $`\mathrm{\Omega }_\alpha `$ is $`\mathrm{\Omega }_\alpha =\sqrt{\lambda _\alpha }`$. The eigenvectors $`\beta _i^\alpha `$ are orthogonal in the Riemannian metric with metric tensor $`𝐌`$, where $`𝐌`$ is a diagonal matrix whose diagonal corresponds to the array $`𝐦`$. For $`\lambda _\alpha >0`$, introducing the mass-weighted coordinates $`q_i^{}=\sqrt{m_i}q_i`$ the eigenvalue problem can be rewritten as
$$\underset{j}{}V_{ij}^{}\beta _{j}^{\alpha }{}_{}{}^{}=\mathrm{\Omega }_\alpha ^2\beta _{i}^{\alpha }{}_{}{}^{}\text{for}\alpha =1,\mathrm{},N,$$
(8)
where now $`V_{ij}^{}=V_{ij}/\sqrt{m_im_j}`$, and the metric tensor is the identity matrix, as for cartesian orthogonal coordinates. The eigenvalue problem is now equivalent to the one of $`N`$ identical ions of unitary mass. The matrix $`\{\beta _{i}^{\alpha }{}_{}{}^{}\}`$ defines an orthogonal transformation, which reduces the system to the principal axes $`\pi _\alpha `$ of $`V_{ij}`$ and of the kinetic term: $`\pi _\alpha =_i\beta _{i}^{\alpha }{}_{}{}^{}q_i^{}`$. In this representation the Lagrangian describes a set of $`N`$ independent harmonic oscillators with frequencies $`\mathrm{\Omega }_\alpha `$. We quantize the motion by associating a quantum mechanical oscillator with each mode. Then, denoting $`a_\alpha `$, $`a_\alpha ^{}`$ the annihilation, creation operators for the mode $`\alpha `$, respectively, the coordinate $`\pi _\alpha `$ associated with the oscillator of frequency $`\mathrm{\Omega }_\alpha `$ is written as $`\pi _\alpha =\sqrt{\mathrm{}/2\mathrm{\Omega }_\alpha }\left(a_\alpha +a_\alpha ^{}\right)`$. Going back to the original set of oblique coordinates $`q_i`$, they have the quantized form
$$q_i=\frac{1}{\sqrt{m_i}}\underset{\alpha }{}\left(\beta _i^\alpha \right)^1\sqrt{\frac{\mathrm{}}{2\mathrm{\Omega }_\alpha }}\left(a_\alpha +a_\alpha ^{}\right).$$
(9)
Some general features can now be recognized. From Eq. (8) it is evident that the eigenmodes depend on the values of the ionic masses. Furthermore, since the matrix $`V_{ij}`$ is symmetrical by exchange of any pair of ions, the properties of the motion will be mainly characterized by the symmetries of the sequence $`𝐦`$. These properties are reflected in the eigenmodes of the motion $`\{\beta _\alpha \}`$, and thus they affect the coupling of the crystal to radiation.
We discuss these points below. First, we consider the properties connected to two different values of the ionic masses, by analysing the case of two ions. Then, we discuss the ones connected to the symmetries of $`𝐦`$ by considering a three-ion crystal. Finally, on the basis of Eq. (9) we study the mechanical effect of radiation on the crystal, and in particular sympathetic cooling of the chain.
### A Two ions of different masses
We analyse here the two-ion crystal, where the ions have masses $`m`$ and $`M=\mu m`$ with $`\mu `$ real parameter, $`\mu >1`$. For this case the secular equation (8) can be solved analytically, and the eigenfrequencies of the motion have the form:
$$\mathrm{\Omega }_\pm ^2=\frac{u_0}{m}\left(1+\frac{1}{\mu }\pm \sqrt{1+\frac{1}{\mu ^2}\frac{1}{\mu }}\right)$$
(10)
with corresponding displacements:
$$q_\pm =N_\pm (\frac{1\mu \sqrt{1+\mu ^2\mu }}{\sqrt{\mu }},\frac{1}{\sqrt{\mu }})$$
(11)
where the first and second components refer to the particles of mass $`m`$ and $`M`$, respectively. Here, $`N_\pm `$ are the normalization factor, according to the scalar product $`q_i(1)q_j(1)+\mu q_i(2)q_j(2)=\delta _{ij}`$ with $`i,j=\pm `$.
In Fig. 1(a) and (b) we plot the eigenfrequencies and the eigenvectors $`q_\pm `$, respectively, as a function of $`\mu `$. The ratio $`\mu =1`$ corresponds to the the well-known case of two ions of equal masses in a linear trap, where the ratio of the eigenfrequencies are in the relation $`1:\sqrt{3}`$. For this value of $`\mu `$, $`\mathrm{\Omega }_{}`$ and $`\mathrm{\Omega }_+`$ correspond to the center of mass (COM) and stretch mode frequencies, respectively, as it can also be verified from Eq. (11). As $`\mu `$ increases, the value of the eigenfrequencies decreases and tends asymptotically to the values $`\mathrm{\Omega }_{}0`$ and $`\mathrm{\Omega }_+\sqrt{2u_0/m}`$. The limit $`\mathrm{\Omega }_{}0`$ corresponds to the case where both ions stand still at their equilibrium position, as it can be seen in Fig. 1(b), while for the limit $`\mathrm{\Omega }_+\sqrt{2u_0/m}`$ the heavy ion does not move, and the light ion oscillates around its equilibrium position.
In the following we will concentrate on the case $`\mu >1`$. As it is apparent from Fig. 1 (b) and Eq. (11) the two modes preserve some characteristics of the case of two ions with equal masses: In the mode of eigenfrequency $`\mathrm{\Omega }_{}`$ the ions oscillate in phase, whereas for $`\mathrm{\Omega }_+`$ they oscillate with opposite phases. The two modes, however, do not correspond to the COM and relative motion any longer. This can be understood by observing that, in absence of interactions, the trap frequency for an ion of mass $`m`$ is proportional to $`1/\sqrt{m}`$. This argument applies to an $`N`$-ion chain of two (or more) species, and can be verified by substituting the vector $`q^{\mathrm{COM}}=(1,1,\mathrm{},1)/\sqrt{N}`$ describing the center of mass motion inside the secular equation (7); one obtains $`_jV_{ij}=u_0=m_i\lambda ^{\mathrm{COM}}`$ for $`i=1,\mathrm{},N`$, which cannot be fulfilled for any value of $`\lambda ^{\mathrm{COM}}`$, unless all masses $`m_i`$ are equal.
The non-separability of the modes into center of mass and relative motion has some consequences on the dynamics. For example, for two ions the anharmonicity (i.e. the corrections to the harmonic approximation of the potential in Eq. (3)) couples the two modes, whereas in the crystals of ions with equal masses the COM motion is an exact eigenmode of the problem. A further consequence is the coupling of both modes to the fluctuations of the electric field at the trap-electrodes, since none of the modes is orthogonal to the COM motion. The strength with which each mode couples to this source of decoherence is a function of the mass-ratio between the species $`\mu `$, as has been discussed in .
### B $`N`$-ion crystal
As discussed above, the characteristic properties of the motion of a two-ion crystal are a function of the mass ratio $`\mu `$. For crystals with $`N>2`$ ions, some further parameters characterize the properties of the motion: The number of ions of each species and the sequence in which they are arranged. These two features are described through the array $`𝐦`$. In one dimension, the relevant symmetry property of the sequence $`𝐦`$ is the symmetry under reflections with respect to the center of the trap (which is also the center of the string). This corresponds to an invariance of the Hamiltonian under spatial parity transformations. Be $`\mathrm{\Pi }^{(N)}`$ the parity operator, defined on the wave functions $`|\varphi (x_1,x_2,\mathrm{},x_N)`$ of the N-ion Hilbert space as $`\mathrm{\Pi }^{(N)}|\varphi (x_1,x_2,\mathrm{},x_N)=|\varphi (x_1,x_2,\mathrm{},x_N)`$. This operator has eigenvalues $`p=+1`$ (even), $`p=1`$ (odd), corresponding to the states with even and odd parity, respectively. If the array $`𝐦`$ is symmetric under reflections, the Hamiltonian for the small oscillations commute with $`\mathrm{\Pi }^{(N)}`$ and the eigenmodes of Eq. (8) are also eigenvectors of $`\mathrm{\Pi }^{(N)}`$ at the eigenvalue either $`p=1`$ or $`p=1`$. From a simple evaluation of the number of degrees of freedom, one can verify that the even modes are $`N/2`$ for an even number $`N`$ of ions and $`(N1)/2`$ for $`N`$ odd. In particular, for symmetry reasons the central ion does not move in the modes of even parity of a chain with an odd number of ions. Thus, these modes are independent of the mass of the central ion, as it can be deduced from Eq. (7).
It is instructive to take a closer look at the eigenfrequencies of a chain of $`N`$ ions as a function of all possible sequences $`𝐦`$. We discuss the case of 3 ions, since it shares some similarities with the normal modes of a triatomic molecule, as discussed in textbooks , and it exhibits features which can be extended to chains of large $`N`$.
In Fig. 2 we plot the eigenfrequencies of a $`N=3`$ chain as a function of all possible sequences of Indium and Magnesium ions, where the sequences have been ordered with increasing total mass $`M_C`$ of the crystal. The mode of frequency $`\mathrm{\Omega }_1`$ (solid line) is characterized by the oscillation in phase of the three ions. For ions of equal mass it corresponds to the center of mass mode, but in all cases still represents the mechanical response of the whole crystal to excitations: In fact, $`\mathrm{\Omega }_11/\sqrt{M_C}`$ and in general it does not show an appreciable dependence on the order in which the ions are arranged.
The properties of the higher excitations depend on $`\mu `$ and on the sequence. In particular, the distance among the eigenenergies changes depending on where the heavy ions are placed in the sequence. In addition, symmetric sequences preserve the properties of a chain of three ions of equal masses. Thus, the eigenmode of frequency $`\mathrm{\Omega }_2`$ (dashed line) is characterized by the out-of-phase oscillation of the external ions, whereas the central ion stands still. Hence, $`\mathrm{\Omega }_2`$ takes the same value for the sequences A and C, and for the sequences D and F. In asymmetric sequences the external ions still oscillate out-of-phase, but the oscillation amplitude of the central ion is large, independently of its mass, provided that $`\mu 1`$.
In the symmetric sequences of the eigenmode $`\mathrm{\Omega }_3`$ the external ions oscillate in phase, whereas the central ion is out-of-phase, and its amplitude is a monotonic function of $`1/\mu `$. For the sequence C, in particular, $`\mathrm{\Omega }_2<\mathrm{\Omega }_3`$. The two frequencies are almost degenerate since both modes correspond to the case where the outer ions move symmetrically with respect to the center, since the central ion in $`\mathrm{\Omega }_3`$ has small displacements. In asymmetric sequences (B,E) the light ions have large oscillation amplitudes, whereas the displacements of the heavy ions are smaller.
These properties have some immediate implications for quantum logic with a two-component chain. For example, the eigenmodes of even parity of symmetric sequences are decoupled from the fluctuations of the electric field at the electrodes, and thus are good candidate for the quantum bus. This issue have been discussed quantitatively for a particular sequence in .
Furthermore, in symmetric sequences the parity operator $`\mathrm{\Pi }^{(N)}`$ commutes with the one-dimensional Hamiltonian $`H`$: Hence, the eigenmodes with odd parity are not coupled via anharmonicity to the eigenmodes with even parity. The axial motion, however, is coupled to the radial motion by non-linearities, and in three-dimension there are no groups of modes which are decoupled from the others. We discuss this point in section III.
Another important implication regards the spacing between the eigenfrequencies. In fact, when choosing a mode for quantum logic, the distance in energy among the eigenfrequencies should be taken into account, since the presence of quasi-degeneracies will lower the efficiency of single mode addressing in quantum logic operations. Thus, it is not convenient to use the mode with eigenfrequency $`\mathrm{\Omega }_3`$ in the sequence C as quantum bus, given its closeness to $`\mathrm{\Omega }_2`$.
Finally, from the mechanical properties of the chain we can infer in which position the cooling ions should be placed for optimal sympathetic cooling of the chain. We analyse this aspect in the next subsection.
### C Interaction of the crystal with light
In the dipole approximation the coupling of the external degrees of freedom of an atom with radiation is represented by the kick operator $`\mathrm{exp}(i𝐤𝐫)`$, where $`𝐤`$ is the wave vector of light and $`𝐫`$ the atomic position. In a Coulomb crystal of atomic ions, optical light couples to the internal degrees of freedom of a single ion, which we consider here a two-level dipole transition, and to the external degrees of freedom of the collective motion. Thus, assuming that the ion $`j`$ scatters a photon and that $`|\psi _i`$ is the motional state of the crystal before the scattering, after the scattering the motion of the crystal is described by the state $`|\psi _f`$ given by:
$`|\psi _f`$ $`=`$ $`\mathrm{exp}(ikx_j)|\psi _i`$ (12)
$`=`$ $`\text{e}^{i\varphi }\mathrm{exp}(ikq_j)|\psi _i`$ (13)
where $`\varphi =kx_j^{(0)}`$ is a real scalar, and $`k`$ is the projection of $`𝐤`$ on the axis of the crystal. The coupling of radiation to the collective modes of the crystal is visible by substituting Eq. (9) into (13). Thus, the kick operator can be written as
$$\mathrm{exp}(ikq_j)=\mathrm{\Pi }_{\alpha =1}^N\text{e}^{i\eta _j^\alpha (a_\alpha ^++a_\alpha )}$$
(14)
where $`\eta _j^\alpha `$ is the Lamb-Dicke parameter for the mode $`\alpha `$ and the ion $`j`$, and is defined as:
$$\eta _j^\alpha =k\beta _j^\alpha \sqrt{\frac{\mathrm{}}{2m_j\mathrm{\Omega }_\alpha }}$$
(15)
Through the Lamb-Dicke parameter we can infer the mechanical response of the crystal to the scattering of a photon. For one ion of mass $`m`$, it corresponds to the square root of the ratio between the recoil frequency $`\omega _R=\mathrm{}k^2/2m`$ and the trap frequency $`\mathrm{\Omega }`$: $`\eta =\sqrt{\omega _R/\mathrm{\Omega }}`$, and it determines how many motional levels are coupled by the scattering of one photon. The Lamb-Dicke regime corresponds to the case in which $`\omega _R\mathrm{\Omega }`$, and mathematically to the condition $`\eta \sqrt{n}1`$, where $`n`$ is the vibrational state number. In this limit the kick operator can be expanded in powers of $`\eta `$, and a change in the quantum motional state due to the incoherent scattering of one photon is of higher order in $`\eta `$ . In this regime an ion can be sideband cooled to the ground state of the vibrational motion .
In presence of more than one ion, the Lamb-Dicke parameter $`\eta _j^\alpha `$ describes how the displacement of the ion $`j`$ couples to the mode $`\alpha `$. In particular, it determines (i) the possibility of addressing a single motional sideband, which appears when scanning a probe beam through the resonance frequency (and thus of exciting one mode selectively) , and (ii) the possibility of laser-cooling a mode to its vibrational ground state, in analogy to the one-ion case.
Let us first consider the response to light of a two-ion crystal, and compare the case where the ions have equal masses, as discussed in , with the case where they have different masses. In the first case, due to the symmetry of the configuration the Lamb-Dicke parameters for each ion are equal (apart for some difference in the sign). In the second case, we can see in Eq. (15) that the Lamb-Dicke parameters depend on the mass of the ion, and given the asymmetry of the crystal the geometrical factor $`\beta _j^\alpha `$ has different values for the two ions. One might be tempted to think that the Lamb-Dicke regime can be more easily accessed by addressing the heavier ions. This, however, is actually not true for all modes. This consideration is particularly applicable to the mode of lowest frequency, which shares some properties with the center of mass motion, and in general describes the response of the crystal as a whole to the mechanical excitation. For this mode, the displacements of the two ions are comparable, and actually the displacement of the heavier ion is slightly larger, as can be seen in Fig. 1(b). On the contrary, the eigenmode with frequency $`\mathrm{\Omega }_+`$ is characterized by smaller displacements of the heavy ion than of the lighter one. For this mode, the heavy ion might be well in the Lamb-Dicke regime, while the light ion is not.
This situation can be visualized by comparing the absorption spectra obtained by shining light on each ion separately. We define the absorption spectrum in a two-level transition with resonant frequency $`\omega _0`$ and driven by a laser of frequency $`\omega _L`$ through the function $`I(\delta )`$, where $`\delta =\omega _L\omega _0`$ is the detuning, which is evaluated by summing all contributions to laser–excited transitions at frequency $`\omega _L`$ :
$$I(\delta )=\underset{E_𝐧E_𝐥=\delta }{}|𝐧|\mathrm{exp}(ikq_j)|𝐥|^2P(𝐧)$$
(16)
Here $`|𝐧=|(n_1,n_2)`$ are the motional states of energy $`E_𝐧=n_1\mathrm{}\mathrm{\Omega }_1+n_2\mathrm{}\mathrm{\Omega }_2`$, with $`n_\alpha `$ occupation of the mode $`\mathrm{\Omega }_\alpha `$ ($`\alpha =1,2`$), and $`P(𝐧)`$ is a normalized distribution over the motional states $`|𝐧`$. The coordinate $`q_j`$ characterizes the driven ion. In Fig. 3 we plot $`I(\delta )`$ as a function of $`\delta `$ for a crystal composed of a Magnesium and an Indium ion. Figure 3 shows the absorption spectrum when driving separately the Magnesium and the Indium ion, for a thermal distribution over the motional states with average total energy $`E=5\mathrm{}\mathrm{\Omega }_{}`$, where for the chosen parameters $`\mathrm{\Omega }_{}=0.552`$ MHz and $`\mathrm{\Omega }_+=1.456`$ MHz. In both cases the motional sidebands of the mode with frequency $`\mathrm{\Omega }_{}`$ are visible, whereas when driving the In<sup>+</sup> ion the sidebands of $`\mathrm{\Omega }_+`$ almost disappear. In fact, as evaluated from (15) the Lamb-Dicke parameters for $`\mathrm{\Omega }_{}`$ are $`\eta _{\mathrm{Mg}}^{}=0.22`$, $`\eta _{\mathrm{In}}^{}=0.38`$, whereas the ones for $`\mathrm{\Omega }_+`$ are $`\eta _{\mathrm{Mg}}^+=0.5`$, $`\eta _{\mathrm{In}}^+=0.06`$, i.e. the weight of the motional sidebands for the mode $`\mathrm{\Omega }_+`$ in Indium is two orders of magnitude smaller than the ones for the mode $`\mathrm{\Omega }_{}`$. It is interesting to compare this result with the case of two ions of equal masses. In that case the absorption spectrum is the same independently of which ion of the chain is driven. Then, if the COM mode is in the Lamb-Dicke regime, the relative motion mode is also in the Lamb-Dicke regime, since in that case the Lamb-Dicke parameter scales simply as the inverse of the squared root of the eigenfrequency.
For crystals with $`N>2`$ ions the factor $`\beta _j^\alpha `$ in (15) reflects the structure of the chain and consequentely how the driven ion couples to the mode to cool. It thus contains some information on where the cooling ions should be placed in the sequence so as to achieve more efficient cooling. This can be theoretically illustrated by a rate equation, which describes cooling of one mode in a chain in the Lamb-Dicke regime . Here, we define the Lamb-Dicke regime with the condition:
$$\text{max}_{\{j\}}|\eta _j^\alpha |^2n_\alpha 1\text{ for }\alpha =1,\mathrm{},N$$
(17)
where $`n_\alpha `$ denotes the occupation of the mode $`\alpha `$ and where the set of ions $`\{j\}`$ represents the set of positions of the cooling ions in the chain. Assuming that the laser interacts with the internal two-level transition of the individual atom, in second order perturbation theory in the parameter $`g/\gamma `$, with $`g`$ Rabi frequency and $`\gamma `$ decay rate of the atomic transition, the excited state can be eliminated, and one obtains a set of equations projected on the electronic ground state, where populations and coherences between different motional states are coupled. In the Lamb-Dicke limit such coupling can be neglected, thus reducing the equation for the $`N`$-ion density matrix $`\rho `$ in the low saturation limit to rate equations . Furthermore, introducing the reduced density matrix $`\rho _\alpha `$ for the mode $`\alpha `$ defined by:
$$\rho _\alpha =\underset{n_{\beta _1}}{}\mathrm{}\underset{n_{\beta _{n1}}}{}n_{\beta _1},\mathrm{},n_{\beta _{n1}}|\rho |n_{\beta _1},\mathrm{},n_{\beta _{n1}},$$
(18)
with $`\beta _1,\mathrm{},\beta _{n1}\alpha `$, one can derive the rate equation for cooling of the mode $`\alpha `$ in one-dimension :
$`{\displaystyle \frac{\text{d}}{\text{d}t}}P(n_\alpha )`$ $`=`$ $`{\displaystyle \frac{g^2}{\gamma }}{\displaystyle \underset{\{j\}}{}}\eta _{j}^{\alpha }{}_{}{}^{2}[(n_\alpha +1)A_{}^\alpha P(n_\alpha +1)`$ (19)
$``$ $`((n_\alpha +1)A_+^\alpha +n_\alpha A_{}^\alpha )P(n_\alpha )+n_\alpha A_+^\alpha P(n_\alpha 1)`$ (20)
the coupling with the other modes being of higher order in the Lamb-Dicke parameter. Here, $`P(n_\alpha )=n_\alpha |\rho _\alpha |n_\alpha `$ and the coefficients $`A_\pm ^\alpha `$ are defined as
$`A_+^\alpha ={\displaystyle \frac{1}{16\mathrm{\Omega }_\alpha ^2/\gamma ^2+1}}+{\displaystyle \frac{2}{5}}{\displaystyle \frac{1}{4\mathrm{\Omega }_\alpha ^2/\gamma ^2+1}},`$ (21)
$`A_{}^\alpha =1+{\displaystyle \frac{2}{5}}{\displaystyle \frac{1}{4\mathrm{\Omega }_\alpha ^2/\gamma ^2+1}}.`$ (22)
where we have assumed that the Rabi coupling is spatially constant over the whole crystal and equal to $`g`$, and that the laser is tuned on the first red sideband of the mode $`\mathrm{\Omega }_\alpha `$. Eq. (19) is the sum of all contributions to cooling of mode $`\alpha `$ from the coupling $`g`$ of the cooling laser to the driven ions. It is fully equivalent to the equation for cooling of one ion in a harmonic trap of frequency $`\mathrm{\Omega }_\alpha `$ in the Lamb–Dicke regime , apart from the scaling factor multiplying the term on the RHS of (19):
$$W_\alpha =\underset{\{j\}}{}|\eta _j^\alpha |^2$$
(23)
Provided that (23) is different from zero, it does not affect the steady state, but scales the rate of cooling of the mode $`\alpha `$. The factor $`W_\alpha `$ represents the contribution of the cooling ions in the array to the speed of the process. Obviously, the largest cooling rate is achieved when all ions are driven by the cooling laser. In that case, $`W_\alpha `$ has the form $`W_\alpha ^{\mathrm{max}}=_{j=1}^N|\eta _j^\alpha |^2=\omega _R/\mathrm{\Omega }_\alpha `$, where $`\omega _R`$ is the recoil frequency of the single ion. The rate of cooling of each mode scales according to the relation $`W_\alpha ^{\mathrm{max}}=W_\alpha ^{}^{\mathrm{max}}\mathrm{\Omega }_\alpha ^{}/\mathrm{\Omega }_\alpha `$, and it scales with the mass $`m`$ of the cooling ion as $`W_\alpha ^{\mathrm{max}}\sqrt{1/m}`$ since $`\omega _R1/m`$ and $`\mathrm{\Omega }_\alpha 1/\sqrt{m}`$.
For quantum logic we are interested in employing only some ions of the chain for cooling. We look thus for the best sequence and mode for cooling, where for “best sequence” we intend a compromise between the highest number of ions for quantum logic (i.e. the lowest number of cooling ions) and the best cooling rate.
In Fig. 4 we plot the factor $`W_\alpha `$ vs all possible sequences of $`N=3`$ ions, made up of Magnesium and/or Indium ions, where the cooling ions are in (a) Indium and (b) Magnesium ions. Note that between the curves in (a) and (b) there is a scaling factor corresponding, as expected, to the squared root of the ratio of the ionic masses. The curves in Fig. 4 can be easily interpreted by considering the properties of the modes, as discussed in the previous subsection. Thus, as the mode of frequency $`\mathrm{\Omega }_1`$ is characterized by an oscillation in phase of all ions, and does not strongly depend on the sequence, the cooling rate increases as the number of cooling ions increases. On the other hand, the mode with frequency $`\mathrm{\Omega }_2`$ is mainly characterized by the oscillation of opposite phase of the ions placed externally. Thus, large rates of cooling are achieved when the cooling ions are placed in the external positions. In particular, the cooling rates of the sequences A and C, and of the sequences D and F are equal. In fact, these configurations are symmetric under reflection and the mode with eigenfrequency $`\mathrm{\Omega }_2`$ has even parity. Thus, the central ion does not contribute to cooling nor to quantum information processing with the mode. Finally, in symmetric configurations, the mode of frequency $`\mathrm{\Omega }_3`$ is characterized by large oscillations of the central ions, and a relatively large rate of cooling is obtained by simply placing the cooling ion in the center, as can be noticed for sequence C in Fig. 4(a) and sequence D in Fig. 4(b). In asymmetric sequences (B,E) the rate of cooling in (a) is small, whereas in (b) is large, as expected from the considerations made in the previous subsection.
For the case here considered, and on the basis of the mechanical properties only, the best sequence of cooling is E provided that the cooling ion has lighter mass (Fig. 4(b)).
In general, we can conclude that by preparing certain sequences one can have efficient sympathetic cooling of some modes using a relatively small number of cooling ions. This characteristic does not depend strongly on the mass-ratio $`\mu `$ between the qubit and the cooling ions. In a crystal with large total mass $`M_C`$, however, the Lamb-Dicke regime condition can be accessed more easily. In this respect, it would be better to use heavier ions for quantum logic.
Finally, sequences with an even number $`N`$ of ions are to be preferred over sequences with odd $`N`$, so that all positions in the chain contribute either to quantum logic or to the cooling process.
## III Effects of the anharmonicity on cooling and quantum logic
The harmonic approximation of the mechanical potential in (1) relies on the assumption that $`V`$ possesses strict local minima, around which the motion is well-localized. In that case, the higher orders of its Taylor expansion are a small correction. For two ions those terms have the form (for $`x_2>x_1`$):
$$\delta V=\underset{n=3}{\overset{\mathrm{}}{}}\delta V^{(n)}=\underset{n=3}{\overset{\mathrm{}}{}}(1)^n\frac{e^2/4\pi ϵ_0}{x_0^{n+1}}\left[q_2q_1\right]^n$$
(24)
where $`\delta V^{(n)}`$ is the $`n`$-th order correction. The effect of these terms, the so-called anharmonicity, consists in causing shifts to the motional energies, and mixing the eigenstates of the normal modes. Such mixing is in general a small correction to the eigenstates of the ideal case, but it may become particularly enhanced because of quasi degeneracies among the motional energies of the states. In fact, the density of motional states of an $`N`$-ion chain in the interval of energy $`[E,E+\delta E]`$ is approximately $`D(E)E^{N1}`$. Thus, as the number $`N`$ of ions increases and/or for larger values of the motional energies, the dynamics of quasi-degenerate states are definetely affected by the anharmonicity. Here, we discuss the effects of the departures from the ideal harmonic system, first on sideband cooling and then on the efficiency of quantum logic gates.
In sideband cooling the laser addresses the motional sideband of the mode to be cooled. Thus, cooling will be efficient as long as the shift in energy caused by the anharmonicity is much smaller than the frequency of one phonon of the mode addressed. On the other hand, the mixing between the eigenstates will constitute a thermalization effect among the modes, and it will not constitute an obstacle to cooling as long as all modes are at sufficiently low temperatures, so that $`\delta V`$ is a small correction to the whole system.
To obtain some estimates, we evaluate the order of magnitude of the shift to the energy in first order perturbation theory, and ask for which range of values of the vibrational numbers sideband cooling may still work. In it has been shown that in the perturbative regime $`\delta V\delta V^{(3)}`$ . Thus:
$$\delta V^{(3)}\frac{e^2/4\pi ϵ_0}{x_0}\left(\frac{a_{0\alpha }}{x_0}\right)^3n_\alpha ^{3/2}$$
(25)
where $`a_{0\alpha }=\sqrt{\mathrm{}/m_i\mathrm{\Omega }_\alpha }`$, $`n_\alpha `$ vibrational number of the mode $`\alpha `$, and $`m_i`$ mass of the lighter ion. In deriving (25) we have assumed $`n_\alpha /\mathrm{\Omega }_\alpha n_\beta /\mathrm{\Omega }_\beta `$ ($`\alpha ,\beta =1,2,\mathrm{},N`$). Taking two ions, one Indium and one Magnesium, $`\mathrm{\Omega }_1=1`$MHz and the mass of Magnesium <sup>25</sup>Mg<sup>+</sup>, then $`|\delta V|/\mathrm{}6\times 10^3n_1^{3/2}\mathrm{\Omega }_1`$, which implies that first order perturbation theory holds for $`n_180`$ ($`n_230`$). In this limit, a laser tuned on the first sideband to the red of the mode $`\mathrm{\Omega }_1`$ cools the system to the ground state. We have verified this conjecture numerically: We have taken a two-ion chain composed of an Indium and a Magnesium ion, and considered sideband cooling in the Lamb-Dicke regime of one of the modes, comparing the case in which the mechanical potential is fully harmonic with the case where the third and fourth orders in the anharmonic expansion have been included. We have not noticed any significant difference between the two cases, and the system is cooled efficiently to the ground state. In particular, no visible effect could be interpreted as due to the anharmonic coupling between quasi-degenerate states. However, according to the above estimates, in the interval of states of the numerical calculation the spectrum of energy is not very “dense”. Then, in order to verify numerically the effect of anharmonocities in presence of quasi-degeneracies we take a system with exact degeneracies, and more specifically with two modes of frequencies $`\mathrm{\Omega }_1=\mathrm{\Omega }`$ and $`\mathrm{\Omega }_2=2\mathrm{\Omega }`$. Here, in the Lamb-Dicke regime a laser sideband-cools the mode of frequency $`\mathrm{\Omega }_1`$. We compare the harmonic with the anharmonic case, where here we simply substitute the chosen values $`\mathrm{\Omega }_1,\mathrm{\Omega }_2`$ in the quantized form of the displacements of Eq. (9). In Fig. 5 the average occupation number for the modes (a) $`\mathrm{\Omega }_1`$ and (b) $`\mathrm{\Omega }_2`$ is plotted as a function of the time. The dashed and solid lines correspond to the harmonic and anharmonic case, respectively. In the harmonic case the mode $`\mathrm{\Omega }_1`$ is cooled and $`\mathrm{\Omega }_2`$ is like “frozen”, since it is coupled to radiation at higher orders in the Lamb-Dicke parameter . In the anharmonic case the rate of cooling of the mode addressed is slowed down, whereas the mode $`\mathrm{\Omega }_2`$ is simultaneously cooled: The system is cooled as a whole, but on a relatively slower time-scale. Thus, the two modes thermalize on a time-scale which is faster than the cooling one.
With respect to quantum information processing, present quantum logic schemes with ions are based on the harmonic properties of the motion. Thus, they are affected both by the shift in energy and the mixing induced by the anharmonicity. Quantum gates which require the preparation of the system in the ground state of the motion will be very weakly affected, since in that part of the spectrum the anharmonic perturbation is extremely small, and there are no quasi-degeneracies. On the other hand, the perturbation may affect the efficiency of hot gates, since they operate on higher-lying motional states. Thus, the speed of hot gates must be faster than the rate of anharmonic coupling between quasi-degenerate states. We can define a time-scale $`\tau _{\mathrm{Anh}}`$ for the anharmonic perturbation $`\tau _{\mathrm{Anh}}\mathrm{}/|\delta V|`$. The typical duration of a quantum gate must be shorter than $`\tau _{\mathrm{Anh}}`$. In presence of degenerate states which are coupled by three-phonon transitions (i.e. for $`\mathrm{\Omega }_\alpha 2\mathrm{\Omega }_\beta `$), from (25) $`\tau _{\mathrm{Anh}}10\mu `$sec given $`\mathrm{\Omega }_\alpha =1`$ MHz and $`n_\alpha 10`$. This estimate, which is rather worrying if compared with the duration of a quantum gate , reflects the worst case, which might occur for certain sequences, large number of ions and large excitations. Note that, if the degenerate states are coupled by a four-phonon transition (i.e. for $`\mathrm{\Omega }_\alpha 3\mathrm{\Omega }_\beta `$), then $`\tau _{\mathrm{Anh}}600\mu `$sec. In general, we expect this problem to arise when working with a large number of ions and for high excitations. It could be minimized by choosing symmetric sequences: in that case the “effective density” of motional states which are coupled by the anharmonicities of the axial potential will decrease, since only states with the same parity will be coupled to each other, whereas the coupling with the radial degrees of freedom is of higher order.
As a general rule, however, one should avoid degeneracies among the radial and the axial frequencies.
## IV Conclusions
We have studied the small oscillations behaviour of a two-component linear crystal, with particular emphasis on the applications to sympathetic cooling and quantum logic with the ions, and have discussed the effect of anharmonicity on the operations of the ion-trap quantum computer. We have seen that higher efficiency in quantum logic and sympathetic cooling are achieved by selecting the right ionic sequences. That raises the issue of how to prepare the desired sequence of ions. A rigorous investigation in this direction should take into account the full non-linear potential in the three-dimensional space and it is subject of on-going investigations.
We would like to thank P. Lambropoulos for many stimulating discussions and the critical reading of this manuscript, and W. Lange, S. Köhler, V. Ludsteck, E. Peik, who are involved in the experimental realization of the ion structure discussed here. This work is supported in parts by the European Commission within the TMR-networks ERB-FMRX-CT96-0087 and ERB-FMRX-CT96-0077.
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# Compatible Poisson brackets of hydrodynamic type
## 1 Introduction
In 1983 Dubrovin and Novikov introduced the Poisson brackets of hydrodynamic type
$$\{F,G\}=\frac{\delta F}{\delta u^i}A^{ij}\frac{\delta G}{\delta u^j}𝑑x$$
(1)
defined by the Hamiltonian operators $`A^{ij}`$ of the form
$$A^{ij}=g^{ij}\frac{d}{dx}+b_k^{ij}u_x^k,b_k^{ij}=g^{is}\mathrm{\Gamma }_{sk}^j.$$
(2)
They proved that in the nondegenerate case $`(detg^{ij}0)`$ the bracket (1), (2) is skew-symmetric and satisfies the Jacobi identities if and only if the metric $`g^{ij}`$ (with upper indices) is flat, and $`\mathrm{\Gamma }_{sk}^j`$ are the Christoffel symbols of the corresponding Levi-Civita connection.
Let us assume that there is a second Poisson bracket of hydrodynamic type defined on the same phase space by the Hamiltonian operator
$$\stackrel{~}{A}^{ij}=\stackrel{~}{g}^{ij}\frac{d}{dx}+\stackrel{~}{b}_k^{ij}u_x^k,\stackrel{~}{b}_k^{ij}=\stackrel{~}{g}^{is}\stackrel{~}{\mathrm{\Gamma }}_{sk}^j,$$
(3)
corresponding to a flat metric $`\stackrel{~}{g}^{ij}`$. Two Poisson brackets (Hamiltonian operators) are called compatible, if their linear combinations $`\stackrel{~}{A}^{ij}+\lambda A^{ij}`$ are Hamiltonian as well. This requirement implies, in particular, that the metric $`\stackrel{~}{g}^{ij}+\lambda g^{ij}`$ must be flat for any $`\lambda `$ (plus certain additional restrictions). The necessary and sufficient conditions of the compatibility were first formulated by Dubrovin , (see , for further discussion). In sect. 2 we reformulate these conditions in terms of the operator $`r_j^i=\stackrel{~}{g}^{is}g_{sj}`$ (Theorem 1) which, in particular, imply the vanishing of the Nijenhuis tensor of the operator $`r_j^i`$:
$$N_{jk}^i=r_j^s_sr_k^ir_k^s_sr_j^ir_s^i(_jr_k^s_kr_j^s)=0$$
(see , ).
Examples of compatible Hamiltonian pairs naturally arise in the theory of Hamiltonian systems of hydrodynamic type — see e.g. , , , , , . Dubrovin developed a deep theory for a particular class of compatible Poisson brackets arising in the framework of the associativity equations , . Compatible Poisson brackets of hydrodynamic type can also be obtained as a result of Whitham averaging (dispersionless limit) from the local compatible Poisson brackets of integrable systems , , , , , , . Some further examples and partial classification results can be found in , , , , , .
If the spectrum of $`r_j^i`$ is simple, the vanishing of the Nijenhuis tensor implies the existence of a coordinate system where both metrics $`g^{ij}`$ and $`\stackrel{~}{g}^{ij}`$ become diagonal. In these diagonal coordinates the compatibility conditions take the form of an integrable reduction of the Lamé equations. We present the corresponding Lax pairs in sect. 3. Another approach to the integrability of this system has been proposed recently by Mokhov by an appropriate modification of Zakharov’s scheme .
The main observation of this paper is the relationship between compatible Poisson brackets of hydrodynamic type and hypersurfaces $`M^{n1}E^n`$ which possess nontrivial deformations preserving the Weingarten operator. For surfaces $`M^2E^3`$ these deformations have been investigated by Finikov and Gambier as far back as in 1933 , . In sect. 4 we demonstrate that the n-orthogonal coordinate system in $`E^n`$ corresponding to the flat metric $`\stackrel{~}{g}^{ij}+\lambda g^{ij}`$ (rewritten in the diagonal coordinates) deforms with respect to $`\lambda `$ in such a way that the Weingarten operators of the coordinate hypersurfaces are preserved up to constant scaling factors. In sect. 5 we discuss surfaces $`M^2E^3`$ which possess nontrivial one-parameter deformations preserving the Weingarten operator and explicitly introduce a spectral parameter in the corresponding Gauss-Codazzi equations.
## 2 Differential-geometric criterion of the compatibility
To formulate the necessary and sufficient conditions of the compatibility we introduce the operator $`r_j^i=\stackrel{~}{g}^{is}g_{sj}`$, which is automatically symmetric
$$r_s^ig^{sj}=r_s^jg^{si},$$
(4)
so that $`\stackrel{~}{g}^{ij}=r_s^ig^{sj}=r_s^jg^{si}=r^{ij}`$. In what follows we use the first metric $`g^{ij}`$ for raising and lowering the indices.
###### Theorem 1
Hamiltonian operators (2), (3) are compatible if and only if the following conditions are satisfied:
1. The Nijenhuis tensor of $`r_j^i`$ vanishes:
$$N_{jk}^i=r_j^s_sr_k^ir_k^s_sr_j^ir_s^i(_jr_k^s_kr_j^s)=0.$$
(5)
2. The metric coefficients $`\stackrel{~}{g}^{ij}=r^{ij}`$ satisfy the equations
$$^i^jr^{kl}+^k^lr^{ij}=^i^kr^{jl}+^j^lr^{ik}.$$
(6)
Here $`^i=g^{is}_s`$ is the covariant differentiation corresponding to the metric $`g^{ij}`$. The vanishing of the Nijenhuis tensor implies the following expression for the coefficients $`\stackrel{~}{b}_k^{ij}`$ in terms of $`r_j^i`$:
$$2\stackrel{~}{b}_k^{ij}=^ir_k^j^jr_k^i+_kr^{ij}+2b_k^{sj}r_s^i$$
(7)
In a somewhat different form the necessary and sufficient conditions of the compatibility were formulated in , , , .
Remark. The criterion of the compatibility of the Hamiltonian operators of hydrodynamic type resembles that of the finite-dimensional Poisson bivectors: two skew-symmetric Poisson bivectors $`\omega ^{ij}`$ and $`\stackrel{~}{\omega }^{ij}`$ are compatible if and only if the Nijenhuis tensor of the corresponding recursion operator $`r_j^i=\stackrel{~}{\omega }^{is}\omega _{sj}`$ vanishes. We emphasize that in our situation operator $`r_j^i`$ does not coincide with the recursion operator.
Proof of Theorem 1:
Recall that in terms of $`g^{ij}`$ and $`b_k^{ij}`$ the conditions for the operator $`A`$ to be Hamiltonian take the form
$$2b_s^{ki}g^{sj}=g^{js}_sg^{ik}+g^{ks}_sg^{ij}g^{is}_sg^{kj}$$
(8)
and
$$g^{js}_sb_n^{ik}g^{is}_sb_n^{jk}+(b_s^{ij}b_s^{ji})b_n^{sk}+b_s^{ik}b_n^{js}b_s^{jk}b_n^{is}=0,$$
(9)
respectively (the last condition follows from the identity $`g^{il}g^{js}R_{nls}^k=0`$ after rewriting it in terms of $`b_k^{ij}`$). Note that (8) is equivalent to a pair of simpler conditions
$$b_k^{ij}+b_k^{ji}=_kg^{ij},b_s^{ik}g^{sj}=b_s^{jk}g^{si}.$$
To write down the compatibility conditions of (2) and (3), we replace $`g^{ij}`$ and $`b_k^{ij}`$ by the linear combinations
$$g^{ij}\lambda g^{ij}+\stackrel{~}{g}^{ij},b_k^{ij}\lambda b_k^{ij}+\stackrel{~}{b}_k^{ij},$$
substitute them into (8), (9), collect the terms with $`\lambda `$ (terms with $`\lambda ^2`$ and $`\lambda ^0`$ vanish since (2) and (3) are Hamiltonian), and equate them to zero. Thus, (8) produces the first compatibility condition
$$\begin{array}{c}2\stackrel{~}{b}_s^{ki}g^{sj}+2b_s^{ki}\stackrel{~}{g}^{sj}=\\ \\ \stackrel{~}{g}^{js}_sg^{ik}+g^{js}_s\stackrel{~}{g}^{ik}+\stackrel{~}{g}^{ks}_sg^{ij}+g^{ks}_s\stackrel{~}{g}^{ij}\stackrel{~}{g}^{is}_sg^{kj}g^{is}_s\stackrel{~}{g}^{kj}.\end{array}$$
(10)
Similarly, (9) produces the second compatibility condition
$$\begin{array}{c}\stackrel{~}{g}^{js}_sb_n^{ik}+g^{js}_s\stackrel{~}{b}_n^{ik}\stackrel{~}{g}^{is}_sb_n^{jk}g^{is}_s\stackrel{~}{b}_n^{jk}+\\ \\ (\stackrel{~}{b}_s^{ij}\stackrel{~}{b}_s^{ji})b_n^{sk}+(b_s^{ij}b_s^{ji})\stackrel{~}{b}_n^{sk}+\stackrel{~}{b}_s^{ik}b_n^{js}+b_s^{ik}\stackrel{~}{b}_n^{js}\stackrel{~}{b}_s^{jk}b_n^{is}b_s^{jk}\stackrel{~}{b}_n^{is}=0.\end{array}$$
(11)
To simplify further calculations, it is convenient to work in the coordinates where the flat metric $`g`$ assumes the constant coefficient form $`g^{ij}=const`$, so that $`b_k^{ij}0`$. In these coordinates the compatibility conditions (10), (11) reduce to
$$2\stackrel{~}{b}_s^{ki}g^{sj}=g^{js}_s\stackrel{~}{g}^{ik}+g^{ks}_s\stackrel{~}{g}^{ij}g^{is}_s\stackrel{~}{g}^{kj}$$
(12)
and
$$g^{js}_s\stackrel{~}{b}_n^{ik}g^{is}_s\stackrel{~}{b}_n^{jk}=0,$$
(13)
respectively. Rewriting the left-hand side of (12) in the form $`2\stackrel{~}{b}_s^{ki}g^{sj}=2\stackrel{~}{b}_s^{ki}\stackrel{~}{g}^{sl}(r^1)_l^j`$ and substituting the expressions for $`\stackrel{~}{b}_s^{ki}\stackrel{~}{g}^{sl}`$ from (8), we arrive at
$$(\underset{¯}{\stackrel{~}{g}^{ls}_s\stackrel{~}{g}^{ik}}+\stackrel{~}{g}^{ks}_s\stackrel{~}{g}^{il}\stackrel{~}{g}^{is}_s\stackrel{~}{g}^{kl})(r^1)_l^j=\underset{¯}{g^{js}_s\stackrel{~}{g}^{ik}}+g^{ks}_s\stackrel{~}{g}^{ij}g^{is}_s\stackrel{~}{g}^{kj}.$$
Cancelling the underlined terms and substituting $`\stackrel{~}{g}^{ij}=r_s^ig^{sj}=r_s^jg^{si}`$, we obtain
$$(r_p^sg^{pk}_s(r_n^lg^{ni})r_p^sg^{pi}_s(r_n^lg^{nk}))(r^1)_l^j=g^{ks}_s(r_n^jg^{ni})g^{is}_s(r_n^jg^{nk}),$$
or
$$(r_p^sg^{pk}g^{ni}_sr_n^lr_p^sg^{pi}g^{nk}_sr_n^l)(r^1)_l^j=g^{ks}g^{ni}_sr_n^jg^{is}g^{nk}_sr_n^j.$$
Contraction with $`r_j^m`$ results in
$$g^{pk}g^{ni}r_p^s_sr_n^mg^{pi}g^{nk}r_p^s_sr_n^m=g^{ks}g^{ni}r_j^m_sr_n^jg^{is}g^{nk}r_j^m_sr_n^j$$
which is equivalent to
$$g^{pk}g^{ni}(r_p^s_sr_n^mr_n^s_sr_p^mr_j^m_pr_n^j+r_j^m_nr_p^j)=0,$$
implying the vanishing of the Nijenhuis tensor.
To establish the second identity (9), we will make use of the formula (7) for the coefficients $`\stackrel{~}{b}_k^{ij}`$ in terms of $`r`$, the proof of which is included in the Appendix (note that this formula is true in arbitrary coordinate system). In the coordinates where $`g^{ij}=const`$ we have $`g^{is}_s=^i,2\stackrel{~}{b}_k^{ij}=^ir_k^j^jr_k^i+_kr^{ij},`$ so that (13) takes the form
$$^j(\underset{¯}{^ir_n^k}^kr_n^i+_nr^{ik})^i(\underset{¯}{^jr_n^k}^kr_n^j+_nr^{jk})=0.$$
Cancellation of the underlined terms and contraction with $`g^{sn}`$ produces (6). This completes the proof of the Theorem.
Remark. If the spectrum of $`r_j^i`$ is simple, condition (6) is redundant: it is automatically satisfied by virtue of (5) and the flatness of both metrics $`g`$ and $`\stackrel{~}{g}`$. This was the motivation for me to drop condition (6) in the compatibility criterion formulated in . However, in this general form the criterion proved to be incorrect: recently it was pointed out by Mokhov that in the case when the spectrum of $`r_j^i`$ is not simple the vanishing of the Nijenhuis tensor is no longer sufficient for the compatibility.
## 3 Compatibility conditions in the diagonal form: the Lax pairs
If the spectrum of $`r_j^i`$ is simple, the vanishing of the Nijenhuis tensor implies the existence of the coordinates $`R^1,\mathrm{},R^n`$ in which the objects $`r_j^i,g^{ij},\stackrel{~}{g}^{ij}`$ become diagonal. Moreover, the $`i`$-th eigenvalue of $`r_j^i`$ depends only on the coordinate $`R^i`$, so that $`r_j^i=diag(\eta _i),g^{ij}=diag(g^{ii}),\stackrel{~}{g}^{ij}=diag(g^{ii}\eta _i)`$ where $`\eta _i`$ is a function of $`R^i`$. This is a generalization of the analogous observation by Dubrovin in the particular case of compatible Poisson brackets originating from the theory of the associativity equations. Introducing the Lame coefficients $`H_i`$ and the rotation coefficients $`\beta _{ij}`$ by the formulae
$$H_i=\sqrt{g_{ii}}=1/\sqrt{g^{ii}},_iH_j=\beta _{ij}H_i,$$
(14)
we can rewrite the zero curvature conditions for the metric $`g`$ in the form
$$_k\beta _{ij}=\beta _{ik}\beta _{kj},$$
(15)
$$_i\beta _{ij}+_j\beta _{ji}+\underset{ki,j}{}\beta _{ki}\beta _{kj}=0.$$
(16)
The zero curvature condition for the metric $`\stackrel{~}{g}`$ imposes the additional constraint
$$\eta _i_i\beta _{ij}+\eta _j_j\beta _{ji}+\frac{1}{2}\eta _i{}_{}{}^{}\underset{¯}{\text{e}}ta_{ij}+\frac{1}{2}\eta _j{}_{}{}^{}\beta _{ji}^{}+\underset{ki,j}{}\eta _k\beta _{ki}\beta _{kj}=0,$$
(17)
resulting from (16) after the substitution of the rotation coefficients $`\stackrel{~}{\beta }_{ij}=\beta _{ij}\sqrt{\eta _i/\eta _j}`$ of the metric $`\stackrel{~}{g}`$. As can be readily seen, equations (16) and (17) already imply the compatibility, so that in the diagonalisable case condition (6) of Theorem 1 is indeed superfluous. Solving equations (16), (17) for $`_i\beta _{ij}`$, we can rewrite (15)–(17) in the form
$$\begin{array}{c}_k\beta _{ij}=\beta _{ik}\beta _{kj},\\ \\ _i\beta _{ij}=\frac{1}{2}\frac{\eta _i^{}}{\eta _j\eta _i}\beta _{ij}+\frac{1}{2}\frac{\eta _j^{}}{\eta _j\eta _i}\beta _{ji}+_{ki,j}\frac{\eta _k\eta _j}{\eta _j\eta _i}\beta _{ki}\beta _{kj}.\end{array}$$
(18)
It can be verified by a straightforward calculation that system (18) is compatible for any choice of the functions $`\eta _i(R^i)`$, and its general solution depends on $`n(n1)`$ arbitrary functions of one variable (indeed, one can arbitrarily prescribe the value of $`\beta _{ij}`$ on the j-th coordinate line). Under the additional ”Egorov” assumption $`\beta _{ij}=\beta _{ji}`$, system (18) reduces to the one studied by Dubrovin in . For $`n3`$ system (18) is essentially nonlinear. Its integrability follows from the Lax pair
$$_j\psi _i=\beta _{ij}\psi _j,_i\psi _i=\frac{\eta _i^{}}{2(\lambda +\eta _i)}\psi _i\underset{ki}{}\frac{\lambda +\eta _k}{\lambda +\eta _i}\beta _{ki}\psi _k$$
(19)
with a spectral parameter $`\lambda `$ (another demonstration of the integrability of system (19) has been proposed recently in by an appropriate modification of Zakharov’s approach ).
Remark. In fact, the Lax pair (19) is gauge-equivalent to the equations
$$(\stackrel{~}{g}^{ik}+\lambda g^{ik})_k_j\psi +(\stackrel{~}{b}_j^{ik}+\lambda b_j^{ik})_k\psi =0$$
for the Casimirs $`\psi 𝑑x`$ of the Hamiltonian operator $`\stackrel{~}{A}^{ij}+\lambda A^{ij}`$.
After the gauge transformation $`\psi _i=\phi _i/\sqrt{\lambda +\eta _i}`$ the Lax pair (19) assumes the manifestly skew-symmetric form
$$_j\phi _i=\sqrt{\frac{\lambda +\eta _i}{\lambda +\eta _j}}\beta _{ij}\phi _j,_i\phi _i=\underset{ki}{}\sqrt{\frac{\lambda +\eta _k}{\lambda +\eta _i}}\beta _{ki}\phi _k,$$
(20)
which is of the type discussed in . Thus, we can introduce an orthonormal frame $`\stackrel{}{\phi }_1,\mathrm{},\stackrel{}{\phi }_n`$ in the Euclidean space $`E^n`$ satisfying the equations
$$_j\stackrel{}{\phi }_i=\sqrt{\frac{\lambda +\eta _i}{\lambda +\eta _j}}\beta _{ij}\stackrel{}{\phi }_j,_i\stackrel{}{\phi }_i=\underset{ki}{}\sqrt{\frac{\lambda +\eta _k}{\lambda +\eta _i}}\beta _{ki}\stackrel{}{\phi }_k,(\stackrel{}{\phi }_i,\stackrel{}{\phi }_j)=\delta _{ij}.$$
(21)
Let us introduce a vector $`\stackrel{}{r}`$ such that
$$_i\stackrel{}{r}=\frac{H_i}{\sqrt{\lambda +\eta _i}}\stackrel{}{\phi }_i$$
(the compatibility of these equations can be readily verified). In view of the formula
$$(_i\stackrel{}{r},_j\stackrel{}{r})=\frac{H_i^2}{\lambda +\eta _i}\delta _{ij},$$
the radius-vector $`\stackrel{}{r}`$ is descriptive of an n-orthogonal coordinate system in $`E^n`$ corresponding to the flat metric
$$\underset{i}{}\frac{H_i^2}{\lambda +\eta _i}(dR^i)^2.$$
Geometrically, $`\stackrel{}{\phi }_i`$ are the unit vectors along the coordinate lines of this n-orthogonal system.
Let us discuss in some more detail the case $`\eta _i=const=c_i`$, in which system (18) takes the form
$$\begin{array}{c}_k\beta _{ij}=\beta _{ik}\beta _{kj},\\ \\ _i\beta _{ij}=_{ki,j}\frac{c_kc_j}{c_jc_i}\beta _{ki}\beta _{kj}.\end{array}$$
(22)
One can readily verify that the quantity
$$P_i=\underset{ki}{}(c_kc_i)\beta _{ki}^2$$
is an integral of system (22), namely, $`_jP_i=0`$ for any $`ij`$, so that $`P_i`$ is a function of $`R^i`$. Utilising the obvious symmetry $`R^is_i(R^i),\beta _{ki}\beta _{ki}/s_i^{}(R^i)`$ of system (22), we can reduce $`P_i`$ to $`\pm 1`$ (if nonzero). Let us consider the simplest nontrivial case $`n=3,P_1=P_2=1,P_3=1`$:
$$\begin{array}{c}P_1=(c_2c_1)\beta _{21}^2+(c_3c_1)\beta _{31}^2=1,\\ \\ P_2=(c_1c_2)\beta _{12}^2+(c_3c_2)\beta _{32}^2=1,\\ \\ P_3=(c_1c_3)\beta _{13}^2+(c_2c_3)\beta _{23}^2=1.\end{array}$$
Assuming $`c_3>c_2>c_1`$ and introducing the parametrization
$$\begin{array}{c}\beta _{21}=\mathrm{sin}p/\sqrt{c_2c_1},\beta _{31}=\mathrm{cos}p/\sqrt{c_3c_1},\\ \\ \beta _{12}=\mathrm{sinh}q/\sqrt{c_2c_1},\beta _{32}=\mathrm{cosh}q/\sqrt{c_3c_2},\\ \\ \beta _{13}=\mathrm{sin}r/\sqrt{c_3c_1},\beta _{23}=\mathrm{cos}r/\sqrt{c_3c_2},\end{array}$$
we readily rewrite (22) in the form
$$\begin{array}{c}_1q=\mu _1\mathrm{cos}p,_1r=\mu _1\mathrm{sin}p,\\ \\ _2p=\mu _2\mathrm{cosh}q,_2r=\mu _2\mathrm{sinh}q,\\ \\ _3p=\mu _3\mathrm{cos}r,_3q=\mu _3\mathrm{sin}r,\end{array}$$
where
$$\mu _1=\sqrt{\frac{c_3c_2}{(c_2c_1)(c_3c_1)}},\mu _2=\sqrt{\frac{c_3c_1}{(c_2c_1)(c_3c_2)}},\mu _3=\sqrt{\frac{c_2c_1}{(c_3c_1)(c_3c_2)}}.$$
After rescaling, this system simplifies to
$$\begin{array}{c}_1q=\mathrm{cos}p,_1r=\mathrm{sin}p,\\ \\ _2p=\mathrm{cosh}q,_2r=\mathrm{sinh}q,\\ \\ _3p=\mathrm{cos}r,_3q=\mathrm{sin}r.\end{array}$$
(23)
Expressing $`p`$ and $`r`$ in the form $`p=\mathrm{arccos}_1q,r=\mathrm{arcsin}_3q,`$ we can rewrite (23) as a triple of pairwise commuting Monge-Ampere equations
$$\begin{array}{c}_1_2q=\mathrm{cosh}q\sqrt{1_1q^2},\\ \\ _1_3q=\sqrt{1_1q^2}\sqrt{1_3q^2},\\ \\ _2_3q=\mathrm{sinh}q\sqrt{1_3q^2}.\end{array}$$
Similar triples of Monge-Ampere equations were obtained in in the classification of quadruples of $`3\times 3`$ hydrodynamic type systems which are closed under the Laplace transformations. However, at the moment there is no explanation of this coincidence.
## 4 Deformations of n-orthogonal coordinate systems inducing rescalings of the Weingarten operators of the coordinate hypersurfaces
We have demonstrated in sect. 3 that the radius-vector $`\stackrel{}{r}(R^1,\mathrm{},R^n)`$ of the n-orthogonal coordinate system in $`E^n`$ corresponding to the flat diagonal metric $`_i\frac{H_i^2}{\lambda +\eta _i}(dR^i)^2`$ satisfies the equations
$$_i\stackrel{}{r}=\frac{H_i}{\sqrt{\lambda +\eta _i}}\stackrel{}{\phi }_i,$$
where the infinitesimal displacements of the orthonormal frame $`\stackrel{}{\phi }_i`$ are governed by
$$_j\stackrel{}{\phi }_i=\sqrt{\frac{\lambda +\eta _i}{\lambda +\eta _j}}\beta _{ij}\stackrel{}{\phi }_j,_i\stackrel{}{\phi }_i=\underset{ki}{}\sqrt{\frac{\lambda +\eta _k}{\lambda +\eta _i}}\beta _{ki}\stackrel{}{\phi }_k.$$
Since our formulae depend on the spectral parameter, we may speak of the ”deformation” of the n-orthogonal coordinate system with respect to $`\lambda `$. To investigate this deformation in some more detail, we fix a coordinate hypersurface $`M^{n1}E^n`$ (say, $`R^n=const`$). Its radius-vector $`\stackrel{}{r}`$ and the unit normal $`\stackrel{}{\phi }_n`$ satisfy the Weingarten equations
$$_i\stackrel{}{\phi }_n=\frac{\beta _{ni}}{H_i}\sqrt{\lambda +\eta _n}_i\stackrel{}{r},i=1,\mathrm{},n1,$$
implying that
$$k^i=\frac{\beta _{ni}}{H_i}\sqrt{\lambda +\eta _n}$$
are the principal curvatures of $`M^{n1}`$. Since $`\eta _n`$ is a constant on $`M^{n1}`$, our deformation preserves the Weingarten operator of $`M^{n1}`$ up to a constant scaling factor $`\sqrt{\lambda +\eta _n}`$ (we point out that the curvature line parametrization $`R^1,\mathrm{},R^{n1}`$ is preserved by a construction). Thus, compatible Poisson brackets of hydrodynamic type give rise to deformations of n-orthogonal systems in $`E^n`$ which, up to scaling factors, preserve the Weingarten operators of the coordinate hypersurfaces. If we follow the evolution of a particular coordinate hypersurface $`M^{n1}`$, this scailing factor can be eliminated by a homothetic transformation of the ambient space $`E^n`$, so that we arrive at the nontrivial deformation of a hypersurface which preserves the Weingarten operator. However, this scaling factor cannot be eliminated for all coordinate hypersurfaces simultaneously.
## 5 Surfaces in $`E^3`$ which possess nontrivial deformations preserving the Weingarten operator
Interestingly enough, the problem of the classification of surfaces $`M^2E^3`$ which possess nontrivial deformations preserving the Weingarten operator has been formulated by Finikov and Gambier as far back as in 1933 , . Among other results, they demonstrated that the only surfaces possessing 3-parameter families of such deformations are the quadrics, conformal transforms of surfaces of revolution and all other surfaces having the same spherical image of curvature lines (if surfaces have the same spherical image of curvature lines or, equivalently, related by a Combescure transformation, they can be deformed simultaneously).
In this section we discuss surfaces which possess 1-parameter families of such deformations. Let $`M^2E^3`$ be a surface parametrized by coordinates $`R^1,R^2`$ of curvature lines. Let
$$G_{11}(dR^1)^2+G_{22}(dR^2)^2$$
(24)
be its third fundamental form (or metric of the Gaussian image, which is automatically of constant curvature 1). Let $`k^1,k^2`$ be the radii of principal curvature satisfying the Peterson-Codazzi equations
$$\frac{_2k^1}{k^2k^1}=_2\mathrm{ln}\sqrt{G_{11}},\frac{_1k^2}{k^1k^2}=_1\mathrm{ln}\sqrt{G_{22}}.$$
(25)
Suppose there exists a flat metric
$$g_{11}(dR^1)^2+g_{22}(dR^2)^2$$
(26)
such that
$$G_{11}=g_{11}/\eta _1,G_{22}=g_{22}/\eta _2,$$
(27)
where $`\eta _1,\eta _2`$ are functions of $`R^1,R^2`$, respectively. One can readily verify that under these assumptions the metric
$$\stackrel{~}{G}_{11}(dR^1)^2+\stackrel{~}{G}_{22}(dR^2)^2=\frac{g_{11}}{\lambda +\eta _1}(dR^1)^2+\frac{g_{22}}{\lambda +\eta _2}(dR^2)^2$$
(28)
has constant curvature $`1`$ for any $`\lambda `$. Since equations (25) are still true if we replace $`G_{ii}`$ by $`\stackrel{~}{G}_{ii}`$, we arrive at a 1-parameter family of surfaces $`M_\lambda ^2`$ with the third fundamental forms (28) (which depend on $`\lambda `$) and the principal curvatures $`k^1,k^2`$ (which are independent of $`\lambda `$). Hence, the Weingarten operators of surfaces $`M_\lambda ^2`$ coincide. The problem of the classification of surfaces which possess 1-parameter families of deformations preserving the Weingarten operator is thus reduced to the classification of metrics (28) which have constant Gaussian curvatures $`1`$ for any $`\lambda `$. Any such metric generates an infinite family of deformable surfaces whose principal curvatures $`k^1,k^2`$ satisfy (25). In terms of the Lame coefficients $`H_1=\sqrt{g_{11}},H_2=\sqrt{g_{22}}`$ and the rotation coefficients $`\beta _{12}=_1H_2/H_1,\beta _{21}=_2H_1/H_2`$ our problem reduces to the nonlinear system
$$\begin{array}{c}_1H_2=\beta _{12}H_1,_2H_1=\beta _{21}H_2,\\ \\ _1\beta _{12}+_2\beta _{21}=0,\\ \\ \eta _1_1\beta _{12}+\eta _2_2\beta _{21}+\frac{1}{2}\eta _1^{}\beta _{12}+\frac{1}{2}\eta _2^{}\beta _{21}+H_1H_2=0,\end{array}$$
(29)
which possesses the Lax pair
$$\begin{array}{c}_1\psi =\left(\begin{array}{ccc}0& \sqrt{\frac{\lambda +\eta _2}{\lambda +\eta _1}}\beta _{21}& \frac{H_1}{\sqrt{\lambda +\eta _1}}\\ \sqrt{\frac{\lambda +\eta _2}{\lambda +\eta _1}}\beta _{21}& 0& 0\\ \frac{H_1}{\sqrt{\lambda +\eta _1}}& 0& 0\end{array}\right)\psi ,\\ \\ _2\psi =\left(\begin{array}{ccc}0& \sqrt{\frac{\lambda +\eta _1}{\lambda +\eta _2}}\beta _{12}& 0\\ \sqrt{\frac{\lambda +\eta _1}{\lambda +\eta _2}}\beta _{12}& 0& \frac{H_2}{\sqrt{\lambda +\eta _2}}\\ 0& \frac{H_2}{\sqrt{\lambda +\eta _2}}& 0\end{array}\right)\psi .\end{array}$$
Geometrically, this Lax pair governs infinitesimal displacements of the orthonormal frame of the orthogonal coordinate system on the unit sphere $`S^2`$, corresponding to the metric (28). In $`2\times 2`$ matrices it takes the form
$$\begin{array}{c}2\sqrt{\lambda +\eta _1}_1\psi =\left(\begin{array}{ccc}i\sqrt{\lambda +\eta _2}\beta _{21}& H_1& \\ H_1& i\sqrt{\lambda +\eta _2}\beta _{21}& \end{array}\right)\psi ,\\ \\ 2\sqrt{\lambda +\eta _2}_2\psi =i\left(\begin{array}{ccc}\sqrt{\lambda +\eta _1}\beta _{12}& H_2& \\ H_2& \sqrt{\lambda +\eta _1}\beta _{12}& \end{array}\right)\psi .\end{array}$$
Remark. In we established a one-to-one correspondence between surfaces possessing nontrivial deformations preserving the Weingarten operator and multi-Hamiltonian systems of hydrodynamic type. Indeed, let us introduce the Hamiltonian operator
$$\delta ^{ij}g^{ii}\frac{d}{dx}+b_k^{ij}R_x^k$$
associated with the flat diagonal metric (26), and the nonlocal Hamiltonian operator (see )
$$\delta ^{ij}G^{ii}\frac{d}{dx}+\stackrel{~}{b}_k^{ij}R_x^k+R_x^i\left(\frac{d}{dx}\right)^1R_x^j,$$
associated with the diagonal metric (24) of constant curvature 1 (these operators are compatible by virtue of (27)). According to the results of Tsarev , , equations (25) imply that the systems of hydrodynamic type
$$R_t^1=k^1(R)R_x^1,R_t^2=k^2(R)R_x^2$$
are automatically bi-Hamiltonian with the respect to both Hamiltonian structures. Characteristic velocities of these systems are the radii of principal curvature of the corresponding surfaces.
The results of this section generalise in a straightforward way to multidimensional hypersurfaces $`M^{n1}E^n`$.
## 6 Appendix: formula for $`\stackrel{~}{b}_k^{ij}`$
To verify formula (7), it suffices to check the identities
$$\stackrel{~}{b}_k^{ij}+\stackrel{~}{b}_k^{ji}=_kr^{ij},$$
(30)
$$\stackrel{~}{b}_s^{ik}r^{sj}=\stackrel{~}{b}_s^{jk}r^{si}.$$
(31)
Substituting the expression for the covariant derivative
$$_kr^{ij}=_kr^{ij}b_k^{si}r_s^jb_k^{sj}r_s^i,$$
into (7), we readily obtain
$$2\stackrel{~}{b}_k^{ij}=(^ir_k^j^jr_k^i+b_k^{sj}r_s^ib_k^{si}r_s^j)+_kr^{ij},$$
where the expression in brackets is skew-symmetric in $`i,j`$. This proves (30).
To verify (31), we first rewrite it in the form
$$(^ir_s^k^kr_s^i+_sr^{ik}+\underset{¯}{2b_s^{lk}r_l^i})r^{sj}=(^jr_s^k^kr_s^j+_sr^{jk}+\underset{¯}{2b_s^{lk}r_l^j})r^{si}.$$
(32)
Since $`b_s^{lk}r_l^ir^{sj}=b_s^{lk}g_{lt}r^{ti}r^{sj}=\mathrm{\Gamma }_{ts}^kr^{ti}r^{sj}`$, the underlined terms cancel in view of the symmetry of $`\mathrm{\Gamma }`$. Contracting (32) with $`g_{pj}g_{mk}g_{ni}`$, we arrive at
$$g_{pj}(_nr_{ms}_mr_{ns})r^{sj}+r_p^s_sr_{mn}=g_{ni}(_pr_{ms}_mr_{ps})r^{si}+r_n^s_sr_{mp},$$
which, by virtue of the identity $`r_{ms}r^{sj}=r_m^sr_s^j`$, transforms to
$$g_{pl}r_s^l(_nr_m^s_mr_n^s)+r_p^s_sr_{mn}=g_{nl}r_s^l(_pr_m^s_mr_p^s)+r_n^s_sr_{mp}.$$
In view of the identity
$$r_s^l(_nr_m^s_mr_n^s)=r_n^s_sr_m^lr_m^s_sr_n^l,$$
manifesting the vanishing of the Nijenhuis tensor (we emphasize that in (5) partial derivatives can be replaced by covariant derivatives with respect to any symmetric affine connection without changing the Nijenhuis tensor), the last equation can be rewritten as follows:
$$g_{pl}(r_n^s_sr_m^lr_m^s_sr_n^l)+r_p^s_sr_{mn}=g_{nl}(r_p^s_sr_m^lr_m^s_sr_p^l)+r_n^s_sr_{mp},$$
or
$$r_n^s_sr_{pm}r_m^s_sr_{pn}+r_p^s_sr_{mn}=r_p^s_sr_{nm}r_m^s_sr_{np}+r_n^s_sr_{mp},$$
which is obviously an identity. This proves formula (7).
## 7 Acknowledgements
I would like to thank O. I. Mokhov for the references , .
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# Tricritical scaling at the 𝑁_𝑡=6 chiral phase transition for 2 flavour lattice QCD with staggered quarks.
## I ACKNOWLEDGEMENTS
These computations were performed on the CRAY C90/J90/SV1’s and T3E at NERSC, and on a T3E at SGI/Cray. We would like to thank J.-F. Lagaë who contributed in the earlier stages of this work. We also thank Dr. M. Stephanov for discussions of tricritical phenomena. This work was supported by the U. S. Department of Energy under contract W-31-109-ENG-38, and the National Science Foundation under grant NSF-PHY96-05199.
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# The heat flow of the CCR algebra
## 1. Discussion, basic results
Consider the canonical operators $`P,Q`$ acting on an appropriate common dense domain in $`L^2()`$
$`P`$ $`={\displaystyle \frac{1}{i}}{\displaystyle \frac{d}{dx}},`$
$`Q`$ $`=\text{multiplication by }x.`$
These operators can be used to define unbounded derivations (say on the dense $``$-algebra $`𝒜`$ of all integral operators having kernels which are smooth and of compact support) by
$$D_P(X)=i(PXXP),D_Q(X)=i(QXXQ),X𝒜.$$
Thinking of these derivations as noncommutative counterparts of $`/x`$ and $`/y`$ we define a “Laplacian” $`L:𝒜𝒜`$ by
$$L=D_P^2+D_Q^2.$$
$`1.1`$
Throughout this paper we will use the term CP semigroup to denote a semigroup $`\varphi =\{\varphi _t:t0\}`$ of normal completely positive linear maps on the algebra $`(H)`$ of all bounded operators on a separable Hilbert space $`H`$, which preserves the unit $`\varphi _t(\mathrm{𝟙})=\mathrm{𝟙}`$, and which is continuous in the natural sense (namely $`\varphi _t(A)\xi ,\eta `$ should be continuous in $`t`$ for fixed $`\xi ,\eta H`$ and $`A(H)`$). The purpose of this section is to exhibit concretely a CP semigroup whose generator can be identified with the operator mapping $`L`$ of (1.1) (see Theorem 1.10).
Let $`U_t=e^{itQ}`$, $`V_t=e^{itP}`$ be the two unitary groups associated with $`Q`$, $`P`$,
$$U_tf(x)=e^{itx}f(x),V_tf(x)=f(x+t),fL^2().$$
These two groups satisfy the Canonical Commutation Relations $`V_tU_s=e^{ist}U_sV_t`$ for $`s,t`$. It is more convenient to make use of the CCRs in Weyl’s form. For every $`z=(x,y)^2`$ the Weyl operator
$$W_z=e^{\frac{ixy}{2}}U_xV_y$$
$`1.2`$
is unitary, it is strongly continuous in $`z`$, and it satisfies the Weyl relations
$$W_{z_1}W_{z_2}=e^{i\omega (z_1,z_2)}W_{z_1+z_2}$$
$`1.3`$
where $`\omega `$ is the symplectic form on $`^2`$ given by
$$\omega ((x,y),(x^{},y^{}))=\frac{1}{2}(x^{}yxy^{}).$$
$`1.4`$
A strongly continuous mapping $`zW_z(H)`$ into the unitary operators on some Hilbert space $`H`$ which satisfies (1.3) is called a Weyl system. It is well known that the Weyl system (1.2) is irreducible, and hence the space of all finite linear combinations of the $`W_z`$ is a unital strongly dense $``$-subalgebra of $`(L^2())`$. The Stone-von Neumann theorem implies that every Weyl system is unitarily equivalent to a direct sum of copies of the concrete Weyl system (1.2).
Proceeding heuristically for a moment, let $`D_P`$ and $`D_Q`$ be the derivations above. After formally differenting the relation $`V_tU_s=e^{ist}U_sV_t`$ we find that
$`D_P(U_x)`$ $`=ixU_x,`$ $`D_P(V_y)`$ $`=0,`$
$`D_Q(U_x)`$ $`=0,`$ $`D_Q(V_y)`$ $`=iyU_y,`$
hence the action of $`L=D_P^2+D_Q^2`$ on the Weyl system (1.2) is given by
$$L(W_z)=(x^2+y^2)W_z=|z|^2W_z,z=(x,y)^2.$$
After formally exponentiating we find that for $`t0`$ the operator mapping $`\varphi _t=\mathrm{exp}(tL)`$ for $`t0`$ can be expected to satisfy
$$\varphi _t(W_z)=e^{t|z|^2}W_z,z^2,t0.$$
$`1.5`$
###### Remark Remarks
A number of authors have considered completely positive semigroups defined on a Weyl system by formulas such as (1.5), using techniques similar to those of Proposition 1.7 below (see pp. 128-129 of , or for two notable examples). We include a full discussion of these basic issues since in section 3 we require details of the construction that are not easily found in the literature.
We also remark that virtually all of the results below have straightforward generalizations to the case in which $`P,Q`$ are replaced with the canonical operators $`P_1,\mathrm{},P_n`$, $`Q_1,\mathrm{},Q_n`$ associated with $`n`$ degrees of freedom. Indeed, the generalization amounts to little more than a reinterpretation of notation. On the other hand, while Proposition 1.7 remains valid (with the same proof) for infinitely many degrees of freedom (c.f. loc cit), we do not know if that is the case for the more precise results of section 3.
In order to define the CCR heat flow rigorously we take (1.5) as our starting point and deduce the existence of the semigroup and its basic properties from the following general result. Consider the Banach space $`M(^2)`$ of all complex-valued measures $`\mu `$ on $`^2`$ having finite total variation $`\mu `$. $`M(^2)`$ is a commutative Banach algebra with unit relative to the usual convolution of measures
$$\mu \nu (S)=_{^2\times ^2}\chi _S(z+w)𝑑\mu (z)𝑑\nu (w).$$
It will be convenient to define the Fourier transform of a measure $`\mu M(^2)`$ in terms of the symplectic form $`\omega `$ of (1.4)
$$\widehat{\mu }(\zeta )=_^2e^{i\omega (\zeta ,z)}𝑑\mu (z).$$
$`1.6`$
###### Remark Remark
While this definition of the Fourier transform differs from the usual one, which involves the Euclidean inner product of $`^2`$
$$(x,y),(x^{},y^{})=xx^{}+yy^{}$$
rather than the symplectic form $`\omega `$, it is equivalent to it in a natural way. Indeed, since $`\omega `$ is nondegenerate there is a unique invertible skew symmetric linear operator $`\mathrm{\Omega }`$ on the two dimensional real vector space $`^2`$ satisfying $`\omega (z,z^{})=\mathrm{\Omega }z,z^{}`$ for all $`z,z^{}^2`$. Hence one can pass back and forth from the usual Fourier transform of a measure to the one above by the invertible linear change-of-variables given by composing the transformed measure with either $`\mathrm{\Omega }`$ or $`\mathrm{\Omega }^1=4\mathrm{\Omega }`$.
###### Proposition 1.7
Let $`\{W_z:z^2\}`$ be an irreducible Weyl system acting on a Hilbert space $`H`$. For every complex measure $`\mu M(^2)`$ there is a unique normal completely bounded linear map $`\varphi _\mu :(H)(H)`$ satisfying
$$\varphi _\mu (W_z)=\widehat{\mu }(z)W_z,z^2.$$
One has $`\varphi _\mu \varphi _\nu =\varphi _{\mu \nu }`$, and $`\varphi _\mu _{cb}\mu `$, where $`\psi _{cb}`$ denotes the completely bounded norm of an operator mapping $`\psi `$. When $`\mu `$ is a positive measure $`\varphi _\mu `$ is a completely positive map.
###### Demonstration proof
Fix $`\mu M(^2)`$. The uniqueness of the mapping $`\varphi _\mu `$ is apparent from the irreducibility hypothesis on the Weyl system, since the set of all linear combinations of the $`W_z`$, $`z^2`$, is a unital $``$-algebra which is weak-dense in $`(H)`$.
For existence, we exhibit $`\varphi _\mu (A)`$ for $`A(H)`$ as a weak integral
$$\varphi _\mu (A)=_^2W_{2^{1/2}\zeta }AW_{2^{1/2}\zeta }^{}𝑑\mu (\zeta ),$$
$`1.8`$
namely the operator defined by the bounded sesquilinear form on the right of (1.9)
$$\varphi _\mu (A)\xi ,\eta =_^2W_{2^{1/2}\zeta }AW_{2^{1/2}\zeta }^{}\xi ,\eta 𝑑\mu (\zeta ),\xi ,\eta H.$$
$`1.9`$
A straightforward estimate shows that $`\varphi _\mu (A)A\mu `$, and after promoting $`\varphi _\mu `$ to $`n\times n`$ matrices over $`(H)`$ a similar estimate shows that $`\varphi _\mu _{cb}\mu `$. $`\varphi _\mu `$ is obviously completely positive when $`\mu `$ is a positive measure.
Formula (1.9), together with a straightforward application of the bounded convergence theorem, implies that when $`A_1,A_2,\mathrm{}`$ is a (necessarily bounded) sequence in $`(H)`$ which converges weakly to $`A`$, one has
$$\underset{n\mathrm{}}{lim}\varphi _\mu (A)\xi ,\eta =\varphi _\mu (A)\xi ,\eta ,\xi ,\eta H.$$
It follows that $`\varphi _\mu `$ is a normal linear map.
Finally, from the commuation relation (1.3) we find that
$$W_{2^{1/2}\zeta }W_zW_{2^{1/2}\zeta }^{}=W_{2^{1/2}\zeta }W_zW_{2^{1/2}\zeta }=e^{\omega (\zeta ,z)}W_z,$$
and hence (1.8) implies that $`\varphi _\mu (W_z)=\widehat{\mu }(z)W_z`$.
I want to thank Daniel Markiewicz for a suggestion that simplified the proof of Proposition 1.7.
###### Theorem 1.10
Let $`W=\{W_z:z^2\}`$ be an irreducible Weyl system. Then there is a unique CP semigroup $`\varphi =\{\varphi _t:t0\}`$ satisfying
$$\varphi _t(W_z)=e^{t|z|^2}W_z,z^2.$$
$`1.11`$
The only bounded normal linear functional $`\rho `$ for which $`\rho \varphi _t=\rho `$ for all $`t0`$ is $`\rho =0`$. In particular, there is no normal state of $`(H)`$ which is invariant under $`\varphi `$.
###### Demonstration proof
For each $`t0`$, $`u_t(z)=e^{t|z|^2}`$ is a continuous function of positive type, which takes the value $`1`$ at $`z=0`$. Thus it is the Fourier transform of a unique probability measure $`\mu _tM(^2)`$. We will require an explicit formula for the Gaussian measure $`\mu _t`$ later on; but for purposes of this section we require nothing more than its existence and uniqueness.
Since $`u_s(z)u_t(z)=u_{s+t}(z)`$ for all $`z^2`$ it follows that $`\mu _s\mu _t=\mu _{s+t}`$. Hence Proposition 1.7 implies that there is a semigroup $`\varphi =\{\varphi _t:t0\}`$ of normal completely positive maps on $`(H)`$ which satisfies (1.11). It is a simple matter to check that the required continuity of $`\varphi _t`$ in $`t`$ follows from the continuity of the right side of (1.11) in $`t`$ for fixed $`z`$.
Suppose now that $`\rho `$ is a normal linear functional which is invariant under $`\varphi `$. Then for every $`z^2`$ and every $`t0`$, the definition of $`\varphi _t`$ implies that
$$\rho (W_z)=\rho (\varphi _t(W_z))=e^{t|z|^2}\rho (W_z)$$
and for fixed $`z0`$, the right side tends to $`0`$ as $`t\mathrm{}`$. Hence $`\rho (W_z)=0`$ for every $`z0`$; by strong continuity on the unit ball it follows that $`\rho (\mathrm{𝟙})=\omega (W_0)=0`$. hence $`\rho `$ vanishes on the irreducible $``$-algebra spanned by $`W_z`$, $`z^2`$ and by normality it follows that $`\rho =0`$.
###### Remark Remarks
We point out that while $`\varphi `$ has no normal invariant states, it does have a normal invariant weight…namely the trace, in that
$$\text{trace}(\varphi _t(A))=\text{trace}(A)$$
for every positive operator $`A(L^2())`$ and every $`t0`$. One sees this immediately from (1.8). It follows that $`\varphi _t`$ leaves the $`C^{}`$-algebra $`𝒦`$ of all compact operators invariant, $`\varphi _t(𝒦)𝒦`$. Since $`K`$ is the $`C^{}`$-algebra associated with the canonical commuation relations (more precisely, $`𝒦`$ is the enveloping $`C^{}`$-algebra of the Banach $``$-algebra of all Weyl integral operators associated with the canonical commutation relations with a finite number of degrees of freedom), this justifies viewing the semigroup of restrictions $`\{\varphi _t_𝒦:t0\}`$ as the heat flow of the canonical commutation relations.
We also remark that one can deduce the existence of other CP semigroups along similar lines. For example, the proof of Theorem 1.10 implies that there is a “Cauchy” semigroup $`\psi =\{\psi _t:t0\}`$ which is defined uniquely by the requirement
$$\psi _t(W_z)=e^{t(|x|+|y|)}W_z,z=(x,y)^2,$$
and which has properties similar to those discussed above for $`\varphi =\{\varphi _t:t0\}`$.
## 2. Harmonic analysis of the commutation relations
A classical theorem of Beurling asserts that singletons obey spectral synthesis. More precisely, if $`G`$ is a locally compact abelian group and $`f`$ is an integrable function on $`G`$ whose Fourier transform vanishes at a point $`p`$ in the dual of $`G`$, then there is a sequence of functions $`f_nL^1(G)`$ such that $`ff_n0`$ as $`n\mathrm{}`$ and such that the Fourier transform of each $`f_n`$ vanishes identically on some open neighborhood $`U_n`$ of $`p`$. The purpose of this section is to present a noncommutative version of that result, which will be required in section 3.
Let $`\{W_z:z^2\}`$ be an irreducible Weyl system acting on a Hilbert space $`H`$ (for example, one may take the Weyl system (1.2) acting on $`L^2(^2)`$). For every trace-class operator $`A^1(H)`$ we consider the following analogue of the Fourier transform $`\widehat{A}:\widehat{}^2`$
$$\widehat{A}(z)=\text{trace}(AW_z),z^2.$$
This transform $`A^1(H)\widehat{A}`$ shares many features in common with the commutative Fourier transform. For example, using the concrete realization (1.2), it is quite easy to establish a version of the Riemann-Lebesgue lemma
$$\underset{|z|\mathrm{}}{lim}\widehat{A}(z)=0,$$
for every $`A^1(H)`$. What we actually require is the following analogue of Beurling’s theorem, which lies somewhat deeper.
###### Theorem 2.1
Let $`A^1(H)`$ and let $`\zeta ^2`$ be such that $`\text{trace}(AW_\zeta )=0`$. There is a sequence $`A_n^1(H)`$ and a sequence of open neighborhoods $`U_n`$ of $`\zeta `$ such that
$$\text{trace}(A_nW_z)=0,zU_n,$$
and such that $`\text{trace}|AA_n|0`$ as $`n\mathrm{}`$.
###### Demonstration proof
By replacing $`A`$ with $`AW_\zeta `$ and making obvious use of the canonical commutation relations (1.3), we may immediately reduce to the case $`\zeta =0`$. We find it more convenient to establish the dual assertion of Theorem 2.1. For that, consider the following linear subspaces of $`(H)`$
$$𝒮_ϵ=\overline{\text{span}}\{W_z:|z|ϵ\},ϵ>0,$$
the closure being taken relative to the weak topology on $`(H)`$. Obviously the spaces $`𝒮_ϵ`$ decrease as $`ϵ`$ decreases, and the identity operator belongs to $`𝒮_ϵ`$ for every $`ϵ>0`$. The pre-annihilator of $`𝒮_ϵ`$ is identified with the space of all trace-class operators $`A`$ satisfying
$$\widehat{A}(z)=\text{trace}(AW_z)=0,|z|ϵ.$$
$`2.2`$
###### Lemma 2.3
Let $`\{W_z:z^2\}`$ be an arbitrary Weyl system acting on a separable Hilbert space $`H`$. Then $`\{𝒮_ϵ:ϵ>0\}=\mathrm{𝟙}.`$
###### Demonstration proof of Lemma 2.3
Let $`𝒮_0`$ denote the intersection $`\{𝒮_ϵ:ϵ>0\}`$. We have already remarked that the inclusion $``$ is obvious. For the opposite one, consider the von Neumann algebra $``$ generated by $`\{W_z:z^2\}`$. $`M`$ is a factor (of type $`I_{\mathrm{}}`$) because of the Stone-von Neumann theorem. We will show that $`𝒮_0`$ is contained in the center of $``$.
For that, choose $`T𝒮_0`$ and consider the operator-valued function $`zW_zTW_z^{}`$. We have to show that this function is constant; equivalently, we will show that for fixed $`\xi `$ and $`\eta `$ in $`H`$, the function
$$z^2W_zTW_z^{}\xi ,\eta $$
$`2.4`$
is constant. Since the function of (2.4) is bounded and continuous, it suffices to show that its spectrum (in the sense of spectral synthesis for functions in $`L^{\mathrm{}}(^2)`$) is the singleton $`\{0\}`$: this is the dual formulation of Beurling’s theorem cited above. Thus we have to show that for every function $`fL^1(^2)`$ whose Fourier transform
$$\widehat{f}(\zeta )=_^2e^{i\omega (z,\zeta )}f(z)𝑑z$$
vanishes throughout a neighborhood of the origin $`\zeta =0`$, we have
$$_^2f(z)W_zTW_z\xi ,\eta 𝑑z=0.$$
$`2.5`$
Fix such an $`fL^1(^2)`$ and choose $`ϵ>0`$ small enough so that $`\widehat{f}(\zeta )=0`$ for all $`\zeta `$ satisfying $`|\zeta |ϵ`$. Since the linear functional
$$X(H)_^2f(z)W_zXW_z^{}\xi ,\eta 𝑑z$$
is weak-continuous and $`T`$ belongs to the weak-closed linear span of operators of the form $`W_\zeta `$ with $`|\zeta |ϵ`$, to prove (2.5) it suffices to show that for every $`\zeta `$ with $`|\zeta |ϵ`$ we have
$$_^2f(z)W_zW_\zeta W_z^{}\xi ,\eta 𝑑z=0.$$
$`2.6`$
Using the canonical commutation relations we can write
$$W_zW_\zeta W_z^{}=e^{i\omega (z,\zeta )}W_{z+\zeta }W_z=e^{\omega (\zeta ,z)}W_\zeta =e^{i\omega (z,\zeta )}W_\zeta .$$
Hence the left side of (2.6) becomes
$$_^2f(z)e^{i\omega (z,\zeta )}W_\zeta \xi ,\eta 𝑑z=\widehat{f}(\zeta )W_\zeta \xi ,\eta ,$$
and the latter term vanishes because $`\widehat{f}(\zeta )=0`$ for $`|\zeta |ϵ`$.
To complete the proof of Theorem 2.1, choose an operator $`A^1(H)`$ satisfying
$$\widehat{A}(0)=\text{trace}(A)=0,$$
and consider the linear functional $`\rho `$ defined on $`(H)`$ by $`\rho (T)=\text{trace}(AT)`$. $`\rho `$ obviously vanishes on $`\mathrm{𝟙}`$. The linear spaces $`𝒮_ϵ`$ are weak-closed and they decrease to $`𝒮_0=\mathrm{𝟙}`$ as $`ϵ`$ decreases to $`0`$, by Lemma 2.3. Since $`\rho `$ is weak-continuous we must have
$$\underset{ϵ0}{lim}\rho _{𝒮_ϵ}=\rho _\mathrm{𝟙}=0.$$
Thus we can choose a sequence $`ϵ_n0`$ so that $`\rho _{𝒮_{ϵ_n}}1/n`$ for every $`n=1,2,\mathrm{}`$. We have already pointed out that the pre-annihilator of $`𝒮_{ϵ_n}`$ is identified with all trace class operators $`B`$ satisfying
$$\widehat{B}(z)=\text{trace}(BW_z)=0,|z|ϵ_n.$$
$`2.7`$
Since $`\rho _{𝒮_{ϵ_n}}`$ is the trace-norm distance from $`A`$ to the pre-annihilator of $`𝒮_{ϵ_n}`$, we conclude that there is a sequence of operators $`B_n^1(H)`$ which satisfy $`\text{trace}(B_nW_z)=0`$ for $`|z|ϵ_n`$, such that $`\text{trace}|AB_n|2/n`$, as asserted.
## 3. Purity and dilation theory
An $`E_0`$-semigroup is a CP semigroup $`\alpha =\{\alpha _t:t0\}`$, acting on $`(H)`$, such the the individual maps are endomorphisms, $`\alpha _t(AB)=\alpha _t(A)\alpha _t(B)`$, $`A,B(H)`$. An $`E_0`$-semigroup $`\alpha `$ is called pure if its “tail” von Neumann algebra is trivial,
$$\underset{t0}{}\alpha _t((H))=\mathrm{𝟙}.$$
$`3.1`$
It is known that an $`E_0`$-semigroup is pure iff for any pair of normal states $`\rho _1,\rho _2`$ of $`(H)`$ we have
$$\underset{t\mathrm{}}{lim}\rho _1\alpha _t\rho _2\alpha _t=0$$
$`3.2`$
see .
If a pure $`E_0`$-semigroup $`\alpha `$ has a normal invariant state $`\omega `$, then the characterization (3.2) implies that $`\omega `$ must be an absorbing state in the sense that for every normal state $`\rho `$ of $`(H)`$ one has
$$\underset{t\mathrm{}}{lim}\rho \alpha _t\omega =0.$$
$`3.3`$
Conversely, if for an arbitrary $`E_0`$-semigroup $`\alpha `$ there is a state $`\omega `$ of $`(H)`$ which is absorbing in the sense that (3.3) is satisfied for every normal state $`\rho `$ of $`(H)`$, then $`\omega `$ must be a normal invariant state, and thus by (3.2) $`\alpha `$ must be a pure $`E_0`$-semigroup.
In the theory of interactions worked out in , pure $`E_0`$-semigroups occupy a central position, especially those for which there is a normal invariant (and therefore absorbing) state. A natural question that emerges from the theory of interactions is whether or not every pure $`E_0`$-semigroup must have a normal invariant state. Now since the state space of $`(H)`$ is weak-compact, a routine application of the Markov-Kakutani fixed point theorem shows shows that every $`E_0`$-semigroup must have invariant states; but invariant states obtained by such methods need not be normal. In this section we exhibit a concrete $`E_0`$-semigroup which is pure but which has no normal invariant states. This is a result which was asserted (without proof) in . This $`E_0`$-semigroup is obtained from the CP semigroup of Theorem 1.10 by a dilation procedure.
In order that the minimal dilation of a CP semigroup to an $`E_0`$-semigroup should satisfy (3.1), it is necessary and sufficient that the CP semigroup should satisfy property (3.2) (see Proposition 3.5). Thus we generalize the definition of pure $`E_0`$-semigroup as follows.
###### Definition 3.4
A CP semigroup $`\varphi `$ acting on $`(H)`$ is called pure if for every pair of normal states $`\rho _1`$, $`\rho _2`$ of $`(H)`$ we have
$$\underset{t\mathrm{}}{lim}\rho _1\varphi _t\rho _2\varphi _t=0.$$
###### Proposition 3.5
Let $`\varphi =\{\varphi _t:t0\}`$ be a pure CP semigroup which has no normal invariant state, and let $`\alpha `$ be its minimal dilation to an $`E_0`$-semigroup. Then $`\alpha `$ satisfies (3.1) and has no normal invariant state.
###### Demonstration proof
The proof is straightforward, but we require results from . Proposition 2.4 of implies that $`\alpha `$ satisfies (3.1).
To see that $`\alpha `$ has no normal invariant state, we can assume that $`\alpha `$ acts on $`(H)`$ for some Hilbert space $`H`$ and that there is a closed subspace $`KH`$ such that $`\varphi `$ is the compression of $`\alpha `$ onto $`(K)=P(H)P`$, $`P`$ denoting the projection of $`H`$ onto $`K`$. We have $`\alpha _t(P)\mathrm{𝟙}`$ because $`\alpha `$ is minimal over $`P`$. So if $`\omega `$ is any normal state of $`(H)`$ which is invariant under $`\alpha `$ then we have
$$\omega (P)=\underset{t\mathrm{}}{lim}\omega (\alpha _t(P))=\omega (\mathrm{𝟙})=1.$$
Thus the restriction of $`\omega `$ to $`(K)=P(H)P`$ defines a normal $`\varphi `$-invariant state on $`(K)`$, contradicting the hypothesis on $`\varphi `$.
In the remainder of this section we show that the CP semigroup defined in Theorem 1.10 is pure. Once that is established, Proposition 3.5 implies that its minimal dilation is an $`E_0`$-semigroup with properties asserted in the discussion above.
###### Theorem 3.6
The CP semigroup $`\varphi `$ defined in (1.11) is pure.
Before giving the proof, we require
###### Lemma 3.7
For each $`t>0`$ let $`\mu _t`$ be the Gaussian measure on $`^2`$ whose Fourier transform (1.6) is given by
$$\widehat{\mu }_t(z)=e^{t|z|^2},z^2,$$
and choose $`\delta >0`$. There is a family $`\nu _t`$, $`t>0`$, of probability measures on $`^2`$ such that
###### Demonstration proof
For $`t>0`$, $`\mu _t`$ is given by $`d\mu _t=u_t(x,y)dxdy`$, where $`u_t`$ is the density
$$u_t(x,y)=\frac{1}{\pi t}e^{\frac{x^2+y^2}{4t}}.$$
Let $`f_t=\sqrt{u_t}`$. $`f_t`$ belongs to $`L^2(^2)`$ and $`f_t_2=1`$. Choose a function $`gL^1(^2)`$ whose Fourier transform
$$\widehat{g}(\zeta )=_^2e^{i\omega (\zeta ,z)}g(z)𝑑z$$
satisfies $`0\widehat{g}(\zeta )1`$ for all $`\zeta `$, and
$$\widehat{g}(\zeta )=\{\begin{array}{cc}1,\hfill & \text{for }0|\zeta |\delta /4\hfill \\ 0,\hfill & \text{for }|\zeta |\delta /2.\hfill \end{array}$$
Consider the convolution $`gf_tL^2(^2)`$ and the positive measure
$$d\nu _t=|gf_t|^2dxdy.$$
$`\nu _t`$ is obviously a positive finite measure. We claim first that the Fourier transform of $`\nu _t`$ lives in the disk $`|\zeta |\delta `$. Indeed, letting $`U_\zeta `$ (resp. $`T_\zeta `$) be the unitary operator on $`L^2(^2)`$ (resp. $`L^2(\widehat{}^2)`$) given by
$$U_\zeta F(z)=e^{i\omega (\zeta ,z)}F(z),T_\zeta G(w)=G(w+\zeta ),$$
we have by the Plancherel theorem
$`\widehat{\nu }_t(\zeta )`$ $`=U_\zeta (gf_t),gf_t_{L^2(^2)}=T_\zeta (\widehat{g}\widehat{f}_t),\widehat{g}\widehat{f}_t_{L^2(\widehat{}^2)}`$
$`={\displaystyle _{\widehat{}^2}}(\widehat{g}\widehat{f}_t)(w+\zeta )\overline{(\widehat{g}\widehat{f}_t)}(w)𝑑w.`$
When $`|\zeta |\delta `$ the integrand on the right vanishes identically in $`w`$ because $`\widehat{g}\widehat{f}_t`$ is supported in the disk of radius $`\delta /2`$. Hence $`\widehat{\nu }_t(\zeta )=0`$ for $`|\zeta |\delta `$.
To establish property (ii), it is enough to show that
$$\underset{t\mathrm{}}{lim}f_tgf_t_2=0,$$
$`3.8`$
since by the Schwarz inequality
$`\mu _t\nu _t`$ $`={\displaystyle _^2}|f_t^2|gf_t|^2|𝑑z{\displaystyle _^2}|f_tgf_t||f_t+|gf_t||𝑑z`$
$`f_tgf_t_2f_t+|gf_t|_2f_tgf_t_2(f_t_2+gf_t_2).`$
To establish (3.8), we use the Plancherel theorem again to write
$$_^2|f_t(z)gf_t(z)|^2𝑑z=_{\widehat{}^2}|\widehat{f}(\zeta )\widehat{g}(\zeta )\widehat{f}_t(\zeta )|^2𝑑\zeta =_^2|1\widehat{g}(\zeta )|^2|\widehat{f}_t|^2𝑑\zeta .$$
The function $`|1\widehat{g}(\zeta )|`$ is bounded above by $`1`$ and it vanishes throughout the disk $`0|\zeta |\delta /4`$. Hence the term on the right is dominated by
$$_{\{|\zeta |\delta /4\}}|\widehat{f}_t(\zeta )|^2𝑑\zeta .$$
$`3.9`$
In order to estimate the integral (3.9) we require the explicit formula
$$f_t(x,y)=\sqrt{u_t(x,y)}=\frac{1}{\sqrt{\pi t}}e^{\frac{x^2+y^2}{8t}}.$$
The Fourier transform of $`f_t`$ has the form
$$\widehat{f}_t(\zeta )=K\sqrt{t}e^{2t|\zeta |^2}$$
where $`K`$ is a positive constant, hence (3.9) evaluates to
$$K^2t_{\{|\zeta |\delta /4\}}e^{4t|\zeta |^2}𝑑\zeta =K^2_{S_t}e^{4(u^2+v^2)}𝑑u𝑑v,$$
where $`S_t=\{(u,v):\sqrt{u^2+v^2}(\delta /4)\sqrt{t}\}`$. As $`t\mathrm{}`$ the sets $`S_t`$ decrease to $`\mathrm{}`$, hence the right side of the previous expression tends to $`0`$, and (3.8) is proved.
The positive measures $`\nu _t`$ are not necessarily probability measures, but in view of the established property (ii), $`\nu _t(^2)`$ must be arbitrarily close to $`\mu _t(^2)=1`$ when $`t`$ is large. Hence we can rescale $`\nu _t`$ in an obvious way to achive $`\nu _t(^2)=1`$ for all $`t>0`$ as well as the properties (i) and (ii) of Lemma 3.7.
###### Demonstration proof of Theorem 3.6
Let $`W_z`$, $`z^2`$ be an irreducible Weyl system acting on a Hilbert space $`H`$, and let $`\varphi =\{\varphi _t:t0\}`$ be the CP semigroup defined by the condition
$$\varphi _t(W_z)=e^{t|z|^2}W_z,z^2.$$
Choose a pair of normal states $`\rho _1`$, $`\rho _2`$ on $`(H)`$, and consider their difference $`\omega =\rho _1\rho _2`$. We have to show that
$$\underset{t\mathrm{}}{lim}\omega \varphi _t=0.$$
$`3.10`$
For that, let $`A`$ be the self-adjoint trace-class operator defined by $`\text{trace}(AT)=\omega (T)`$, $`T(H)`$ and choose $`ϵ>0`$. $`A`$ has trace zero, so by Theorem 2.1, we can find a self-adjoint trace-class operator $`A_0`$ such that $`\text{trace}(A_0W_z)=0`$ for every $`z`$ in some neighborhood $`U`$ of $`z=0`$, and $`\text{trace}|AA_0|ϵ`$. It follows that the normal linear functional $`\omega _0(T)=\text{trace}(A_0T)`$ satisfies $`\omega \omega _0ϵ`$ and $`\omega _0(W_z)=0`$ for $`zU`$.
By Lemma 3.5 we can find probability measures $`\nu _t`$, $`t>0`$, such that $`\widehat{\nu }_t(z)`$ vanishes for $`zU`$ and $`\mu _t\nu _t`$ tends to $`0`$ as $`t\mathrm{}`$. For each $`t>0`$ let $`\psi _t`$ be the completely positive map defined by Proposition 1.7,
$$\psi _t(W_z)=\widehat{\nu }_t(z)W_z,z^2.$$
In order to prove (3.10) we decompose the linear functional $`\omega \varphi _t`$ into a sum of three terms as follows
$$\omega \varphi _t=(\omega \omega _0)\varphi _t+\omega _0(\varphi _t\psi _t)+\omega _0\psi _t.$$
$`3.11`$
The third term on the right of (3.11) is zero because for every $`z^2`$ we have
$$\omega _0(\psi _t(W_z))=\omega _0(\widehat{\nu }_t(z)W_z)=\widehat{\nu }_t(z)\omega _0(W_z)=0,$$
since $`\widehat{\nu }_t(z)`$ vanishes when $`zU`$ and $`\omega _0(W_z)`$ vanishes when $`zU`$ (recall that the linear span of the $`W_z`$ for $`z^2`$ is a strongly dense $``$-subalgebra of $`(H)`$). The first term on the right of (3.11) is estimated for arbitrary $`t`$ by
$$(\omega \omega _0)\varphi _t\omega \omega _0ϵ.$$
In order to estimate the second term, note that
$$\varphi _t\psi _t\mu _t\nu _t$$
$`3.12`$
for every $`t>0`$. Indeed, considering the measure $`\sigma _t=\mu _t\nu _tM(^2)`$, we can write
$$\varphi _t(W_z)\psi _t(W_z)=\widehat{\mu }_t(z)W_z\widehat{\nu }_t(z)W_z=\widehat{\sigma }_t(z)W_z.$$
It follows from Proposition 1.7 that the completely bounded norm of the operator mapping $`\varphi _t\psi _t`$ is at most $`\sigma _t=\mu _t\nu _t`$, hence (3.12).
From (3.11) and these estimates we may conclude that
$$\underset{t\mathrm{}}{lim\; sup}\omega \varphi _tϵ+\underset{t\mathrm{}}{lim}\mu _t\nu _t+0=ϵ.$$
Since $`ϵ`$ is arbitrary the limit (3.10) is proved.
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# Extracting |𝑉_{𝑢𝑏}| from the Hadronic Mass Spectrum of Inclusive 𝐵 decays
## I Introduction
An accurate determination of $`V_{ub}`$ would fill a gaping hole in our understanding of the Cabibbo-Kobayashi-Maskawa (CKM) sector of the standard model. Indeed, while it is not clear that we will be able to cleanly determine the angles of the unitarity triangle, an accurate measurement of $`|V_{ub}|`$ could determine whether or not the CKM description is correct by measuring the lengths of the sides.
While in principle one could measure $`|V_{ub}|`$ quite easily in a systematic approximation to QCD from an inclusive measurement of the uncharmed semi-leptonic branching ratio, in practice, it seems to be much more difficult. The problem lies in rejecting the charmed decay background. This cut causes the canonical expansions in $`\alpha _s`$ and $`\mathrm{\Lambda }/m_b`$ to breakdown. The imposition of the cut changes the expansion parameters to $`\alpha _s\mathrm{log}\rho `$ and $`\mathrm{\Lambda }/(m_b\rho )`$, where $`\rho `$ parameterizes the cut and is numerically small. It may be the case that one or both of these parameters can be of order 1, thus necessitating a reorganization of the expansion. Here we will not review the systematics of these expansions, as they have been thoroughly discussed in the literature .
The nature of the abovementioned series depends on the kinematic variable with which we choose to cut. Each choice has advantages and disadvantages. The cut lepton energy spectrum, $`E_l>(M_B^2M_D^2)/(2M_B)`$, is relatively simple experimentally since there is no need for neutrino reconstruction. Theoretically, this choice of cuts is challenging due to the breakdown of the non-perturbative series in $`\mathrm{\Lambda }/(m_b2E_l)`$. Furthermore, the fraction of the $`bu`$ transitions included after the cut is only about $`10\%`$. The cut hadronic invariant mass spectrum is more challenging from an experimental viewpoint, but has its virtues. First of all, while this observable smears over the same range of hadronic masses as the cut lepton spectrum, it is weighted more towards states of larger invariant mass. We thus would expect local duality to be more effective in this case. Furthermore, this cut rate includes around $`4080\%`$ of the $`bu`$ transitions . In both of these cases the non-perturbative series must be resummed. The necessity for resummation of the perturbative series can only be determined a posteriori. For example, in the case of radiative decay with an experimental cut of $`2.1`$ GeV on the photon energy, it was found that the resummation of the perturbative series was not necessary. In it was shown that, ignoring the non-perturbative structure function, the perturbative resummation is crucial, and that the cut hadronic rate can not be calculated in a systematic way. However, as was emphasized in , this conclusion may not stop us from extracting $`|V_{ub}|`$ since the inclusion of the structure function can greatly soften the effects of perturbative resummation.
In it was shown how the endpoint data from radiative decay can be used to eliminate our ignorance of the Fermi motion of the heavy quark (i.e., the non-perturbative structure function) in the electron end point spectrum. In this paper we will apply these ideas to the hadronic mass spectrum. The next section is a synopsis of the results given in for the resummed hadronic invariant mass spectrum, which is followed by a section dedicated to applying the formalism of to derive a closed form expression for $`|V_{ub}|`$ in terms of the hadronic mass spectrum. The final section of the paper includes the results and a discussion of the errors.
## II Infrared Factorization and Resummation
The factorization and resummation for inclusive $`B`$ decays in leptonic variables are discussed in Ref. . It was extended to hadronic variables in Ref. . Here we will briefly review the results. Consider the inclusive semi-leptonic decay of the $`b`$ quark into a lepton pair with momenta $`q=(p_l+p_\nu )`$ and a hadronic jet of momenta $`p_h`$. We define the following leptonic kinematic variables
$`y_0`$ $`=`$ $`{\displaystyle \frac{2vq}{m_b}},`$ (1)
$`y`$ $`=`$ $`{\displaystyle \frac{q^2}{m_b^2}},`$ (2)
$`x`$ $`=`$ $`{\displaystyle \frac{2vp_l}{m_b}},`$ (3)
where $`v=(1,\stackrel{}{0})`$, and partonic kinematic variables
$`\widehat{s}_0`$ $`=`$ $`{\displaystyle \frac{s_0}{m_b^2}}={\displaystyle \frac{p_h^2}{m_b^2}},`$ (4)
$`h`$ $`=`$ $`{\displaystyle \frac{2vp_h}{m_b}}.`$ (5)
In terms of the leptonic variables, $`\widehat{s}_0=(1y_0+y)`$ and $`h=2y_0`$, one can see that in the end point region of the electron energy spectrum when $`x1`$ with $`y<1`$, the invariant mass of the jet approaches zero with its energy held fixed. In addition, the jet hadronizes at a much later time in the rest frame of the $`B`$ meson, due to the time dilation. Factorization exploits this and separates the particular differential rate under consideration into subprocesses with disparate scales. This factorization fails when the jet energy vanishes in the dangerous region $`yx1`$. However, this problematic region of phase space is suppressed because the rate to produce soft massless fermions vanishes at tree level.
In terms of $`y_0`$ and $`y`$, the triply differential rate, which factorizes into hard, jet and soft subprocesses , may be written as
$$\frac{1}{\mathrm{\Gamma }_0}\frac{d^3\mathrm{\Gamma }}{dy_0dydx}=\mathrm{\hspace{0.33em}12}(y_0x)(xy)_\xi ^{x_m}𝑑zS(z)m_b^2J[m_b^2h(z\xi ),\mu ]H(m_bh/\mu ),$$
(6)
$$\mathrm{\Gamma }_0=\frac{G_F^2}{192\pi ^3}|V_{ub}|^2m_b^5,$$
(7)
where $`x_m=M_B/m_b`$, $`\mu `$ is a factorization scale, and $`\xi =(1y)/(2y_0)`$ is analogous to the Bjorken scaling variable in deep inelastic scattering. $`z=1+k_+/m_b`$, where $`k_+`$ is the heavy quark light cone residual momentum. $`S(z)`$ essentially describes the probability for the $`b`$ quark to carry light cone momentum fraction $`z`$ and allows for a leakage past the partonic endpoint, as can be seen explicitly in the upper limit of $`z`$.
It is convenient to change variables from $`y`$ to $`\xi `$ and perform the $`x`$ integration to yield
$$\frac{1}{\mathrm{\Gamma }_0}\frac{d^2\mathrm{\Gamma }}{dy_0d\xi }=2(2y_0)^3(2y_01)_\xi ^{x_m}𝑑zS(z)m_b^2J[m_b^2(2y_0)(z\xi )]H[m_b(2y_0)].$$
(8)
To proceed, we take the $`N`$th moment with respect to $`\xi `$ in the large $`N`$ limit. In the region $`\widehat{s}_00`$ and $`z\xi 1`$, one can replace $`J[m_b^2h(z\xi )]`$ in Eq. (8) with $`J[m_b^2h(1\xi /z)]`$. This replacement is permissible to the order we are working. We then obtain
$`M_N`$ $`=`$ $`{\displaystyle _0^{x_m}}𝑑\xi \xi ^{N1}{\displaystyle \frac{1}{\mathrm{\Gamma }_0}}{\displaystyle \frac{d^2\mathrm{\Gamma }}{d\xi dy_0}}`$ (9)
$`=`$ $`2(2y_0)^2(2y_01)S_NJ_N[m_b^2(2y_0)/\mu ^2]H[m_b(2y_0)/\mu ]+𝒪(1/N),`$ (10)
$`J_N(m_b^2/\mu ^2)`$ $`=`$ $`m_b^2{\displaystyle _0^1}𝑑yy^{N1}J[m_b^2(1y),\mu ],`$ (11)
$`S_N`$ $`=`$ $`{\displaystyle _0^{x_m}}𝑑zz^NS(z).`$ (12)
The soft moment $`S_N`$ further decomposes into a perturbative piece, which accounts for soft gluon radiation and a non-perturbative piece which incorporates bound state dynamics and serves as the boundary condition for the renormalization group equation .
$`S(z)`$ $`=`$ $`{\displaystyle _z^{x_m}}{\displaystyle \frac{ds}{s}}f[m_b(s1)]\sigma (z/s),`$ (13)
$`S_N`$ $`=`$ $`f_N\sigma _N,`$ (14)
where $`f(s)=B(v)|\overline{b}_v\delta (siD_+)b_v|B(v)`$ is the non-perturbative structure function defined in Ref. , which resums an infinite number of operators all in the same order of $`\mathrm{\Lambda }/m_b`$. A similar expression holds for the inclusive $`BX_s\gamma `$ in the end point region of the photon spectrum
$`M_N^\gamma `$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{\Gamma }_0^\gamma }}{\displaystyle _0^{x_m}}𝑑x_\gamma x_\gamma ^{N1}{\displaystyle \frac{d\mathrm{\Gamma }^\gamma }{dx_\gamma }}`$ (15)
$`=`$ $`S_NJ_N^\gamma H^\gamma ,`$ (16)
$`\mathrm{\Gamma }_0^\gamma `$ $`=`$ $`{\displaystyle \frac{G_F^2|V_{ts}^{}V_{tb}|^2\alpha C_7^2m_b^5}{32\pi ^4}}.`$ (17)
Here $`C_7`$ is the Wilson coefficient of the $`O_7`$ operator evaluated at the scale $`m_b`$ and $`x_\gamma =2E_\gamma /m_b`$.
In moment space the soft and jet functions have been calculated to next-to-leading (NLL) logarithmic order and are given by
$`\sigma _NJ_N`$ $`=`$ $`\mathrm{exp}[\mathrm{log}(N)g_1(\chi )+g_2^\gamma (\chi )+g_{sl}(\chi ,y_0)]`$ (18)
$`\sigma _NJ_N^\gamma `$ $`=`$ $`\mathrm{exp}[\mathrm{log}(N)g_1(\chi )+g_2^\gamma (\chi )],`$ (19)
where $`\chi =\alpha _s(m_b^2)\beta _0\mathrm{log}N`$, and $`g_1`$, $`g_2^\gamma `$ and $`g_{sl}`$ can be found in . The hard parts are given in :
$`H(y_0)`$ $`=`$ $`1{\displaystyle \frac{2\alpha _s}{3\pi }}\left[4\mathrm{log}^2(2y_0)+{\displaystyle \frac{810y_0}{2y_01}}\mathrm{log}(2y_0)+2\mathrm{L}\mathrm{i}_2(y_01)+{\displaystyle \frac{5}{2}}+{\displaystyle \frac{2\pi ^2}{3}}\right],`$ (20)
$`H^\gamma `$ $`=`$ $`1{\displaystyle \frac{2\alpha _s}{3\pi }}\left({\displaystyle \frac{13}{2}}+{\displaystyle \frac{2\pi ^2}{3}}\right).`$ (21)
To get back the physical spectra from the moment space, the inverse Mellin transform has to be evaluated at NLL accuracy as well. If we are willing to ignore the effects of the structure function we can get a closed form expression for the resummed rate. To this end, we use the results derived in the Appendix of Ref.
$`{\displaystyle \frac{1}{\mathrm{\Gamma }_0}}{\displaystyle \frac{d^2\mathrm{\Gamma }}{dy_0d\xi }}`$ $`=`$ $`2(2y_0)^2(2y_01)H(y_0)`$ (23)
$`\times {\displaystyle \frac{d}{d\xi }}\left\{\theta \left(1\xi \eta \right){\displaystyle \frac{e^{lg_1(\alpha _s\beta _0l)+g_2(\alpha _s\beta _0l)}}{\mathrm{\Gamma }\left[1g_1(\alpha _s\beta _0l)\alpha _s\beta _0lg_1^{}(\alpha _s\beta _0l)\right]}}\right\},`$
where $`l=\mathrm{log}\left(\mathrm{log}\xi \right)\mathrm{log}(1\xi )`$ and $`g_2=g_2^\gamma +g_{sl}`$. The $`\theta `$-function defines the differential rates in a distribution sense, as $`\eta 0`$, and turn the singular terms into the ’$`+`$’ distributions.
This result would only be useful in making a physical prediction if we had a measured structure function with which we could perform a convolution. Instead, following we will take the inverse Mellin transform of the ratio between the moments of the semi-leptonic and radiative decay rates.
## III Extraction of $`|V_{ub}|`$
The key observation to extract $`|V_{ub}|`$ is that the soft function, Eq. (13), is universal in inclusive $`B`$ decays. Using (15) – (19) we can rewrite (9) as
$$M_N=2(2y_0)^2(2y_01)\frac{H(y_0)}{H^\gamma }e^{g_{sl}(\chi ,y_0)}M_N^\gamma $$
(24)
Taking the inverse Mellin transform we arrive at
$`{\displaystyle \frac{\mathrm{\Gamma }_0^\gamma }{\mathrm{\Gamma }_0}}{\displaystyle \frac{d^2\mathrm{\Gamma }}{dy_0d\xi }}`$ $`=`$ $`2(2y_0)^2(2y_01)`$ (26)
$`\times {\displaystyle _\xi ^1}\stackrel{~}{H}(y_0)du{\displaystyle \frac{d\mathrm{\Gamma }^\gamma }{du}}(z{\displaystyle \frac{d}{dz}})\left[\theta (1z\eta )e^{g_{sl}(\chi _z,y_0)}\right],`$
where
$`\stackrel{~}{H}(y_0)`$ $`=`$ $`{\displaystyle \frac{H(y_0)}{H^\gamma }}`$ (27)
$`=`$ $`1{\displaystyle \frac{2\alpha _s}{3\pi }}\left[4\mathrm{log}^2(2y_0)+{\displaystyle \frac{810y_0}{2y_01}}\mathrm{log}(2y_0)+2\mathrm{L}\mathrm{i}_2(y_01)4\right]`$ (28)
$`g_{sl}(\chi ,y_0)`$ $`=`$ $`{\displaystyle \frac{4}{3\pi \beta _0}}\mathrm{log}(2y_0)\mathrm{log}(1\chi ),`$ (29)
$`\chi _z=\alpha _s\beta _0\mathrm{log}(1z)`$ and $`z=\xi /u`$. We have now rescaled to “physical” variables which are normed to the meson mass $`M_B`$. We have also dropped a factor of $`\overline{\mathrm{\Lambda }}/m_b`$ which would multiply $`y_0`$ in these renormed variables, as it would only be relevant in the region where factorization breaks down, $`y_02`$, which, as previously mentioned is the suppressed part of phase space.
Now to get the cut integrated rate we change the order of integration so that integral over photon energy is the last to be performed. This leads to a result which can most easily be integrated over the data for the radiative decay. We find the following expression for the rate with a cut $`c=s_c/M_B^2`$ on the hadronic invariant mass spectrum:
$`\delta \mathrm{\Gamma }(c)`$ $`=`$ $`({\displaystyle _{1\frac{\sqrt{c}}{2}}^{\frac{1}{1+c}}}du{\displaystyle _{\frac{1\sqrt{c}/2}{u}}^1}dz{\displaystyle _{2\frac{c}{1uz}}^{4uz2}}dy_0+{\displaystyle _{\frac{1}{1+c}}^1}du{\displaystyle _{\frac{1\sqrt{c}/2}{u}}^{\frac{1}{u(1+c)}}}dz{\displaystyle _{2\frac{c}{1uz}}^{4uz2}}dy_0`$ (31)
$`+{\displaystyle _{\frac{1}{1+c}}^1}du{\displaystyle _{\frac{1}{u(1+c)}}^1}dz{\displaystyle _{2\frac{1}{uz}}^{4uz2}}dy_0){\displaystyle \frac{\mathrm{\Gamma }_0}{\mathrm{\Gamma }_0^\gamma }}I(u,y_0,z)`$
$`I(u,y_0,z)`$ $`=`$ $`2(2y_0)^2(2y_01)\stackrel{~}{H}(y_0){\displaystyle \frac{d\mathrm{\Gamma }^\gamma }{du}}\left(z{\displaystyle \frac{d}{dz}}\right)\left[\theta (1z\eta )e^{g_{sl}(\chi _z,y_0)}\right].`$ (32)
We can estimate the effects of resummation on the cut rate by using a model for the structure function. In Figure 1, we show the effects of the resummation of the end point logs as a function of the cut using the same model as in . We see that the effects of resummation are small, and thus the end point logs do not form a dominant sub-series, and the unknown piece of the two loop result dominates the perturbative uncertainties. Note that this is not inconsistent with the results of , where it was found that the effects of resummation were large, as the structure function mollifies the effects of the resummation.
If we now expand $`g_{sl}`$ to order $`\alpha _s`$, we can write down a simple expression for $`|V_{ub}|`$ as a convolution of the data from radiative decay. We may use the approximation
$$\frac{d}{dz}\left(\theta (1z)e^{g_{sl}}\right)\delta (1z)+\frac{4\alpha _s}{3\pi }\left(\frac{\mathrm{log}(2y_0)}{1z}\right)_+.$$
(33)
We then find
$$\delta \mathrm{\Gamma }(c)=\frac{\mathrm{\Gamma }_0}{\mathrm{\Gamma }_0^\gamma }[I_0(c)+I_+(c)],$$
(34)
$`I_0(c)`$ $`=`$ $`\left({\displaystyle _{1\frac{\sqrt{c}}{2}}^{\frac{1}{1+c}}}𝑑u{\displaystyle _{2\frac{c}{1u}}^{4u2}}𝑑y_0+{\displaystyle _{\frac{1}{1+c}}^1}𝑑u{\displaystyle _{2\frac{1}{u}}^{4u2}}𝑑y_0\right)2(2y_0)^2(2y_01){\displaystyle \frac{d\mathrm{\Gamma }^\gamma }{du}}\stackrel{~}{H}(y_0),`$ (35)
$`I_+(c)`$ $`=`$ $`{\displaystyle 𝑑u𝑑z𝑑y_0\mathrm{\hspace{0.17em}2}(2y_0)^2(2y_01)z\frac{d\mathrm{\Gamma }^\gamma }{du}\frac{4\alpha _s}{3\pi }\left(\frac{\mathrm{log}(2y_0)}{1z}\right)_+}.`$ (36)
The plus distribution is defined as
$$_a^1𝑑x\frac{f(x)}{(1x)_+}=_a^1𝑑x\frac{f(x)f(1)}{1x}_0^a𝑑x\frac{f(1)}{1x}.$$
(37)
The regions of integration in Eq. (36) are the same as those in Eq. (31). The integrals over $`y`$ and $`z`$ can all be done analytically. The final result is rather large and complicated, so we choose not to include it in print. The resulting expression is an integral which involves a known function of $`u`$ and $`c`$ and the radiative decay endpoint data. Thus we may write
$$\frac{|V_{ub}|}{|V_{ts}|}=\left\{\frac{6\alpha C_7(m_b)^2\delta \mathrm{\Gamma }(c)}{\pi [I_0(c)+I_+(c)]}\right\}^{\frac{1}{2}}.$$
(38)
The dominant source of error, modulo the usual uncertainties inherent in all inclusive predictions, will come from unknown effects at order $`\alpha _s^2`$, and $`\mathrm{\Lambda }^2/(cm_b^2)`$, each of which should contribute only at the few percent level.
## IV Results, Conclusions and Outlook
The extraction of $`|V_{ub}|`$ using the cut hadronic mass spectrum has the theoretical advantage over the lepton spectrum since the mass cut includes a larger fraction of the $`bu`$ transition. In addition, we expect local parton-hadron duality to work better since more resonances will contribute to the cut mass spectrum. The real issue for the viability of this method is the resolution. In particular, how good will the “best” resolution be for the invariant mass? This will be limited by the inability to detect neutrals, as well as the charge particle detection inefficiencies. Presently, the resolution is not good enough to eliminate the charmed background, which has to be modeled. Hopefully, it will be possible to eventually push the cut further. If the resolution does become good enough to really eliminate the charmed background, then there is an additional option available. It has been pointed out that the leptonic mass spectrum is effectively less sensitive to the Fermi motion of the $`b`$ quark . Recently, this fact has been used to determine an expression for $`|V_{ub}|`$ in terms of the cut lepton mass spectrum, without needing to use the data from radiative decay. However, the resulting expansion parameter is $`\mathrm{\Lambda }/M_D`$, and grows as the cut is raised above $`(M_BM_D)^2`$. The authors in calculated the leading non-perturbative correction and found it to be anomalously small. They made an estimate for the sub-leading corrections and found that they grow rapidly as the cut is raised from its smallest possible value.
In the end it is clear that no one single extraction of $`|V_{ub}|`$ should be trusted, given the fact that our expansion parameters are never as small as we would like. We will only gain confidence in the extractions after there is convergence among several independent extractions.
###### Acknowledgements.
This work was supported in part by the Department of Energy under grant numbers DOE-ER-40682-143 and DE-FG03-92-ER 40701, and also by the NSF under grant number PHY94-07194. The authors would like to thank the theory group at Caltech for their hospitality. I. L. is supported in part by an ITP Graduate Fellowship.
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# 1 Introduction, SCRMs and Lovász-Schrijver procedures
## 1 Introduction, SCRMs and Lovász-Schrijver procedures
Since 1960’s, complementarity problems attracted a very significant attention in the theory as well as applications of operations research. See, for instance, the book on LCP . In this paper, we consider various complementarity problems in the context of Successive Convex Relaxation Methods (SCRMs) proposed by the authors . Since these methods can be used to compute the convex hull of any compact subset of an Euclidean space described by a system of quadratic inequalities and a compact convex set, they can be used to attack many complementarity problems from several angles.
In the special case of 0-1 optimization problems over convex sets, or more specially polytopes, there are many Successive Convex Relaxation Methods (SCRMs) based on lift-and-project techniques. We also discuss some of the relationships of general SCRMs and these more specialized algorithms in solving LCPs.
Let $`F`$ be a compact set in the $`n`$-dimensional Euclidean space $`R^n`$. SCRMs take as input, a compact convex subset $`C_0`$ of $`R^n`$ and a set $`𝒫_F`$ of quadratic functions which induce a description of $`F`$ such that
$$F=\{𝒙C_0:qf(𝒙;\gamma ,𝒒,𝑸)0,qf(;\gamma ,𝒒,𝑸)𝒫_F\}.$$
Here we denote by $`qf(;\gamma ,𝒒,𝑸)`$, the quadratic function $`(\gamma +2𝒒^T𝒙+𝒙^T𝑸𝒙)`$. Note that the variable $`𝒙`$ is irrelevant outside a context and it will always be clear what the variable vector is, from the context.
Let $`\mathrm{}`$ be an integer such that $`1<2\mathrm{}m`$, $`𝒅R^m`$, and let $`A`$ be a compact convex subset of $`R^m`$. Consider the convex optimization problem with complementarity conditions:
$$\begin{array}{cc}\text{maximize}\hfill & 𝒅^T𝒖\hfill \\ \text{subject to}\hfill & 𝒖A,0u_i,0u_{i+\mathrm{}},u_iu_{i+\mathrm{}}=0,i\{1,2,\mathrm{},\mathrm{}\}.\hfill \end{array}\}$$
(1)
First of all, it is clear that LCP, with a known upper bound on a solution of it, is a special case of (1) (we can take $`m=2\mathrm{}`$ and $`A`$ as an affine subspace intersected with a large enough ball). Secondly, it is very elementary to formulate this problem as a mixed 0-1 optimization problem with convex constraints:
$$\begin{array}{ccc}\text{maximize}\hfill & 𝒄^T𝒗\hfill & \\ \text{subject to}\hfill & 𝒗C_0,v_i\{0,1\},i\{m+1,m+2,\mathrm{},n\},\hfill & \end{array}\}$$
(2)
where
$`C_0`$ $``$ $`\{𝒗=\left(\begin{array}{c}𝒖\\ v_{m+1}\\ \mathrm{}\\ v_n\end{array}\right)R^{m+\mathrm{}}:\begin{array}{ccc}𝒖A,\hfill & & \\ 0u_irv_{m+i},\hfill & & \\ 0u_{i+\mathrm{}}r(1v_{m+i}),\hfill & & \\ i\{1,2,\mathrm{},\mathrm{}\}\hfill & & \end{array}\},𝒄\left(\begin{array}{c}𝒅\\ \text{0}\end{array}\right)R^{m+\mathrm{}},`$
$`n`$ $``$ $`m+\mathrm{},r\underset{i}{\mathrm{max}}\left\{\mathrm{max}\{u_i:𝒖A\}\right\}.`$
In general, we allow $`C_0`$ to be an arbitrary compact convex set in $`R^n`$. There are various successive convex relaxation methods that can be applied to such a problem.
We can represent the feasible region $`FR^n`$ of (2) as
$$F=\{𝒗C_0:p(𝒗)0,p()𝒫_F\},$$
where $`𝒫_F`$ denotes a set consisting of quadratic functions
$$(v_i^2v_i),(v_i^2+v_i),i\{m+1,m+2,\mathrm{},n\}$$
on $`R^n`$.
In connection with the SCRMs and also the Lovász-Schrijver procedures (see ), it seems convenient to introduce the following notation: For every compact convex relaxation $`CC_0`$ of $`F`$ and every subset $`D`$ of $`\overline{D}\{𝒅R^n:𝒅=1\}`$,
$`𝒫^2(C,D)`$ $``$ $`\{(𝒅^T𝒗\alpha (C_0,𝒅))(\overline{𝒅}^T𝒗\alpha (C,\overline{𝒅})):𝒅D,\overline{𝒅}\overline{D}\},`$
$`\widehat{𝒩}(C,D)`$ $``$ $`\{𝒗C_0:\begin{array}{c}𝑽𝒮^n\text{such that }\hfill \\ \gamma +2𝒒^T𝒗+𝑸𝑽0,\hfill \\ qf(;\gamma ,𝒒,𝑸)𝒫_F𝒫^2(C,D)\hfill \end{array}\}`$
$`\text{(a Semi-Infinite LP relaxation of }F\text{)},`$
$`\widehat{𝒩}_+(C,D)`$ $``$ $`\{𝒗C_0:\begin{array}{c}𝑽𝒮^n\text{such that }\left(\begin{array}{cc}1\hfill & 𝒗^T\hfill \\ 𝒗\hfill & 𝑽\hfill \end{array}\right)𝒮_+^{1+n},\hfill \\ \gamma +2𝒒^T𝒗+𝑸𝑽0,\hfill \\ qf(;\gamma ,𝒒,𝑸)𝒫_F𝒫^2(C,D)\hfill \end{array}\}`$
$`\text{(an SDP relaxation of }F\text{)},`$
where $`\alpha (C,𝒅)\mathrm{max}\{𝒅^T𝒗:𝒗C\}`$ for every $`𝒅\overline{D}`$. Let $`𝒮^n`$ and $`𝒮_+^{1+n}`$ denote the set of $`n\times n`$ symmetric matrices and the set of $`(1+n)\times (1+n)`$ symmetric positive semidefinite matrices, respectively. The corresponding variants of Successive Semi-Infinite LP Relaxation Method (SSILPRM) and Successive SDP Relaxation Method (SSDPRM) can be written as follows.
###### Algorithm 1.1
(SSILPRM)
Step 0: Choose a $`D_0\overline{D}`$. Let $`k=0`$.
Step 1: If $`C_k=`$ (the convex hull of $`F`$), then stop.
Step 2: Let $`C_{k+1}=\widehat{𝒩}(C_k,D_0)`$.
Step 3: Let $`k=k+1`$, and go to Step 1.
###### Algorithm 1.2
(SSDPRM)
Steps 0, 1 and 3: The same as the Steps 0, 1 and 3 of Algorithm 1.1.
Step 2: Let $`C_{k+1}=\widehat{𝒩}_+(C_k,D_0)`$.
To connect these algorithms to the Lovász-Schrijver procedures, we need to introduce some additional notation. For every pair of closed convex cones $`𝒦`$ and $`𝒯`$ in $`R^{1+n}`$, define
$`(𝒦,𝒯)`$ $``$ $`\{𝒀=\left(\begin{array}{cc}\lambda \hfill & \lambda 𝒗^T\hfill \\ \lambda 𝒗\hfill & \lambda 𝑽\hfill \end{array}\right):\begin{array}{c}\lambda 0,𝒗C_0,𝑽𝒮^n,\hfill \\ v_i=V_{ii},i\{m+1,m+2,\mathrm{},n\},\hfill \\ 𝒗^T𝒀𝒘0,𝒗𝒯^{},𝒘𝒦^{}\hfill \end{array}\},`$
$`_+(𝒦,𝒯)`$ $``$ $`\{𝒀=\left(\begin{array}{cc}\lambda \hfill & \lambda 𝒗^T\hfill \\ \lambda 𝒗\hfill & \lambda 𝑽\hfill \end{array}\right):\begin{array}{c}\lambda 0,𝒗C_0,𝑽𝒮^n,𝒀𝒮_+^{1+n}\hfill \\ v_i=V_{ii},i\{m+1,m+2,\mathrm{},n\},\hfill \\ 𝒗^T𝒀𝒘0,𝒗𝒯^{},𝒘𝒦^{}\hfill \end{array}\}.`$
Let $`𝒯_0`$ and $`𝒦_0`$ be closed convex cones given by
$`𝒯_0^{}`$ $`=`$ $`\text{c.cone}\left(\{\left(\begin{array}{c}\alpha (C_0,𝒅)\\ 𝒅\end{array}\right)R^{1+n}:𝒅D_0\}\right),`$
$`𝒦_0`$ $`=`$ $`\{\left(\begin{array}{c}\lambda \\ \lambda 𝒗\end{array}\right)R^{1+n}:𝒗C_0,\lambda 0\}.`$
(Note that $`𝒯_0`$ itself is defined as the dual of $`𝒯_0^{}`$.) If $`CC_0`$ is a compact convex relaxation of $`F`$ and
$$𝒦=\{\left(\begin{array}{c}\lambda \\ \lambda 𝒗\end{array}\right)R^{1+m}:𝒗C,\lambda 0\},$$
then
$`\widehat{𝒩}(C,D_0)`$ $`=`$ $`\{𝒗R^n:\left(\begin{array}{cc}1\hfill & 𝒗^T\hfill \\ 𝒗\hfill & 𝑽\hfill \end{array}\right)(𝒦,𝒯_0)\},`$
$`\widehat{𝒩}_+(C,D_0)`$ $`=`$ $`\{𝒗R^n:\left(\begin{array}{cc}1\hfill & 𝒗^T\hfill \\ 𝒗\hfill & 𝑽\hfill \end{array}\right)_+(𝒦,𝒯_0)\}.`$
Algorithms 1.1 and 1.2 specialized to (2) with $`𝒫_F=\{v_i^2v_i,v_i^2+v_i,i\{m+1,m+2,\mathrm{},n\}\}`$ can be stated in the following forms, which are essentially the Lovász-Schrijver procedures.
Algorithm 1.1H (Homogeneous form of Algorithm 1.1)
Step 0: Choose a $`D_0\overline{D}`$. Define $`𝒯_0`$ and $`𝒦_0`$ as above. Let $`k=0`$.
Step 1: If $`𝒦_k=\text{c.cone}\left(\{\left(\begin{array}{c}1\\ 𝒗\end{array}\right):𝒗F\}\right)`$ then stop.
Step 2: Let $`𝒦_{k+1}=\{𝒀𝒆_0:𝒀(𝒦_k,𝒯_0)\}`$.
Step 3: Let $`k=k+1`$, and go to Step 1.
Algorithm 1.2H (Homogeneous form of Algorithm 1.2)
Steps 0, 1 and 3: The same as Steps 0, 1 and 3 of Algorithm 1.2H, respectively.
Step 2: Let $`𝒦_{k+1}=\{𝒀𝒆_0:𝒀_+(𝒦_k,𝒯_0)\}`$.
In this paper $`𝒆_j`$ denotes the $`j`$th unit vector and $`𝒆`$ denotes the vector of all ones (the dimensions of the vectors will be clear from the context).
## 2 SCRMs applied to LCP with an á priori bound
Let $`𝑴R^\mathrm{}\times \mathrm{}`$, $`𝒒R^{\mathrm{}}`$ be given. Consider the LCP in the following form.
(LCP) Find $`𝒙`$, $`𝒔`$ such that $`𝑴𝒙+𝒒=𝒔,`$
$`𝒙\text{0},𝒔\text{0},`$
$`x_is_i=\mathrm{\hspace{0.17em}\hspace{0.17em}0},i\{1,2,\mathrm{},\mathrm{}\}.`$
Suppose we are given $`(𝝃,r)\{𝒖R^2\mathrm{}:𝒖𝝃r\}`$, an Euclidean ball containing a solution of the LCP. (In the case of rational data $`(𝑴,𝒒)`$, we can take $``$ centered at the origin with the radius bounded above by a polynomial function of the “bit size” of the data $`(𝑴,𝒒)`$.) For the rest of this section, we assume that the Euclidean ball with center $`𝝃\text{0}`$ and the radius $`r`$ ($`r`$ is assumed given) contains some solution of the LCP.
Under the boundedness assumption above, it is particularly easy to model any LCP as a 0-1 mixed integer programming problem, since the only nonlinear constraints of LCP can be expressed as
$$x_i=\mathrm{\hspace{0.17em}\hspace{0.17em}0},\text{ or }s_i=\mathrm{\hspace{0.17em}\hspace{0.17em}0},i\{1,2,\mathrm{},\mathrm{}\}.$$
Balas’ method can be directly applied to such formulations. We can also apply some variants of the Lovász-Schrijver procedures to the mixed integer programming feasibility problem:
Find $`𝒙`$, $`𝒔`$ and $`𝒛`$ such that $`𝑴𝒙+𝒒=𝒔,`$
$`\text{0}𝒙r𝒛,\text{0}𝒔r(𝒆𝒛),`$
$`𝒛\{0,1\}^{\mathrm{}}.`$
Note that we can eliminate the variable vector $`𝒔`$ from the formulation and apply the SSILPR and SSDPR Methods to the following formulation:
$$\begin{array}{c}\text{0}𝑴𝒙+𝒒r(𝒆𝒛),\hfill \\ \text{0}𝒛𝒆,\text{0}𝒙r𝒛,\hfill \\ z_i^2z_i\mathrm{\hspace{0.17em}\hspace{0.17em}0},z_i^2+z_i\mathrm{\hspace{0.17em}\hspace{0.17em}0},i\{1,2,\mathrm{},\mathrm{}\}.\hfill \end{array}$$
To apply the SCRMs, we can take
$`C_0\{𝒗=\left(\begin{array}{c}𝒙\\ 𝒛\end{array}\right)R^n:\begin{array}{c}\text{0}𝑴𝒙+𝒒r(𝒆𝒛),\hfill \\ \text{0}𝒛𝒆,\text{0}𝒙r𝒛\hfill \end{array}\},`$
$`m\mathrm{},n2\mathrm{},`$
$`𝒫_F\left\{(v_i^2v_i),(v_i^2+v_i),i\{m+1,m+2,\mathrm{},n\}\right\}.`$
Both algorithms, SSILPRM and SSDPRM presented in Section 1, terminate in at most $`\mathrm{}`$ steps. This fact can be proved easily, using the results of Balas , Sherali and Adams , Lovász and Schrijver , or Kojima and Tunçel . For computational experience on similar algorithms for similar problems see , . In the next section, we give the details of a proof of such a convergence result when the methods are applied to a formulation of Pardalos and Rosen .
## 3 SCRMs applied to Pardalos-Rosen formulation of LCP
We will illustrate the convergence proof on a formulation of (LCP) by Pardalos and Rosen . They homogenize the vector $`𝒒`$ with a new continuous variable $`\alpha `$, then they maximize $`\alpha `$.
$`\begin{array}{ccc}(\text{MIP}_\alpha )\hfill & \text{maximize}\hfill & \alpha \hfill \\ & \text{subject to}\hfill & \text{0}𝑴𝒙+𝒒\alpha 𝒆𝒛,\hfill \\ & & \text{0}𝒙𝒛,0\alpha \mathrm{\hspace{0.17em}\hspace{0.17em}1},𝒛\{0,1\}^{\mathrm{}}.\hfill \end{array}\}`$
Note that $`\left(\begin{array}{c}\overline{\alpha }\\ \overline{𝒙}\\ \overline{𝒛}\end{array}\right)\text{0}`$ is feasible in (MIP<sub>α</sub>) and, it is easy to see that (MIP<sub>α</sub>) has an optimal solution with $`\alpha ^{}>0`$ iff the (LCP) has a solution (or solutions) . Moreover, if $`\left(\begin{array}{c}\alpha ^{}\\ 𝒙^{}\\ 𝒛^{}\end{array}\right)`$ is an optimal solution of (MIP<sub>α</sub>) with $`\alpha ^{}>0`$ then $`{\displaystyle \frac{𝒙^{}}{\alpha ^{}}}`$ solves the (LCP) . One advantage of (MIP<sub>α</sub>) is that it does not require the introduction of large, data dependent constants (such as $`r`$ in the previous section) or their a priori estimates. Now, we take
$`C_0\{𝒗=\left(\begin{array}{c}\alpha \\ 𝒙\\ 𝒛\end{array}\right)R^{1+2\mathrm{}}:\begin{array}{c}\text{0}𝑴𝒙+𝒒\alpha 𝒆𝒛,\hfill \\ \text{0}𝒙𝒛,0\alpha \mathrm{\hspace{0.17em}\hspace{0.17em}1}\hfill \end{array}\},`$
$`m\mathrm{}+1,n2\mathrm{}+1,`$
$`𝒫_F\left\{(v_i^2v_i),(v_i^2+v_i),i\{m+1,m+2,\mathrm{},n\}\right\}.`$
We have an analog of a very elementary but also a key lemma (Lemma 1.3 of ) of Lovász and Schrijver (and their proof technique is adapted here). In what follows, we refer to the vectors in the space of $`𝒦_k`$ by $`𝒗`$. At the same time, we refer to different subvectors of $`𝒗`$ by different names, such as $`𝒙`$, $`\alpha `$ etc., to keep the correspondence of elements of $`𝒗`$ and the original formulation of $`F`$ clearer. The proof of Lemma 1.3 of leads to the following analogous result in our case.
###### Lemma 3.1
Let $`D_0\{\pm 𝐞_{m+1},\pm 𝐞_{m+2},\mathrm{},\pm 𝐞_n\}.`$ Then the sequence of convex cones $`\{𝒦_k:k0\}`$ given by Algorithm 1.1H satisfies
$$𝒦_{k+1}\left(𝒦_k\{𝒗:x_i=0\}\right)+\left(𝒦_k\{𝒗:(𝐌𝐱+𝒒\alpha )_i=0\}\right),$$
for every $`i\{1,2,\mathrm{},\mathrm{}\},`$ and for every $`k0.`$
Proof: Let $`𝒘\left(\begin{array}{c}1\\ \overline{\alpha }\\ \overline{𝒙}\\ \overline{𝒛}\end{array}\right)𝒦_{k+1}`$. Fix $`j\{1,2,\mathrm{},\mathrm{}\}`$ arbitrarily. By the definition of $`D_0`$ and $`𝒯_0`$, the unit vector $`𝒆_0`$ is in $`𝒯_0`$. Hence, by the definition of $`(𝒦_k,𝒯_0),`$ $`𝒦_{k+1}𝒦_k`$ for every $`k0`$. Therefore, $`𝒘𝒦_k`$. If $`\overline{x}_j=0`$ or $`(𝑴𝒙+𝒒\alpha )_j=0`$ then the statement of the lemma clearly holds. So, without loss of generality, we assume $`\overline{x}_j>0`$ and $`(𝑴𝒙+𝒒\alpha )_j>0`$. Let $`𝒀(𝒦_k,𝒯_0)`$ such that $`𝒘=𝒀𝒆_0`$. By our choice of the cone $`𝒯_0`$, we conclude that $`𝒀𝒆_{n+j}`$ and $`𝒀(𝒆_0𝒆_{n+j})`$ are both in $`𝒦_k`$. Note that
$$𝒘=\widehat{𝒘}+\stackrel{~}{𝒘},$$
where $`\widehat{𝒘}𝒀𝒆_{n+j}`$ and $`\stackrel{~}{𝒘}𝒀(𝒆_0𝒆_{n+j})`$. We will refer to the $`𝒙`$ and $`𝒛`$ parts of the vector $`\widehat{𝒘}`$ by $`\widehat{𝒙}`$, $`\widehat{𝒛}`$ etc. (Similarly for $`\stackrel{~}{𝒘}.`$) First, since by the definition of $`(𝒦_k,𝒯_0)`$, $`v_i=V_{ii}`$ for every $`i\{m+1,m+2,\mathrm{},n\}`$, we have $`\stackrel{~}{z}_j=0`$ which implies $`\stackrel{~}{x}_j=0`$. Therefore, $`\stackrel{~}{𝒘}`$ lies in the cone $`\left(𝒦_0\{𝒗:x_j=0\}\right)`$. Second, since $`\overline{x}_j>0`$, $`\overline{z}_j`$ must be positive. Therefore, $`(1/\overline{z}_j)\widehat{𝒘}𝒦_0`$. Since $`v_i=V_{ii}`$ for every $`i\{m+1,m+2,\mathrm{},n\}`$, $`\widehat{z}_j=\overline{z}_j`$. So,
$$\frac{1}{\overline{z}_j}\left(\begin{array}{c}\widehat{\alpha }\\ \widehat{𝒙}\\ \widehat{𝒛}\end{array}\right)C_k,$$
with its $`z_j`$ entry equal to 1. Thus, $`(𝑴\widehat{𝒙}+𝒒\widehat{\alpha })_j=0`$. Hence, $`\widehat{𝒘}`$ is in the cone
$`\left(𝒦_k\{𝒗:(𝑴𝒙+𝒒\alpha )_j=0\}\right)`$. Since the argument above is independent of the index $`j`$ the proof is complete.
Note that the conclusion of the above lemma also applies to the SSDPR Method since SSDPR Method yields at least as tight relaxations as the SSILPR Method.
###### Theorem 3.2
Both algorithms, Algorithm 1.1H and 1.2H terminate in $`\mathrm{}`$ iterations when applied to the formulation (MIP<sub>α</sub>) with our choice of $`𝒫_F`$, $`C_0`$ and $`D_0`$ above.
Proof: First note that
$$\text{c.hull}(F)\{\left(\begin{array}{c}\alpha \\ 𝒙\\ 𝒛\end{array}\right)R^n:\left(\begin{array}{c}1\\ \alpha \\ 𝒙\\ 𝒛\end{array}\right)𝒦_k\},k0.$$
Next, let $`i,j\{1,2,\mathrm{},\mathrm{}\}`$, $`ij`$. Since $`𝒙\text{0}`$ and $`𝑴𝒙+𝒒\alpha \text{0}`$, for all $`𝒗𝒦_k`$, for every $`k0`$,
$`\left[\left(𝒦_k\{𝒗:x_i=0\}\right)+\left(𝒦_k\{𝒗:(𝑴𝒙+𝒒\alpha )_i=0\}\right)\right]\{𝒗:x_j=0\}`$
$`=`$ $`\left(𝒦_k\{𝒗:x_i=0,x_j=0\}\right)+\left(𝒦_k\{𝒗:x_j=0,(𝑴𝒙+𝒒\alpha )_i=0\}\right).`$
Similarly, for the intersection with $`\{𝒗:(𝑴𝒙+𝒒\alpha )_j=0\}.`$ Now, we apply Lemma 3.1 repeatedly to conclude that $`𝒦_{\mathrm{}}`$ is the homogenization of the convex hull of all solutions of the LCP that lie in the original relaxation $`C_0`$.
## 4 SCRMs applied to the smaller formulation of LCP with explicit treatment of disjunctive constraints
Now, we consider a formulation with fewer variables and constraints.
$`\begin{array}{ccc}(\text{LCP}_\alpha )\hfill & \text{maximize}\hfill & \alpha \hfill \\ & \text{subject to}\hfill & 𝑴𝒙+𝒒\alpha \text{0},𝒙\text{0},\alpha \mathrm{\hspace{0.17em}\hspace{0.17em}0},\hfill \\ & & 𝒆^T(𝑴+𝑰)𝒙+(𝒆^T𝒒+1)\alpha \mathrm{\hspace{0.17em}\hspace{0.17em}1},\hfill \\ & & x_i(𝑴𝒙+𝒒\alpha )_i=0,i\{1,2,\mathrm{},\mathrm{}\}.\hfill \end{array}`$
It is easy to see that $`\left(\begin{array}{c}\overline{𝒙}\\ \overline{\alpha }\end{array}\right)\text{0}`$ is feasible in (LCP<sub>α</sub>), and it is also easy to observe that (LCP<sub>α</sub>) has an optimal solution with $`\alpha ^{}>0`$ iff the (LCP) has a solution(or solutions). Moreover, if $`\left(\begin{array}{c}𝒙^{}\\ \alpha ^{}\end{array}\right)`$ is an optimal solution of (LCP<sub>α</sub>) with $`\alpha ^{}>0`$ then $`{\displaystyle \frac{𝒙^{}}{\alpha ^{}}}`$ solves the (LCP). Note that the inclusion in Lemma 3.1 can sometimes be strict for the SSILPR and SSDPR Methods.
We explicitly include the variable vector $`𝒔`$ in our discussion in this section, for the sake of presentation. Let
$$C_0\{𝒗=\left(\begin{array}{c}𝒙\\ 𝒔\\ \alpha \end{array}\right)R^{2\mathrm{}+1}:\begin{array}{c}𝒔=𝑴𝒙+𝒒\alpha \text{0},𝒙\text{0},\alpha \mathrm{\hspace{0.17em}\hspace{0.17em}0},\hfill \\ 𝒆^T(𝑴+𝑰)𝒙+(𝒆^T𝒒+1)\alpha \mathrm{\hspace{0.17em}\hspace{0.17em}1},\hfill \end{array}\}.$$
In this section, we will describe another Successive Convex Relaxation Method based on the ideas of Balas , Lovász and Schrijver . This method will use only Linear Programming (LP) relaxations. We describe the method in the original space of $`F`$ and $`C_0`$. Let $`(C_0)`$ denote the set of facet defining inequalities for $`C_0`$. $`(C_0)`$ is the input of the algorithm which we introduce now.
###### Algorithm 4.1
Step 0. $`k0.`$
Step 1. $`(C_{k+1})(C_k).`$
Step 2. For every inequality
$$\underset{i=1}{\overset{\mathrm{}}{}}\left(u_ix_i+u_{\mathrm{}+i}s_i\right)u_{2\mathrm{}+1}\alpha u_0$$
in $`(C_k)`$ and every $`j\{1,2,\mathrm{},\mathrm{}\}`$ solve the LP problems
$`\begin{array}{ccc}(P_j)\hfill & \text{minimize}\hfill & 𝒖^T𝝃^{(j)}\hfill \\ & \text{subject to}\hfill & \xi _j^{(j)}=1,\xi _{\mathrm{}+j}^{(j)}=0,𝝃^{(j)}𝒦_k,\hfill \end{array}`$
and
$`\begin{array}{ccc}(P_{\mathrm{}+j})\hfill & \text{minimize}\hfill & 𝒖^T𝝃^{(\mathrm{}+j)}\hfill \\ & \text{subject to}\hfill & \xi _j^{(\mathrm{}+j)}=0,\xi _{\mathrm{}+j}^{(\mathrm{}+j)}=1,𝝃^{(\mathrm{}+j)}𝒦_k.\hfill \end{array}`$
If $`(P_j)`$ is infeasible then add the equation $`x_j=0`$ (or the inequality $`x_j0`$, since the inequality $`x_j0`$ is already included) to $`(C_{k+1})`$. If $`(P_{\mathrm{}+j})`$ is infeasible then add the equation $`s_j=0`$ to $`(C_{k+1})`$. Otherwise, let $`(𝝃^{(j)})^{}`$ and $`(𝝃^{(\mathrm{}+j)})^{}`$ denote the optimal solutions of $`(P_j)`$ and $`(P_{\mathrm{}+j})`$ respectively. Define $`y_ju_j𝒖^T(𝝃^{(j)})^{}`$, $`y_{\mathrm{}+j}u_{\mathrm{}+j}𝒖^T(𝝃^{(\mathrm{}+j)})^{}`$. Add the inequality
$$\underset{ij}{}\left(u_ix_i+u_{\mathrm{}+i}s_i\right)y_jx_jy_{\mathrm{}+j}s_ju_{2\mathrm{}+1}\alpha u_0$$
to $`(C_{k+1})`$.
Step 3. Let $`k=k+1`$, and go to Step 1.
Note that in iteration $`k`$, the algorithm solves $`(2\mathrm{}|(C_k)|)`$ LP problems.
###### Theorem 4.2
Let $`C_k`$, $`k\{1,2,\mathrm{}\}`$ be the sequence of convex relaxations generated by Algorithm 4.1. Then $`C_{\mathrm{}}=\text{c.hull}(F)`$.
Proof: We think of $`𝒦_k`$ for all $`k0`$, as a subset of $`R^{1+(2\mathrm{}+1)}`$, with the $`0`$th component being the homogenizing variable, the next $`\mathrm{}`$ components representing $`𝒙`$, the next $`\mathrm{}`$ components representing $`𝒔`$ and the last component representing $`\alpha `$. Note that
$$𝒦_1\left(𝒦_0\{𝒗:x_j=0\}\right)+\left(𝒦_0\{𝒗:s_j=0\}\right)$$
iff
$`𝒦_1^{}\left(𝒦_0^{}+\{𝒆_j\}\right)\left(𝒦_0^{}+\{𝒆_{\mathrm{}+j}\}\right).`$ (19)
(We used the fact that $`𝒦_0R_+^{1+(2\mathrm{}+1)}`$.) Therefore, if we ensure the latter inclusion, then Theorem 3.2 applies and we can conclude the convergence of the method in $`\mathrm{}`$ iterations. Recall that every vector $`𝒖𝒦_0^{}`$ represents a valid inequality
$$\underset{i=1}{\overset{\mathrm{}}{}}\left(u_ix_i+u_{\mathrm{}+i}s_i\right)u_{2\mathrm{}+1}\alpha u_0$$
for $`C_0`$. To ensure the inclusion (19), it suffices to prove:
$$\text{“For every }𝒖,𝒘𝒦_0^{}\text{ such that }u_i=w_i,i\{j,\mathrm{}+j\};u_jw_j,u_{\mathrm{}+j}w_{\mathrm{}+j},$$
$$\text{we have }𝒚𝒦_1^{},\text{ where }y_iu_i,ij;y_jw_j\text{.”}$$
This is equivalent to proving the fact that if the two inequalities
$`{\displaystyle \underset{i=1}{\overset{\mathrm{}}{}}}\left(u_ix_i+u_{\mathrm{}+i}s_i\right)u_{2\mathrm{}+1}\alpha u_0,\text{ and}`$
$`{\displaystyle \underset{ij}{}}\left(u_ix_i+u_{\mathrm{}+i}s_i\right)w_jx_jw_{\mathrm{}+j}s_ju_{2\mathrm{}+1}\alpha u_0`$
are valid for $`C_0`$, then
$$\underset{ij}{}\left(u_ix_i+u_{\mathrm{}+i}s_i\right)w_jx_ju_{\mathrm{}+j}s_ju_{2\mathrm{}+1}\alpha u_0$$
is valid for $`C_1`$. To compute all such inequalities defining $`C_1`$, we solve for every valid inequality
$$\underset{i=1}{\overset{\mathrm{}}{}}\left(u_ix_i+u_{\mathrm{}+i}s_i\right)u_{2\mathrm{}+1}\alpha u_0$$
for $`C_0`$ and every $`j\{1,2,\mathrm{},\mathrm{}\}`$, the linear programming problems
$`\begin{array}{ccc}& \text{maximize}\hfill & \beta \hfill \\ & \text{subject to}\hfill & \beta 𝒆_j+\delta 𝒆_{\mathrm{}+j}_{𝒦_0^{}}𝒖,\hfill \end{array}`$
and
$`\begin{array}{ccc}& \text{maximize}\hfill & \gamma \hfill \\ & \text{subject to}\hfill & \kappa 𝒆_j+\gamma 𝒆_{\mathrm{}+j}_{𝒦_0^{}}𝒖.\hfill \end{array}`$
Here, $`_{𝒦_0^{}}`$ denotes the partial order induced by the convex cone $`𝒦_0^{}`$ (that is, $`𝒖^1_{𝒦_0^{}}𝒖^2`$ iff $`(𝒖^2𝒖^1)𝒦_0^{}`$). Note that both problems are always feasible. Therefore, each of them either has an optimal solution or is unbounded. If both LPs have optimal solutions, say $`\beta ^{}`$ and $`\gamma ^{}`$ then we set $`w_ju_j\beta ^{}`$ and $`u_{\mathrm{}+j}u_{\mathrm{}+j}\gamma ^{}`$. Since the above two problems are LPs, we can equivalently solve their duals. Namely, we solve the LPs:
$`\begin{array}{ccc}(P_j)\hfill & \text{minimize}\hfill & 𝒖^T𝝃^{(j)}\hfill \\ & \text{subject to}\hfill & \xi _j^{(j)}=1,\xi _{\mathrm{}+j}^{(j)}=0,𝝃^{(j)}𝒦_0,\hfill \end{array}`$
and
$`\begin{array}{ccc}(P_{\mathrm{}+j})\hfill & \text{minimize}\hfill & 𝒖^T𝝃^{(\mathrm{}+j)}\hfill \\ & \text{subjcet to}\hfill & \xi _j^{(\mathrm{}+j)}=0,\xi _{\mathrm{}+j}^{(\mathrm{}+j)}=1,𝝃^{(\mathrm{}+j)}𝒦_0.\hfill \end{array}`$
These latter two linear programming problems are precisely the ones used by Algorithm 4.1. Notice that since their duals are either unbounded or have optimal solutions, these LP problems either have optimal solutions or are infeasible. When $`(P_j)`$ is infeasible, the equality $`x_j=0`$ is valid for $`F`$ and the algorithm adds this equality to the describing inequalities of $`C_k`$. Similarly, when $`(P_{\mathrm{}+j})`$ is infeasible, $`s_j=0`$ is valid for $`F`$ and the algorithm behaves correctly in this instance. (In either instance, the inclusion (19) is obviously satisfied for $`j`$.) However, the proof is not yet complete; because, the arguments so far ensure the inclusion (19) when the algorithm is ran for every valid inequality of $`C_0`$. So, next we prove that what the algorithm does (using only the facets of $`C_0`$) suffices. To see this, we need to prove that to derive the facets of $`𝒦_1`$, it suffices to start with a facet $`𝒖`$ of $`𝒦_0`$ in the above procedure. Suppose $`𝒖,𝒘𝒦_0^{}`$ satisfy the above conditions but $`𝒖`$ is not facet inducing for $`𝒦_0^{}`$. (We will prove that the valid inequality derived from $`𝒖`$ and $`𝒘`$ is implied by some other inequalities derived from some facets $`𝒖^1,𝒖^2,\mathrm{},𝒖^{\mathrm{}}`$ of $`𝒦_0`$.) Since $`𝒖`$ is not facet inducing for $`𝒦_0`$, $`𝒖`$ is not an extreme ray of $`𝒦_0^{}`$. Hence, there exist extreme rays $`𝒖^1,𝒖^2,\mathrm{},𝒖^{\mathrm{}}`$ of $`𝒦_0^{}`$ such that for some $`\lambda _r>0,`$ $`r\{1,2,\mathrm{},\mathrm{}\}`$, $`_{r=1}^{\mathrm{}}\lambda _r=1`$ the following conditions are satisfied:
$`𝒖`$ $`=`$ $`{\displaystyle \underset{r=1}{\overset{\mathrm{}}{}}}\lambda _r𝒖^r,u_0^r=u_0,r\{1,2,\mathrm{},\mathrm{}\}.`$
Note that $`𝒖^r`$ is facet inducing for each $`r`$. Let $`𝝃^r`$ be the optimal solution of $`(P_j)`$ above for the objective function vector $`𝒖^r`$. Let $`𝝃^{}`$ be an optimal solution of $`(P_j)`$ when the objective function vector is $`𝒖`$. We claim that there exists $`\stackrel{~}{𝝃}𝒦_0`$ such that
$`(𝒖^r)^T\stackrel{~}{𝝃}`$ $`=`$ $`(𝒖^r)^T𝝃^r,r\{1,2,\mathrm{},\mathrm{}\},\stackrel{~}{\xi }_j=1,\stackrel{~}{\xi }_{\mathrm{}+j}=0,\stackrel{~}{𝝃}𝒦_0.`$
(This claim follows from Farkas’ Lemma, using the facts that $`𝒖^r𝒦_0^{},r`$ and $`𝝃^r𝒦_0,r`$.) Thus, we have
$$\underset{r=1}{\overset{\mathrm{}}{}}\lambda _r(𝒖^r)^T𝝃^r=𝒖^T\stackrel{~}{𝝃}𝒖^T𝝃^{}.$$
Therefore, the inequality obtained from $`𝒖`$ is equivalent to or dominated by a nonnegative combination of the inequalities obtained from $`𝒖^r`$ which induce facets of $`𝒦_0`$. The proof is complete.
We illustrated a derivation and convergence proof for a successive relaxation method (closely related to Balas’ approach and analogous to a suggestion of Lovász and Schrijver ) based on Lemma 3.1 and Theorem 3.2. Algorithm 4.1 is an analog of a method based on relaxations $`N_0^k(𝒦)`$ from (which is concerned with the case of 0-1 integer programming). For the relationship of the methods of and , see Balas, Ceria and Cornuejols . (Balas’ method , in essence, corresponds to defining
$$𝒦_{k+1}(𝒦_k\{𝒗:x_{k+1}=0\})+(𝒦_0\{𝒗:(𝑴𝒙+𝒒\alpha )_{k+1}=0\}).)$$
Let $`C_k^{(4)}`$, $`k0`$ denote the projection of $`C_k`$ generated by Algorithm 4.1 onto the coordinates $`\left(\begin{array}{c}𝒙\\ \alpha \end{array}\right)`$. Let $`C_k^{(3)}`$, $`k0`$ denote the projection of $`C_k`$, generated by Algorithm 1.1, as used in Section 3, onto the coordinates $`\left(\begin{array}{c}𝒙\\ \alpha \end{array}\right)`$. Let $`𝒦_k^{(4)}`$ denote the convex cone associated with $`C_k^{(4)}`$. From the proof of Theorem 4.2, it is easy to see that
$$𝒦_{k+1}^{(4)}=\underset{i=1}{\overset{\mathrm{}}{}}\left[\left(𝒦_k^{(4)}\{𝒗:x_i=0\}\right)+\left(𝒦_k^{(4)}\{𝒗:s_i=0\}\right)\right].$$
Therefore, the proofs of Theorems 3.2 and 4.2 imply that
$$\text{if }C_0^{(4)}C_0^{(3)}\text{ then }C_k^{(4)}C_k^{(3)}\text{ for all }k0.$$
Thus, the SSILPR Method (Algorithm 1.1) as applied in Section 3 to ($`\text{MIP}_\alpha `$) converges at least as fast as Algorithm 4.1 applied to ($`\text{LCP}_\alpha `$).
## 5 SCRMs applied to the smaller formulation of LCP with an implicit treatment of the disjunctive constraints
We have already seen various ways of applying SCRMs to LCP problems. Since the methods proposed in only require a formulation of the feasible solutions by quadratic inequalities, we are also interested in applying the methods of to the following formulation:
$$C_0\{\left(\begin{array}{c}\alpha \\ 𝒙\end{array}\right)R^{\mathrm{}+1}:\begin{array}{c}𝑴𝒙+𝒒\alpha \text{0},𝒙\text{0},\alpha \mathrm{\hspace{0.17em}\hspace{0.17em}0},\hfill \\ 𝒆^T(𝑴+𝑰)𝒙+(𝒆^T𝒒+1)\alpha \mathrm{\hspace{0.17em}\hspace{0.17em}1}\hfill \end{array}\},$$
and
$$𝒫_F\left\{x_i(𝑴𝒙+𝒒\alpha )_i\mathrm{\hspace{0.17em}\hspace{0.17em}0},i\{1,2,\mathrm{},\mathrm{}\}\right\}.$$
The general theory of Kojima-Tunçel implies that their SSDPR and SSILPR Methods converge. It would be interesting to characterize the conditions under which the Algorithms 3.1 and 3.2 of converge in at most $`\mathrm{}`$ iterations for the above description of $`𝒫_F`$ and $`C_0`$. Also see , where the authors derived some necessary and some sufficient conditions for the finite convergence of SCRMs.
## 6 A general linear complementarity problem
Let $`𝒜:R^{\mathrm{}}R^{\mathrm{}}`$, a linear transformation, $`𝒒R^{\mathrm{}}`$ and $`𝒦R^{\mathrm{}}`$ a pointed, closed convex cone with nonempty interior, be given. Consider the Complementarity Problem (CP):
(CP) Find $`𝒙`$, $`𝒔`$ such that $`𝒜(𝒙)+𝒒=𝒔,`$
$`𝒙𝒦,𝒔𝒦^{},𝒙,𝒔=\mathrm{\hspace{0.17em}\hspace{0.17em}0},`$
where $`𝒦^{}`$ is the dual of $`𝒦`$:
$$𝒦^{}\{𝒔R^{\mathrm{}}:𝒙,𝒔\mathrm{\hspace{0.17em}\hspace{0.17em}0},𝒙𝒦\}.$$
Since $`𝒦`$ is a pointed, closed convex cone with nonempty interior, so is $`𝒦^{}`$. Such problems were studied recently, in the context of interior-point methods . We pick $`𝜼\text{int}(𝒦)`$, $`\overline{𝜼}\text{int}(𝒦^{})`$ and we can solve instead the optimization problem
$`\begin{array}{ccc}(\text{CP}_\alpha )\hfill & \text{maximize}\hfill & \alpha \hfill \\ & \text{subject to}\hfill & 𝒙𝒦,\left[𝒜(𝒙)+𝒒\alpha \right]𝒦^{},\alpha \mathrm{\hspace{0.17em}\hspace{0.17em}0},\hfill \\ & & \overline{𝜼},𝒙+𝜼,𝒜(𝒙)+𝒒\alpha +\alpha \mathrm{\hspace{0.17em}\hspace{0.17em}1},\hfill \\ & & 𝒙,𝒜(𝒙)+\alpha 𝒒=\mathrm{\hspace{0.17em}\hspace{0.17em}0}.\hfill \end{array}`$
We choose
$$C_0\{\left(\begin{array}{c}\alpha \\ 𝒙\end{array}\right)R^{\mathrm{}+1}:\begin{array}{c}𝒙𝒦,\left[𝒜(𝒙)+𝒒\alpha \right]𝒦^{},\alpha \mathrm{\hspace{0.17em}\hspace{0.17em}0},\hfill \\ \overline{𝜼},𝒙+𝜼,𝒜(𝒙)+𝒒\alpha +\alpha \mathrm{\hspace{0.17em}\hspace{0.17em}1}\hfill \end{array}\}.$$
Note that $`C_0`$ is always a compact convex set (see the next theorem). We also pick
$$𝒫_F\left\{\begin{array}{c}𝒙,𝒜(𝒙)+\alpha 𝒒,𝒙,𝒜(𝒙)+\alpha 𝒒\end{array}\right\}.$$
###### Theorem 6.1
* $`C_0`$ is a compact convex set.
* (CP<sub>α</sub>) has an optimal solution with $`\alpha ^{}>0`$ iff (CP) has a solution (or solutions).
* If $`\left(\begin{array}{c}\alpha ^{}\\ 𝒙^{}\end{array}\right)`$ is an optimal solution of (CP<sub>α</sub>) with $`\alpha ^{}>0`$ then the pair of vectors
$`({\displaystyle \frac{𝒙^{}}{\alpha ^{}}},{\displaystyle \frac{1}{\alpha ^{}}}𝒜(𝒙^{})+𝒒)`$ solves (CP).
Proof:
* We only need to show that $`C_0`$ is bounded; because, $`C_0`$ is a closed and convex subset of $`R^{\mathrm{}+1}`$ by definition. Assume on the contrary that we can take an unbounded direction $`\left(\begin{array}{c}\mathrm{\Delta }\alpha \\ \mathrm{\Delta }𝒙\end{array}\right)\text{0}`$ in $`C_0`$;
$`\left(\begin{array}{c}\mathrm{\Delta }\alpha \\ \mathrm{\Delta }𝒙\end{array}\right)\text{0},\mathrm{\Delta }𝒙𝒦,\mathrm{\Delta }\alpha 0,\left[𝒜(\mathrm{\Delta }𝒙)+𝒒\mathrm{\Delta }\alpha \right]𝒦^{},`$
$`\overline{𝜼},\mathrm{\Delta }𝒙+𝜼,𝒜(\mathrm{\Delta }𝒙)+𝒒\mathrm{\Delta }\alpha +\mathrm{\Delta }\alpha \mathrm{\hspace{0.17em}\hspace{0.17em}0}.`$
Since each term in the left hand side of the last inequality is nonnegative, we have
$`\overline{𝜼},\mathrm{\Delta }𝒙=0\text{ and }\mathrm{\Delta }\alpha =0.`$
Since $`\overline{𝜼}\text{int}(𝒦^{})`$ and $`\mathrm{\Delta }𝒙𝒦`$, the first identity above implies that $`\mathrm{\Delta }𝒙=\text{0}`$. Thus, we have a contradiction to $`\left(\begin{array}{c}\mathrm{\Delta }\alpha \\ \mathrm{\Delta }𝒙\end{array}\right)\text{0}`$.
* Suppose (CP<sub>α</sub>) has an optimal solution $`\left(\begin{array}{c}\alpha ^{}\\ 𝒙^{}\end{array}\right)`$ with $`\alpha ^{}>0`$. Then $`\overline{𝒙}{\displaystyle \frac{x^{}}{\alpha ^{}}}𝒦`$, $`\overline{𝒔}{\displaystyle \frac{1}{\alpha ^{}}}𝒜(𝒙^{})+𝒒𝒦^{}.`$ We have
$$\overline{𝒙},\overline{𝒔}=\overline{𝒙},𝒜(\overline{𝒙})+𝒒=\frac{1}{(\alpha ^{})^2}𝒙^{},𝒜(𝒙^{})+\alpha ^{}𝒒=0.$$
Therefore, $`(\overline{𝒙},\overline{𝒔})`$ solves (CP). For the converse, let $`(\overline{𝒙},\overline{𝒔})`$ be a solution of (CP). Let
$$\zeta \overline{𝜼},\overline{𝒙}+𝜼,\overline{𝒔}\mathrm{\hspace{0.17em}\hspace{0.17em}0},\alpha ^{}=\frac{1}{\zeta +1}\text{ and }𝒙^{}=\frac{\overline{𝒙}}{\zeta +1}.$$
Then $`\left(\begin{array}{c}\alpha ^{}\\ 𝒙^{}\end{array}\right)`$ is a feasible solution of (CP<sub>α</sub>). But the feasible region of (CP<sub>α</sub>) is compact and nonempty, its objective function is linear, hence, (CP<sub>α</sub>) has optimal solution (or solutions). Since we already showed a solution with positive objective value, the optimum value is positive.
* This claim follows from the proof of (ii).
Theorem 6.1 shows that we can apply SCRMs to (CP<sub>α</sub>) with the above $`C_0`$ and $`𝒫_F`$ and solve the original, general problem (CP).
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# Quantum Monte Carlo calculations of pion scattering from Li
## Abstract
We show that the neutron and proton transition densities predicted by recent Quantum Monte Carlo calculations for $`A=6,7`$ nuclei are consistent with pion scattering from <sup>6</sup>Li and <sup>7</sup>Li at energies near the $`\mathrm{\Delta }`$ resonance. This has provided a microscopic understanding of the enhancement factors for quadrupole excitations, which were needed to describe pion inelastic scattering within the nuclear shell model of Cohen and Kurath.
Quantum Monte Carlo (QMC) methods have been successfully developed to predict the properties of low-lying states of light nuclei starting with realistic two- and three-nucleon potentials. While the reproduction of the energy spectra is the essential first step in such an effort, the dynamical content of the resulting nuclear wave functions must be tested against various reactions. This has been achieved so far mainly by considering electroweak processes. In this paper we report on an additional test by using pion elastic and inelastic scattering.
Let us first briefly review the status of our understanding of pion inelastic scattering at medium energies( 80 MeV $`<E_{lab}<`$ 300 MeV). Because of the excitation of the $`\mathrm{\Delta }`$ resonance, the pion-nucleus interactions in this energy region are dominated by the strong absorption mechanism. Consequently, the pion-nucleus inelastic scattering leading to discrete final nuclear states can be described by the Distorted Wave Impulse Approximation (DWIA). This has been well established in very extensive investigations of the data from meson factories. Following the momentum-space approach, the inelastic scattering amplitude can be written as
$$T_{fi}(\stackrel{}{k}_0^{},\stackrel{}{k}_0)=𝑑\stackrel{}{k}^{}𝑑\stackrel{}{k}\chi _{\stackrel{}{k}_0,f}^{()}(\stackrel{}{k}^{})U_{fi}(\stackrel{}{k}^{},\stackrel{}{k})\chi _{\stackrel{}{k}_0,i}^{(+)}(\stackrel{}{k}),$$
(1)
where the $`\stackrel{}{k}`$’s are pion-nucleus relative momenta in the pion-nucleus center of mass frame, and the distorted waves $`\chi ^{(\pm )}`$ are generated from an optical potential which is adjusted to fit the pion-nucleus elastic scattering.
The nuclear excitations are contained in the transition potential $`U_{fi}`$. It can be calculated from the $`\pi N`$ scattering t-matrix and nuclear transition form factors. The details are given in Ref.. To simplify the presentation, the spin-isospin variables will be suppressed here. Then the transition potential can be written as
$`U_{fi}(\stackrel{}{k}^{},\stackrel{}{k})=t_{\pi n}(\stackrel{}{k}^{},\stackrel{}{k},\omega _0)f_{fi}^{1/2}(\stackrel{}{q})+t_{\pi p}(\stackrel{}{k}^{},\stackrel{}{k},\omega _0)f_{fi}^{1/2}(\stackrel{}{q}),`$ (2)
where $`\stackrel{}{q}=\stackrel{}{k}^{}\stackrel{}{k}`$, $`t_{\pi n}(t_{\pi p})`$ is an appropriately parameterized pion-neutron (pion-proton) scattering amplitude, and $`\omega _0`$ is the collision energy calculated from using the fixed-scatterer approximation. With $`t_3=1/2,1/2`$ denoting neutron and proton respectively, the nuclear transition form factors are defined by
$`f_{fi}^{t_3}(\stackrel{}{q})={\displaystyle 𝑑\stackrel{}{r}e^{i\stackrel{}{r}\stackrel{}{q}}\rho _{fi}^{t_3}(\stackrel{}{r})},`$ (3)
with the transition densities defined by
$`\rho _{fi}^{t_3}(\stackrel{}{r})=<\mathrm{\Psi }_f{\displaystyle \frac{1}{A}}{\displaystyle \underset{i=1,A}{}}\delta (\stackrel{}{r}\stackrel{}{r}_i){\displaystyle \frac{1+2t_3\tau _z(i)}{2}}\mathrm{\Psi }_i>,`$ (4)
where $`\tau _z`$ is the z-component of the nucleon isospin operator, and $`\mathrm{\Psi }_i`$ and $`\mathrm{\Psi }_f`$ are the initial and final nuclear states respectively. The transition dynamics can be better understood from the multipole expansion of transition densities
$`\rho _{fi}^{t_3}(\stackrel{}{r})={\displaystyle \underset{KM}{}}{\displaystyle \frac{1}{r^2}}Y_{KM}(\widehat{r})<J_fM_fJ_iKM_iM>F_{K0,t_3}^{fiK}(r),`$ (5)
Here we recall the notation of Ref. to include possible spin transitions $`S=0,1`$ in $`F_{KS,t_3}^{fiJ}(r)`$, where $`K`$ denotes the orbital angular momentum transition and $`\stackrel{}{J}=\stackrel{}{K}+\stackrel{}{S}`$ is the total angular momentum transfer.
In the DWIA study for 1p-shell nuclei, the transition densities were calculated from the shell model of Cohen and Kurath:
$`F_{KS,t_3}^{fiJ}(r)`$ $`=`$ $`{\displaystyle \underset{\alpha \beta }{}}A_{J(KS)t_3}(\alpha ,\beta ;fi)(4\pi j_\alpha )^{1/2}<(l_\alpha 1/2)j_\alpha [Y_K(\widehat{r})\times \sigma _S]_J(l_\beta 1/2)j_\beta >`$ (6)
$`\times `$ $`R_{n_\alpha l_\alpha }(r)R_{n_\beta ,l_\beta }(r)`$ (7)
with $`\sigma _0=1`$, $`\sigma _1=\stackrel{}{\sigma }`$, and
$`A_{J(KS)t_3}(\alpha ,\beta ;fi)=<\mathrm{\Psi }_f[b_{\alpha t_3}^{}\times h_{\beta t_3}^{}]_{J[KS]}\mathrm{\Psi }_i>,`$ (8)
where $`\alpha (n_\alpha ,l_\alpha ,j_\alpha )`$ denotes the single particle orbitals, $`R_\alpha (r)`$ is the radial wave function, and $`b_\alpha ^{}`$ and $`h_\beta ^{}`$ are the creation operators for the particle and hole states respectively.
It was found that the pion inelastic scattering from 1p-shell nuclei can be described by using the above shell-model input only when the quadrupole excitation component $`J(KS)=2(20)`$ is enhanced by a factor $`E_N2`$. In Ref., these enhancement factors were estimated from a systematic analysis of $`B(E2)`$ transitions from 1p-shell nuclei. For proton excitation, this enhancement factor is consistent with what is needed for explaining $`B(E2)`$ values. For $`Z=N`$ nuclei, one can assume that the neutron excitation also has the same enhancement because of isospin invariance. However, for the $`NZ`$ nuclei, such as <sup>7</sup>Li, the enhancement factors for neutron excitations can not be obtained without making some additional assumptions. The predicted pion inelastic cross sections thus are not well justified theoretically. Furthermore, it would be desirable if the calculations do not include any enhancement factors. This however is very difficult, if not impossible, in practice within the shell model since the collective quadrupole excitations can only be described by a very large model space.
In this work, we calculate the transition densities for a given multipolarity by using wave functions from the recent QMC calculations for light nuclei. The input, Eq.(6), to the DWIA calculations for pion inelastic scattering is then defined by the following matrix element
$`F_{KS,t_3}^{fiJ}(r)={\displaystyle \frac{<\mathrm{\Psi }_{J_fM_f}_{i=1,A}\delta (rr_i)r_i^K[Y_K(\widehat{r}_i)\times \sigma _S]_{JM}\frac{1+2t_3\tau (i)}{2}\mathrm{\Psi }_{J_iM_i}>}{<J_fM_fJ_iJM_iM>}}`$ (9)
The QMC calculations use a realistic Hamiltonian containing the Argonne $`v_{18}`$ two-nucleon and Urbana IX three-nucleon potentials, which we refer to as the AV18/UIX model. Both variational (VMC) and Green’s function (GFMC) Monte Carlo calculations have been made for light nuclei . The AV18/UIX model reproduces the experimental binding energies and charge radii of <sup>3</sup>H, <sup>3</sup>He, and <sup>4</sup>He, in the numerically exact GFMC calculations, but underbinds <sup>6</sup>Li and <sup>7</sup>Li by 2–5%. The variational Monte Carlo (VMC) energies are 2% above the GFMC results for $`A=3,4`$ nuclei and 10% above for $`A=6,7`$. However, the known excitation spectra are well reproduced by both the VMC and GFMC calculations, as are the charge radii. The VMC and GFMC calculations also produce very similar one-body densities, while two-nucleon density distributions differ by less than 10%.
The VMC wave functions have been used successfully to describe the elastic and transition electromagnetic form factors for <sup>6</sup>Li without introducing effective charges. They have also given an excellent absolute prediction for the spectroscopic factors in <sup>7</sup>Li($`e,e^{}p`$) reaction . Consequently we expect the VMC wave functions to give a good estimate for both the elastic and transition densities required in pion scattering calculations.
The variational wave function for $`A=6,7`$ nuclei used here is the trial wave function, $`\mathrm{\Psi }_T`$, that serves as the starting point for the GFMC calculations. It has the general form
$$|\mathrm{\Psi }_T=\left[1+\underset{i<j<k}{}\stackrel{~}{U}_{ijk}^{TNI}\right]\left[𝒮\underset{i<j}{}(1+U_{ij})\right]|\mathrm{\Psi }_J,$$
(10)
where $`U_{ij}`$ and $`\stackrel{~}{U}_{ijk}^{TNI}`$ are two- and three-body correlation operators and the Jastrow wave function $`|\mathrm{\Psi }_J`$ is given by
$`|\mathrm{\Psi }_J`$ $`=`$ $`𝒜\{{\displaystyle \underset{i<j<k4}{}}f_{ijk}^c{\displaystyle \underset{i<j4}{}}f_{ss}(r_{ij}){\displaystyle \underset{k4<lA}{}}f_{sp}(r_{kl})`$ (12)
$`{\displaystyle \underset{LS}{}}(\beta _{LS[n]}{\displaystyle \underset{4<l<mA}{}}f_{pp}^{LS[n]}(r_{lm})|\mathrm{\Phi }_A(LS[n]JMTT_3)_{1234:56\mathrm{}A})\}.`$
The $`𝒮`$ and $`𝒜`$ are symmetrization and antisymmetrization operators, respectively. The central pair and triplet correlations $`f_{xy}(r_{ij})`$ and $`f_{ijk}^c`$ are functions of relative positions only; the subscripts $`xy`$ denote whether the particles are in the s- or p-shell. The $`|\mathrm{\Phi }_A(LS[n]JMTT_3)`$ is a single-particle wave function with orbital angular momentum $`L`$, spin $`S`$, and spatial symmetry $`[n]`$ coupled to total angular momentum $`J`$, projection $`M`$, isospin $`T`$, and charge state $`T_3`$:
$`|\mathrm{\Phi }_A(LS[n]JMTT_3)_{1234:56\mathrm{}A}=|\mathrm{\Phi }_\alpha (0000)_{1234}{\displaystyle }_{4<lA}\varphi ^{LS}_p(R_{\alpha l})`$ (13)
$`\{[{\displaystyle \underset{4<lA}{}}Y_{1m_l}(\mathrm{\Omega }_{\alpha l})]_{LM_L[n]}\times [{\displaystyle \underset{4<lA}{}}\chi _l({\displaystyle \frac{1}{2}}m_s)]_{SM_S}\}_{JM}\times [{\displaystyle \underset{4<lA}{}}\nu _l({\displaystyle \frac{1}{2}}t_3)]_{TT_3}.`$ (14)
Particles 1–4 are placed in an $`\alpha `$ core with only spin-isospin degrees of freedom, denoted by $`\mathrm{\Phi }_\alpha (0000)`$, while particles 5–A are placed in $`p`$-wave orbitals $`\varphi _p^{LS}(R_{\alpha l})`$ that are functions of the distance between the center of mass of the $`\alpha `$ core and particle $`l`$. Different amplitudes $`\beta _{LS[n]}`$ are mixed to obtain an optimal wave function by means of a small-basis diagonalization. For <sup>6</sup>Li, the $`(J^\pi ;T)=(1^+;0)`$ ground state is predominantly a <sup>3</sup>S amplitude, with small admixtures of <sup>3</sup>D and <sup>1</sup>P components, while the $`(3^+;0)`$ first excited state is pure <sup>3</sup>D. For <sup>7</sup>Li, the $`(J^\pi ;T)=(\frac{3}{2}^{};\frac{1}{2})`$ ground and $`(\frac{1}{2}^{};\frac{1}{2})`$ first excited states are predominantly <sup>2</sup>P, with small admixtures of <sup>2,4</sup>P, <sup>2,4</sup>D, and <sup>2</sup>S components. The $`(J^\pi ;T)=(\frac{7}{2}^{};\frac{1}{2})`$ and $`(\frac{5}{2}^{};\frac{1}{2})`$ excited states are predominantly <sup>2</sup>F, again with small admixtures of <sup>2,4</sup>P and <sup>2,4</sup>D components. Mixing parameter values are given in Ref. .
The two-body correlation operator $`U_{ij}`$ is defined as:
$$U_{ij}=\underset{p=2,6}{}\left[\underset{ki,j}{}f_{ijk}^p(𝐫_{ik},𝐫_{jk})\right]u_p(r_{ij})O_{ij}^p,$$
(15)
where the $`O_{ij}^{p=2,6}`$ = $`𝝉_i𝝉_j`$, $`𝝈_i𝝈_j`$, $`𝝈_i𝝈_j𝝉_i𝝉_j`$, $`S_{ij}`$, and $`S_{ij}𝝉_i𝝉_j`$. The six radial functions $`f_{ss}(r)`$ and $`u_{p=2,6}(r)`$ are obtained from two-body Euler-Lagrange equations with variational parameters . The $`f_{sp}`$ and $`f_{pp}^{LS[n]}`$ correlations are similar to $`f_{ss}`$ for small separations, but include parameterized long-range tails. The parameters used in constructing these two-body correlations, as well as the description of the three-body correlation operator $`\stackrel{~}{U}_{ijk}^{TNI}`$ and the operator-independent three-body correlations $`f_{ijk}^c`$ and $`f_{ijk}^p`$ are given in Ref. .
In Ref. it was found that inelastic transitions induced by pion scattering are dominated by the quadrupole transition $`J(KS)=2(20)`$. In our QMC calculations, we evaluate the quadrupole transition density:
$$\rho _{E2}^{t_3}(r)=\frac{\sqrt{2J_f+1}\mathrm{\Psi }_{J_fM_f}_{i=1,A}\delta (rr_i)r_i^2Y_2^M(\widehat{r}_i)\frac{1+2t_3\tau (i)}{2}\mathrm{\Psi }_{J_iM_i}}{J_fM_fJ_i2M_iM}.$$
(16)
These neutron and proton transition densities are shown in Fig. 1 for four transitions in <sup>6</sup>Li and <sup>7</sup>Li. The integrated $`B(E2,t_3)`$ values,
$$B(E2,t_3)=\frac{\rho _{E2}^{t_3}(r)d^3r^2}{2J_i+1},$$
(17)
are given in Table I, where they are compared to the experimental proton values obtained from $`(e,e^{})`$ scattering and Coulomb excitation experiments . Our evaluations are made with 160 000 Monte Carlo samples for transitions in <sup>6</sup>Li and 120 000 Monte Carlo samples for transitions in <sup>7</sup>Li. This number of samples is sufficient to give statistical uncertainties that are as small or smaller than the errors on the experimental $`B(E2)`$ values. Other $`J(KS)`$ transition amplitudes can contribute to pion inelastic scattering, particularly in the <sup>6</sup>Li case, and are evaluated in a similar manner.
We see from Fig. 1 and Table I that the predicted differences between the neutron and proton excitations in <sup>7</sup>Li are very significant. Such differences can be most effectively verified by using an important characteristic of pion scattering at energies near the $`\mathrm{\Delta }`$ excitation. At a typical energy $`E_\pi =164`$ MeV, one finds that the $`\pi N`$ amplitude in Eq.(2) has an interesting ratio $`t_{\pi ^+p}/t_{\pi ^+n}=t_{\pi ^{}n}/t_{\pi ^{}p}3`$. Consequently, the $`\pi ^+`$ scattering is dominated by the proton excitations while $`\pi ^{}`$ scattering is dominated by neutron excitations. The agreement with both the $`\pi ^+`$ and $`\pi ^{}`$ data will be a nontrivial test of the QMC wave functions. Thus the present study is complementary to that of Ref. using electron scattering which mainly probes the proton excitations.
We first investigate pion scattering from <sup>6</sup>Li. Here data for $`E_\pi =100`$, 180 and 240 MeV are available for testing the energy-dependence of our predictions. The pion optical potential and the $`\pi N`$ t-matrix are taken from Ref.. The $`\pi N`$ amplitudes we employ are taken from the Karlsruhe-Helsinki analysis, and differ only slightly with the more recent VPI analysis, mainly in the $`S_{11}`$ partial wave (as discussed in Ref.). This partial wave and the other non-$`P_{33}`$ partial waves are much weaker than the $`P_{33}`$ channel in the energy region of interest near the $`\mathrm{\Delta }`$ excitation. For the present exploratory investigation, we therefore do not make any effort to improve the optical potential employed in Ref.. Such an improvement is probably needed in the future when the data at low energies and for the spin observales are investigated.
With transition densities calculated from QMC, there are no adjustable parameters in our DWIA calculations. In Table I, we see that the calculated $`B(E2)`$ for the transition to the $`(3^+,T=0)`$ state is in excellent agreement with the data. This is a significant improvement over the shell-model prediction which required a large enhancement factor $`E_p=E_n=2.5`$ to reproduce the $`B(E2)`$ data, as discussed in Ref.. We thus expect a similar improvement in pion scattering calculations.
Our results for <sup>6</sup>Li are displayed in Fig. 2. We see good agreement with the differential cross sections for elastic scattering. However, the discrepancies in reproducing the diffractive minima at 180 MeV indicate some deficiencies of the simple optical potential we have employed. We also see general agreement with the inelastic scattering to the $`(3^+;T=0)`$ excited state, although some noticeable discrepancies are seen, particularly at 100 MeV. Nevertheless, it is fair to say that our results agree with the differential cross sections to a very large extent in both absolute magnitude and energy-dependence. The agreement seen in Fig. 2 is consistent with the $`B(E2)`$ values listed in Table I. The overall agreement is not surprising in view of the ability of the QMC wave functions to reproduce the elastic and transition form factors in electron scattering experiments . To further improve the agreement with the data and to account for the spin observables, it would be necessary to improve the reaction model. For example, we may have to consider coupled-channel effects and refine the optical potential, but this is beyond the scope of the present investigation.
We next investigate the very old data for 164 MeV $`\pi ^+`$ and $`\pi ^{}`$ inelastic scattering from a <sup>7</sup>Li target . The final states we consider are $`(J^\pi ;T)=(\frac{1}{2}^{};\frac{1}{2})`$ at 0.478 MeV, $`(\frac{7}{2}^{};\frac{1}{2})`$ at 4.63 MeV, and $`(\frac{5}{2}^{};\frac{1}{2})`$ at 6.68 MeV. Our results are shown in Fig. 3. We see that the predicted cross sections (solid curves) are in excellent agreement with the data. In the same figure, we also show the contributions obtained using just the proton excitations. The agreement with both the $`\pi ^+`$ and $`\pi ^{}`$ data is evidently due to the delicate interplay between the neutron and proton excitations. If the shell model input given in Ref. is used without the enhancement factors $`E_n=1.75`$ and $`E_p=2.5`$ for the quadrapole transition $`J(KS)=2(20)`$, the predicted cross sections will be a factor of about 5 lower than the data for all of the cases considered here.
In conclusion, we have performed calculations of pion scattering from <sup>6</sup>Li and <sup>7</sup>Li using the nuclear transition densities predicted by the recent QMC calculations for light nuclei starting with realistic two-nucleon and three-nucleon potentials. The predicted cross sections are in very good agreement with the data. In contrast with the previous calculations using densities from nuclear shell model, the calculation does not include any enhancement factors. Because of the strong isospin dependence of the $`\pi N`$ scattering t-matrix, the present investigation has probed critically the predicted neutron transition densities which are not well tested in electron scattering studies. Our results suggest that the wave functions predicted by the QMC calculations are accurate for investigating various nuclear reactions.
It is highly desirable to extend the present work to re-investigate the very extensive data of pion-nucleus scattering on larger 1p-shell targets and other more complex processes such as pion absorption and double-charge-exchange reactions. We expect that QMC wave functions for $`A=9,10`$ nuclei will become available in the next year. With the nuclear correlations correctly accounted for by using the wave functions predicted by QMC calculations, one now can hope to resolve many long-standing problems in intermediate-energy pion-nucleus reactions.
###### Acknowledgements.
We wish to thank D. Kurath and S. C. Pieper for many useful comments. Our work is supported by the U. S. Department of Energy, Nuclear Physics Division, under contract No. W-31-109-ENG-38.
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# Chaos in a well : Effects of competing length scales
## I Introduction
The construction and study of area preserving mappings has led to a deeper understanding of apparently complex dynamics, especially of Hamiltonian chaos. The maps range from abstract models such as the cat maps and the baker map to more generic ones. One of the most well studied of such generic mappings is the standard map , which has also been investigated extensively in its quantum version . The classical map on the cylinder displays a range of dynamical behaviours, and when completely chaotic, diffusive random walks in momentum take place. Recent experiments using trapped ultra-cold sodium atoms in pulsed laser fields have probed this model, as a kicked rotator, and verified the central phenomenon of quantum localization of momentum diffusion .
In this Letter we study the classical dynamics of a simple generalization of the standard map that naturally arises from the dynamics of a particle trapped in a one dimensional infinitely deep well. The virtue of this generalization is that on the introduction of a competing length scale novel dynamical behaviour manifests, leading to large scale diffusive processes even for small external field strengths. This places the standard map within a family of generally discontinuous area preserving mappings; the points where there is continuity corresponding to the usual standard map.
We study the novel dynamical changes by simple methods exploiting separation of time scales in non-autonomous equations. The failure of the Poincaré-Birkhoff theorem leads to the formation of entirely stable islands and cantori or entirely unstable and chaotic orbits. It is also to be pointed out that this can be one of the methods of controlling or enhancing chaos.
The classical motion of a particle in an infinite square well potential in the presence of a uniform monochromatic external field has been studied previously . In this case diffusion was observed over a limited range of energies and there is no issue of competing length scales. Also, no analytic form of the stroboscopic mapping is derivable. We have considered below a space variation of the external field that is also pulsed in time, the pulsing leading to analytically derivable mappings. More recently a special case of this system was studied and quantized, where (in the chaotic regime) an important effect was found, namely delocalization of eigenstates as opposed to the well known exponentially localized states of the kicked rotator. Our independent and parallel studies are a generalization and in this Letter we present the classical aspects as we feel that these deserve a more complete understanding.
Recent advances in semiconductor physics, makes it possible to construct wells on atomic scales. Also, developments in investigating the quantum nature of electrons in a finite quantum well have been achieved, using tunnel-current spectroscopy, when the corresponding classical system shows chaotic behaviour (see references therein). We hope that the study of novel dynamical behaviours in simple models, as the generalization of the standard map, will reflect on some aspects of such experimentally realizable systems.
## II Kicked particle in a well
We consider a particle of mass $`m`$, trapped in an one dimensional infinite square well potential $`V_{sq}(x)`$ of width $`2a`$. There is an external field $`V(x)`$ that is periodically pulsed with period $`T`$. The Hamiltonian is
$$H=\frac{p^2}{2m}+V_{sq}(x)+V(x)\underset{j=\mathrm{}}{\overset{\mathrm{}}{}}\delta \left(j\frac{t}{T}\right)$$
(1)
where $`x,p`$ are the position and momentum of the particle respectively. We consider below $`V(x)=ϵ\mathrm{cos}(2\pi x/\lambda )`$. Here $`ϵ`$ and $`\lambda `$ are the amplitude and wavelength of the external field respectively. The effect of the pulsed field in $`H`$ is the same as that of an infinite number of travelling waves with identical amplitudes and frequencies which are multiples of $`2\pi /T`$ (the pulse frequency).
The stroboscopic map, relating the dynamical variables immediately after successive kicks, can be derived in a standard manner and results in the following dimensionless form:
$$\begin{array}{ccc}X_{n+1}\hfill & =\hfill & (1)^{M_n}(X_n+P_n)+(1)^{M_n+1}\text{Sgn}(P_n)M_n\hfill \\ P_{n+1}\hfill & =\hfill & (1)^{M_n}P_n+(K/2\pi )\mathrm{sin}(2\pi RX_{n+1}).\hfill \end{array}$$
(2)
where $`M_n=\left[\text{Sgn}(P_n)(X_n+P_n)+\frac{1}{2}\right]`$, is the number of bounces of the particle at the walls during the interval between the $`n`$th kick and the $`(n+1)`$th kick. Here Sgn(..) and $`[..]`$ stand for sign and integer part of the argument respectively. The dimensionless quantities are defined as
$$X_n=\frac{x_n}{2a},P_n=\frac{p_nT}{2am},K=\frac{2ϵ\pi ^2T^2}{am\lambda },R=\frac{2a}{\lambda }.$$
(3)
while the dynamical variables $`x_n`$, $`p_n`$ are the position and momentum of the particle just after the $`n`$th kick. Although Sgn($`P_n`$) is discontinuous and undefined for $`P_n=0`$, one can see easily from the map that either of the values ($`\pm 1`$) can be taken for $`\text{Sgn}(P_n)`$ as this does not alter the dynamics. We refer to the mapping (2) as the well map. Note that $`K`$ and $`R`$ are the only two effective parameters of the well map. Here $`K`$ is related to the strength of the external field; $`R`$ is the ratio of the width of the well to the wavelength of the external field. We do not impose any constraints on $`R`$, for example those while considering standing waves in a cavity, but allow it all values. Evidently, $`|X_n|1/2`$ or $`|x_n|a`$ as the particle motion is confined between two rigid walls.
We now relate the well map to a generalized standard map. The motion of a free particle inside the potential well $`V_{sq}(x)`$ is related to the motion of a free rotator, the difference being one of boundary conditions. Both free motions are same unless the particle hits the walls. The correspondence between them is made explicit in the following way. Let us consider the dynamics of a free rotator in discrete unit time steps, which is the map relating successive angles $`\theta _n`$ and angular momenta $`J_n`$.
$$\theta _{n+1}=\theta _n+J_n\text{(mod 1)};J_{n+1}=J_n.$$
(4)
Here the motion is confined to a cylinder $`[1/2,1/2)\times (\mathrm{},\mathrm{})`$. Similarly the discretized free motion in a well of unit width is
$$X_{n+1}=(1)^{M_n}(X_n+P_n)+(1)^{M_n+1}\text{Sgn}(P_n)M_n;P_{n+1}=(1)^{M_n}P_n$$
(5)
Let us denote $`S_n(\theta _n,J_n)`$, and $`W_n(X_n,P_n)`$ . If we have the same initial condition for both the maps (4) and (5) i.e., $`W_0=S_0`$, we have the following relation after the $`n`$th iterate:
$$W_n=(1)^MS_n$$
(6)
where $`M=_{i=0}^{n1}M_i`$, is the total number of bounces of the particle after $`(n1)`$ iterates. In the rest of the paper we consider the mapping
$$\begin{array}{ccc}J_{n+1}\hfill & =\hfill & J_n+(K/2\pi )\mathrm{sin}(2\pi R\theta _n)\hfill \\ \theta _{n+1}\hfill & =\hfill & \theta _n+J_{n+1}(\text{mod}\mathrm{\hspace{0.33em}\hspace{0.33em}1}).\hfill \end{array}$$
(7)
It is easy to see that we have the same relation (6) between the well map (2) and the map defined in (7) (more precisely the time reversal of this map).
We refer to the map defined in (7) as the Generalized Standard Map (GSM). From the above description it is clear that quantitative dynamical features are same for the well map and the GSM. In attempting to analyze the dynamics of the well map, it is sufficient to understand the dynamics of the GSM. Also the map (7) is easier to handle than the well map, as it is only a slightly generalized version of the well studied standard map ($`R=1`$) . When it is stated that the standard map is a one parameter system the implicit assumption is that there is only one length scale. Here there are naturally two length scales whose ratio is the cause of many interesting effects. Both the maps (2) and (7) have reflection symmetry about their origin. Below we discuss various dynamical behaviour of the GSM in more detail. Throughout this paper we consider $`R`$ to assume positive real values only.
## III Dynamical features
For the standard map (GSM with $`R=1`$), when $`K1`$, the dynamics is quasi-regular and the phase space is filled with a large number of KAM tori which are rotationally invariant circles extending across the phase plane. These closed loops $`𝒞`$ encircle the cylinder and are invariant, meaning that $`𝒯𝒞=𝒞`$, where $`𝒯`$ is an area preserving continuous transformation on the cylinder (for instance Eqn.(7) with $`R=1`$). KAM tori are the principal barriers for the unstable orbits to diffuse in the phase space. As $`K`$ increases there is a smooth transition from regular to chaotic behaviour with a reduction in the number of KAM tori. Dynamical changes of KAM tori with increase of $`K`$ are discussed in . At $`K1`$ all the KAM tori disappear from the phase space ; the dynamics is chaotic and diffusive for $`K1`$. Note that the GSM is periodic in $`J`$ and $`\theta `$ with unit period. For integer values of $`R`$, it is continuous and takes the form of standard map with $`KKR`$. In this case dynamical transitions are identical with that of the standard map. As the effective perturbation here is $`KR`$, higher integer values of $`R`$ also make the dynamics more chaotic.
When $`R`$ assumes non-integer values the situation is entirely different from the earlier case, particularly when $`K`$ is small. Shown in Fig. 1 are typical phase space structures of the GSM for $`R1`$ and $`0<R<2`$ at small $`K`$. It is seen that no KAM tori exist even when $`R`$ is slightly away from unity. In general, no KAM tori appear in the phase space when $`R`$ departs from integer values. This is of course due to the discontinuity in the map when $`R`$ assumes non-integer values. In other words, there are no closed loops $`𝒞`$ encircling the cylinder such that $`𝒯^{}𝒞=𝒞`$, where $`𝒯^{}`$ is the transformation defined by GSM with non-integer values of $`R`$. As can be seen easily the discontinuity in $`𝒯^{}`$ arises at $`\theta =1/2(\text{or}1/2)`$.
This discontinuous dynamics implies the abrupt breakup of KAM tori and a failure of the Poincaré- Birkhoff scenario. The modifications that may occur are little understood, and Chirikov in points out that the piecewise linear sawtooth map $`(p_{n+1}=p_n+Kx_n,x_{n+1}=x_n+p_{n+1}\text{mod}\mathrm{\hspace{0.33em}1})`$ when $`K`$ is a non-integer such that $`4<K<0`$ is one such system where extremely complicated locally stable motions occur and “it is not at all clear what could be a meaningful description, if any, of this apparently trivial model”. The particle in a well naturally gives rise to a non-linear generalization of the sawtooth map and below we show how a simple method may give us significant local stability information that explains the dramatic differences observed in Fig. 1.
In addition Fig. 1 distinguishes the dynamics between the cases $`R\begin{array}{c}<\\ \end{array}1`$ and $`R\begin{array}{c}>\\ \end{array}1`$. The former case is more chaotic and the phase space is filled with unstable orbits, while many regular and chaotic regions are seen in the latter. Thus the entire standard map may be regarded as being poised at a point of bifurcation when regarded as a function of the parameter $`R`$. However, it is seen that when $`K`$ is large, there is no such apparent difference in dynamics between the above two cases since the phase space becomes chaotic and diffusive in both the cases.
### A Stability of the tori
In the absence of the general KAM theorem for irrational tori and the Poincaré-Birkhoff theorem for rational tori, we expect to apply some kind of local analysis to study deviations from integer values of $`R`$. Indeed we will see that the integer values of $`R`$ could be thought of as bifurcation points for entire orbits. We noted earlier that even for small values of $`K`$ no KAM tori exist in the phase space of the GSM when $`R`$ departs from unity. To explain the stability changes we look at the behaviour of a single irrational KAM torus, which exists at $`R=1`$, as $`R`$ varies. A typical observation (Fig. 2) shows that KAM torus becomes unstable and chaotic for $`R\begin{array}{c}<\\ \end{array}1`$ while it is replaced by a chain of stable islands for $`R\begin{array}{c}>\\ \end{array}1`$.
We begin by performing a linear stability analysis for a KAM tori with respect to $`R`$. We assume below that $`K`$ is small enough for such an orbit to exist. Let us denote a GSM (with $`R=k`$, positive integer) KAM orbit as $`\{\stackrel{~}{\theta }_n,\stackrel{~}{J}_n\}`$. This satisfies the following mapping:
$$\begin{array}{ccc}\stackrel{~}{J}_{n+1}\hfill & =\hfill & \stackrel{~}{J}_n+(K/2\pi )\mathrm{sin}(2\pi k\stackrel{~}{\theta }_n)\hfill \\ \stackrel{~}{\theta }_{n+1}\hfill & =\hfill & \stackrel{~}{\theta }_n+\stackrel{~}{J}_{n+1}(\text{mod}\mathrm{\hspace{0.33em}\hspace{0.33em}1}).\hfill \end{array}$$
(8)
Also consider an orbit $`\{\theta _n,J_n\}`$ of the GSM with $`R=k+\mu `$ and $`|\mu |1`$, with the same initial conditions as the above KAM orbit. Let us introduce the differences $`\mathrm{\Delta }\theta _n=\theta _n\stackrel{~}{\theta }_n`$; $`\mathrm{\Delta }J_n=J_n\stackrel{~}{J}_n`$. Then it is possible to write an exact mapping equation for $`\mathrm{\Delta }\theta `$, $`\mathrm{\Delta }J`$ as
$$\begin{array}{ccc}\mathrm{\Delta }J_{n+1}\hfill & =\hfill & \mathrm{\Delta }J_n+(K/2\pi )\left\{\mathrm{sin}\left(2\pi (k+\mu )(\stackrel{~}{\theta }_n+\mathrm{\Delta }\theta _n)\right)\mathrm{sin}(2\pi k\stackrel{~}{\theta }_n)\right\}\hfill \\ \mathrm{\Delta }\theta _{n+1}\hfill & =\hfill & \mathrm{\Delta }\theta _n+\mathrm{\Delta }J_{n+1}.\hfill \end{array}$$
(9)
Expanding in terms of $`\mathrm{\Delta }\theta _n`$ and retaining first order terms leads to a time dependent force and the non-autonomous linear set of equations:
$$\begin{array}{ccc}\mathrm{\Delta }J_{n+1}\hfill & =\hfill & \mathrm{\Delta }J_n+K(k+\mu )\mathrm{cos}(2\pi (k+\mu )\stackrel{~}{\theta }_n)\mathrm{\Delta }\theta _n+A_n\hfill \\ \mathrm{\Delta }\theta _{n+1}\hfill & =\hfill & \mathrm{\Delta }\theta _n+\mathrm{\Delta }J_{n+1}\hfill \end{array}$$
(10)
with
$`A_n={\displaystyle \frac{K}{2\pi }}\left\{\mathrm{sin}(2\pi (k+\mu )\stackrel{~}{\theta }_n)\mathrm{sin}(2\pi k\stackrel{~}{\theta }_n)\right\}.`$
The behaviour of such linear non-autonomous equations can be quite complex (compare for instance the Mathieu differential equation which also arises in linear stability analysis). For our analysis on the stability of the tori we wish be intuitive and derive rough but useful estimates. Although we never expand $`R`$ about an integer, here we assume that this excursion is small, so that we can treat the non-autonomous stability equations perturbatively. Appealing to the method of averaging we simply replace the time dependent force by its time average. The limitation of this is pointed out further ahead.
Also since the motion on the KAM torus is ergodic in $`\stackrel{~}{\theta }`$, we replace the time average by an equivalent space average with uniform measure. Thus
$$\mathrm{cos}(2\pi R\stackrel{~}{\theta }_n)g(R)=_{1/2}^{1/2}\mathrm{cos}(2\pi Rx)𝑑x$$
(11)
This procedure leads us to:
$$\begin{array}{ccc}\mathrm{\Delta }J_{n+1}\hfill & =\hfill & \mathrm{\Delta }J_n+(K/\pi )\mathrm{sin}(\pi (k+\mu ))\mathrm{\Delta }\theta _n+A_n\hfill \\ \mathrm{\Delta }\theta _{n+1}\hfill & =\hfill & \mathrm{\Delta }\theta _n+\mathrm{\Delta }J_{n+1}.\hfill \end{array}$$
(12)
Eqn.(12) is a simple linear map and its stability can be seen from the corresponding Jacobian matrix which gives the stability condition $`|\mathrm{\hspace{0.33em}2}+(K/\pi )\mathrm{sin}(\pi (k+\mu ))|<2`$. This leads to the following window of stability:
$$\frac{4\pi }{K}<(1)^k\mathrm{sin}(\pi \mu )<\mathrm{\hspace{0.33em}\hspace{0.33em}0}.$$
(13)
The above stability window implies that if $`k`$ is odd KAM tori which exist at $`R=k`$ are stable for $`\mu \begin{array}{c}>\\ \end{array}0`$ and unstable for $`\mu \begin{array}{c}<\\ \end{array}0`$; converse is the case if $`k`$ is even. Although $`\mu `$ is small in the above analysis, our numerical observations satisfy the stability condition (13) even for large values of $`\mu `$ (see Fig. 2, which is with $`k=1`$). This explains large phase space regions in Fig. 1.
Before dwelling more in this region of phase space we divert attention to those where this procedure has failed. The chaotic regions around $`(J=0,\theta =0)`$ in Fig. 1 for instance presents such a case. In the interior of the original $`0/1`$ resonance (at $`R=1`$) even a reversal of the stability criteria is observed. The hyperbolic fixed point at $`R=1`$ persists for $`R1`$ and is the genesis of the chaos. The smooth stable and unstable manifolds of the hyperbolic fixed point that exist for $`R=1`$ and small $`K`$ cannot exist for $`R1`$ as they would then represent rotationally invariant curves. Therefore the presence of the hyperbolic points necessarily implies chaos.
Within the framework of the stability analysis done earlier what has failed around the hyperbolic fixed points is the assumption of uniform measure. If we think of the simple pendulum we are moving from energetic rotational motions (KAM) to slow oscillations of large amplitude and clearly the time is spent preferentially around the turning points. Also the replacement with time averages will fail as the time scales involving fast $`\theta `$ motions and slow $`\mathrm{\Delta }J`$ motions become comparable. The latter effect becomes important as we move into the $`0/1`$ resonance region. We note that the question of the existence of two different time scales is initial value dependent and is present in the standard map as well, i.e., it does not arise as a result of the presence of the additional parameter $`R`$.
The modified measure can be approximated well by considering the simple pendulum. We now take $`k=1`$ for simplicity and consider only excursions of $`R`$ from unity. The term $`\mathrm{cos}(2\pi R\stackrel{~}{\theta }_n)`$ in Eq. (10) is replaced with the ergodic average:
$$\mathrm{cos}(2\pi R\stackrel{~}{\theta }_n)g(R,\theta _c)=_{\theta _c}^{1/2}\frac{\mathrm{cos}(2\pi Rx)dx}{\sqrt{E\mathrm{cos}(2\pi x)}}/_{\theta _c}^{1/2}\frac{dx}{\sqrt{E\mathrm{cos}(2\pi x)}}$$
(14)
Here $`1E1`$ is the scaled energy and $`\theta _c`$ is the turning point given by $`\mathrm{cos}(2\pi \theta _c)=E`$. This is shown in Fig. 3 and the point where $`g`$ crosses zero from above ought to be a point where stability is recovered.
This then explains the uniform chaos around the hyperbolic fixed point for $`R1`$. This also explains the stability around the wall $`(\theta =1/2,J=0)`$ in all cases. However the interior of the resonance is not accessible to the simple theory. Remarkably Fig. 3 predicts a recovery of stability at $`\theta _c0.15`$ which is seen when $`R=0.95`$ while stability is recovered at a much higher value of $`\theta _c`$ for $`R=1.05`$ and is not seen in Fig. 3. This fact we attribute to the failure of the averaging procedure in analyzing non-autonomous equations and have no other recourse than Eq. (10) itself.
As we have noted, hyperbolic fixed points generate chaos when $`R1`$. The lack of chaos in the “meat” of the phase space for $`R\begin{array}{c}>\\ \end{array}1`$ where there were KAM tori are therefore regions where there is a total absence of hyperbolic points. The Poincaré-Birkhoff theorem concerning the breakup of rational tori into an equal number of hyperbolic and elliptic periodic points is clearly violated. In fact the modified scenario seems to be the creation of only elliptic fixed points for $`R\begin{array}{c}>\\ \end{array}1`$ and only hyperbolic points for $`R\begin{array}{c}<\\ \end{array}1`$. To illustrate this we have shown in Fig. 4 the fate of points on the lines $`J=1/3`$ and $`J=(\sqrt{5}1)/2`$. While in the former case three elliptic islands emerge prominently, in the latter one can see chains of elliptic islands with the clear hierarchy of the number of islands $`5,8,13,21,34\mathrm{}`$ deriving from the Fibonacci sequences generating closer approximations to the golden mean.
These island chains alternate as the ratios approach the golden mean. As a limiting case we may think of a golden mean chain, but we have of course not rigorously proved its existence. There appear to be two types of orbits in the stable regions: those that form an island chain and those that meander in the interstitial spaces between island chains on presumably a fractal set. These would then be examples of strange non-chaotic sets in Hamiltonian mechanics. More investigations into these simple systems seems warranted.
We have thus acquired a rather detailed understanding, using simple methods, of the dramatic stability changes that accompanies the introduction of an additional length scale or discontinuity. We may also speculate that this may prove a good way of controlling or enhancing chaos. While our study has been for Hamiltonian systems dissipative systems may also display such a behaviour.
### B Lyapunov Exponents
The Lyapunov exponents can be calculated using the Jacobian matrix of the GSM, which at $`n`$th iterate we denote as $`𝐌_n`$. The stability changes induced by $`R`$ ought to reflect on the exponent somewhat like a phase transition effect on an order parameter. Shown in Fig. 5 is a typical calculation of the positive Lyapunov exponent $`\mathrm{\Lambda }_+`$ for $`K=0.1\pi `$. The orbit chosen is a KAM torus for $`R=1`$ and we, from our earlier discussion, expect it to be unstable for $`R<1`$ and stable for $`R>1`$ but not too large. It is seen that $`\mathrm{\Lambda }_+0`$ for $`R\begin{array}{c}>\\ \end{array}1`$ in the range shown, which corresponds to the regular motion due to the stable islands formed from the orbit; for $`R<1`$, $`\mathrm{\Lambda }_+>0`$ which indicates the chaotic behaviour of the orbit. Also the orbit has a maximum value of the Lyapunov exponent at $`R0.5`$, and is hence maximally unstable at this point. We now explore this analytically.
The Jacobian has trace
$$|\text{Tr}(𝐌_n)|=|2+KR\mathrm{cos}(2\pi R\theta _n)|$$
(15)
and hence $`|\text{Tr}(𝐌_n)|>2`$ for $`R<1/2`$ as $`|\theta _n|1/2`$. We remind the readers that this range of $`\theta _n`$ is the entire configuration space and $`\theta _n=\pm 1/2`$ represent the wall boundaries. Since the Jacobian is only a function of the angle at a given phase space point (and not also the action), we conclude that the Jacobian has real eigenvalues throughout the phase space. In other words, the system is hyperbolic for $`R<1/2`$. This implies that there are contracting and expanding real directions or alternatively stable and unstable manifolds throughout phase space. It is worth remarking that very few dynamical models are known exactly to be completely chaotic or hyperbolic. These include the sawtooth maps , the baker maps, flows on surfaces of constant negative curvature. The standard map even for large values of the parameter $`K`$ is not proven to be completely hyperbolic. Thus in this context, we place the fact that GSM has a parameter range for which it is completely hyperbolic for all positive values of the parameter $`K`$.
The similarity of the GSM, when $`R<1/2`$, with the piecewise linear sawtooth maps is made clear by the following linear approximation. We approximate the force $`f(\theta _n)=\mathrm{sin}(2\pi R\theta _n)`$ in the GSM by a monotonic linear function $`2\mathrm{sin}(\pi R)\theta _n`$. This leads to $`|\text{Tr}(𝐌_n)||2+(K/\pi )\mathrm{sin}(\pi R)|`$. The Lyapunov exponents in this approximation are given by
$$\mathrm{\Lambda }_\pm =\mathrm{ln}(r\pm \sqrt{r^21})$$
(16)
where $`r=1+(K/2\pi )\mathrm{sin}(\pi R)`$. It is clear from Fig. 5 that the approximated exponent fits fairly well. Similar agreement was seen for a wide range of initial conditions and different values of $`K`$ as well. This shows the validity of the above linear approximation for $`f(\theta _n)`$ when $`R<1/2`$, which enables us to understand the gross behaviour of the Lyapunov exponents in this regime. We note that the usual linear approximation $`f(\theta _n)2\pi R\theta _n`$ near the origin was found not to be as good as the above approximation.
Although this study has been essentially that of another area-preserving map, we feel that these novel effects are worth emphasizing, especially since the standard map has been realized in the laboratory in the form of kicked cold sodium atoms, and the transitions to chaos in the quantum regime have been studied as noted above. Implementation of the competing length scales in this experiment could lead to novel observable transport effects. Theoretical studies of the quantum dynamics of this system are currently underway.
Figure Captions
Fig 1 : Phase space of the GSM (with $`R1`$) at $`K=0.1\pi `$. No KAM tori are seen in any of the cases.
Fig 2 : Shown is an orbit with initial condition $`S_0=(0,0.34)`$, at $`K=0.1\pi `$, which corresponds to a KAM torus when $`R=1`$. For $`R\begin{array}{c}<\\ \end{array}1`$, the orbit becomes unstable and chaotic, while for $`R\begin{array}{c}>\\ \end{array}1`$ it becomes chain of stable islands.
Fig 3 : The function $`g(R,\theta _c)`$ which when positive determines the instability of some GSM orbits.
Fig 4 : Orbits with initial conditions on the line $`J=1/3`$ (top) and $`J=(\sqrt{5}1)/2`$ (bottom) for $`K=0.1`$ and $`R=1.05`$.
Fig 5 : Positive Lyapunov exponent (dots) of the orbit shown in Fig. 2. The exponent is positive for $`R<1`$ with maximum at $`R0.5`$ and zero for $`1<R<2`$. Solid curve is the approximation (16) for $`R<1/2`$.
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# References
On the effective action
of local composite fields
S.A.Garnov
Department of Physics, Institute of
Transport Engineering, Tashkent, Uzbekistan
The gauge invariant method for calculation of the effective action of the local composite fields in QFT is proposed. The effective action of the local composite fields in QED is studied up to 2-loop level. The graph rules for the local composite fields are derived. On the basis of these rules the problem of one-particle irreducibility is discussed.
1. Composite fields have their applications in various domains of the theoretical physics (see, for example, and references there). In order to study the composite fields in the framework of the functional approach we should calculate the corresponding generating functionals. The effective action is one of such functionals which is particularly usefull to derive the dynamical information from a system.
The effective action was thoroughly studied for the case of non-local composite fields in , . Another interesting problem is the investigation of the gauge-invariant local composite fields. There were attempts to use well-studied effective action of non-local composite fields as a starting point when study the local composite fields (, ) but this is not quite correct because the discussed approaches were not gauge-invariant by themselves . It is important, therefore, to develop a special method for calculation of the effective action of local composite fields which preserves the gauge-invariance.
The existing methods of calculating the effective action (, , ) applied to composite fields lead to substantial difficulties because of the necessity to sum over various classes of graphs corresponging to the chosen composite field. The ”classic” composite field which is the functional argument in the effective action can be derived this way. This program is implemented in in the framework of ”inversion” technic.
It is obvious, that in order to avoid the summing over graphs we should work with such equations that contain the full expressions of ”classic” field, not its series representation. Becides, for the reason of investigation of the one-particle irreducibility (1PI) we should completely remove the explicit dependence on sources from the effective action.
2. We want to propose a method of calculation of the effective action which is complied with the above-mentioned conditions. As an example we will calculate the effective action of a local composite field, namely EM current in QED.
The first step will be the calculation of the effective action for the ordinary field. The obtained expression will be used to receive the equation for the composite field. Having resolved the last one we will get the effective action for the local composite field.
We shall follow here the work and only at the end of this section we make some generalisations on the calculating traces of $`\gamma `$\- matrixes which are important for higher-loop calculations.
The generating functional for the spinor electrodynamics with Lagrangian
$$L=\frac{1}{4}F_{\mu \nu }F^{\mu \nu }+\overline{\psi }\left[i\left(\gamma ^\mu _\mu ie\gamma ^\mu A_\mu \right)+m\right]\psi \frac{1}{2\alpha }\left(_\mu A^\mu \right)^2$$
(1)
has the form (the source is introduced only for the photon field)
$$Z\left[J_\mu \right]=DAD\psi D\overline{\psi }exp\left[iL𝑑x+i\left(JA\right)\right]$$
(2)
Having integrated over spinor variables we get
$$Z\left[J_\mu \right]=DA\mathrm{exp}\left[iS_{eff}+i\left(JA\right)\right]$$
(3)
where
$$S_{eff}=S_A+i\mathrm{}T$$
(4)
$$\begin{array}{c}S_A=\frac{1}{2}A_\mu D_{\mu \nu }^1A_\nu ,T=\text{Tr }\mathrm{ln}K,\hfill \\ K^1=i\gamma ^\mu _\mu +e\gamma ^\mu A_\mu m,D_{\mu \nu }=\frac{1}{^2}\left(g_{\mu \nu }\left(1\alpha \right)\frac{_\mu _\nu }{^2}\right)\hfill \end{array}$$
(5)
Let us define the effective action in the following way
$$\mathrm{\Gamma }\left[A_\mu \right]=W\left[J_\mu \right]\left(J_\nu A^\nu \right)$$
(6)
where
$$A_\mu =Z^1DAA_\mu \mathrm{exp}\left[iS_{eff}+i\left(JA\right)\right]$$
(7)
is so-called ”classical field” and
$$W\left[J_\mu \right]=i\mathrm{ln}Z$$
(8)
The effective action defined this way is the generating functional of the Green’s functions which are 1PI with respect to photon propagator. ¿From (3) and (6) we have
$$\frac{\delta \mathrm{\Gamma }}{\delta A_\mu }=\frac{\delta S_{eff}}{\delta A_\mu }$$
(9)
Further we will use the DeWitt’s approch . For the models containing fermion and vector fields the general form of DeWitt’s formula is
$$Q[A,\overline{\psi },\psi ]=:\mathrm{exp}\left(\widehat{G}\right):Q[A,\overline{\psi },\psi ]$$
(10)
where
$$\begin{array}{c}\widehat{G}=\frac{i}{\mathrm{}}\underset{n=2}{\overset{\mathrm{}}{}}\frac{\left(i\mathrm{}\right)^n}{n!}C_{ijk}^n(1)^jG^{\mu _1\mathrm{}\mu _i\alpha _1\mathrm{}\alpha _j\beta _1\mathrm{}\beta _k}\times \hfill \\ \times \frac{\delta ^n}{\delta A_{\mu _1}\mathrm{}\delta \psi _{\beta _k}\mathrm{}\delta \overline{\psi }_{\alpha _j}\mathrm{}}\hfill \end{array}$$
(11)
$$G^{\mu _1\mathrm{}\mu _i\alpha _1\mathrm{}\alpha _j\beta _1\mathrm{}\beta _k}=\frac{\delta ^nW}{\delta J_{\mu _1}\mathrm{}\delta \eta _{\beta _k}\mathrm{}\delta \overline{\eta }_{\alpha _j}\mathrm{}}$$
(12)
are the connected Green’s functions with i photon and j+k fermion legs; the colons mean that derivatives act only on $`Q`$; the Plank constant $`\mathrm{}`$ is restored and
$$C_{ijk}^n=\frac{n!}{i!j!k!},i+j+k=n$$
(13)
Having applied (10) to the rhs of (9) we get
$$\frac{\delta \mathrm{\Gamma }}{\delta A_\mu }=:\mathrm{exp}\left(\widehat{G}\right):\frac{\delta S_{eff}\left[A_\mu \right]}{\delta A_\mu }$$
(14)
In our case we have for $`\widehat{G}`$
$$\begin{array}{c}\widehat{G}=\frac{i\mathrm{}}{2}G^{\mu \nu }\frac{\delta ^2}{\delta A^\mu \delta A^\nu }\frac{\mathrm{}^2}{6}G^{\mu \nu \lambda }\frac{\delta ^3}{\delta A^\mu \delta A^\nu \delta A^\lambda }+\mathrm{},\hfill \end{array}$$
(15)
where
$$G^{\alpha _1\mathrm{}\alpha _n}=\frac{\delta ^nW}{\delta J_{\alpha _1}\mathrm{}\delta J_{\alpha _n}}$$
(16)
The effective action for the field $`A`$ can be derived from this equation. Up to the 2-loop level the effective action was obtained in :
$$\mathrm{\Gamma }_{ord}=\frac{1}{2}A_\mu D^{1\mu \nu }A_\nu +i\mathrm{}T+\frac{1}{2}\mathrm{}^2\text{Tr }\left(D^{\mu \nu }T_{\mu \nu }\right),$$
(17)
where
$$T^{\alpha _1\mathrm{}\alpha _s}=\left(1\right)^s\frac{\delta ^sT}{\delta A_{\alpha _1}\mathrm{}\delta A_{\alpha _s}}$$
(18)
The condensed notations used up to now imply summarizing over discrete variables and integrating over continuos ones in all the expressions. Nevertheless the question of how to calculate traces over $`\gamma `$-matrices was out of discussion. Therefore the method described above should be expanded with the appropriate rule in order to satisfy the common rules of diagrammatic technique. Namely, the calculations of traces must be carried out along every fermionic loop separately ( in our case the loops are constructed not out of the simple propagators but out of the propagators in ”external field” ).
3. Now we start investigating the composite fields. We define the generating functional as following
$$\mathrm{\Gamma }\left[\mathrm{\Phi }\right]=W\left[J\right]𝑑xJ_\mu \left(x\right)\mathrm{\Phi }^\mu \left(x\right),$$
(19)
where
$$\begin{array}{c}W\left[J\right]=i\mathrm{}\mathrm{ln}DAD\overline{\psi }D\psi \mathrm{exp}\left[\frac{i}{\mathrm{}}L𝑑x+\frac{i}{\mathrm{}}J_\mu \left(\overline{\psi }\gamma ^\mu \psi \right)\right],\hfill \end{array}$$
(20)
$$\mathrm{\Phi }^\mu \left(x\right)=\frac{\delta W\left[J\right]}{\delta J}=\overline{\psi }\gamma ^\mu \psi _J.$$
(21)
Let us introduce an auxiliary functional
$$\begin{array}{c}\stackrel{~}{W}[J,K]=i\mathrm{}\mathrm{ln}DAD\overline{\psi }D\psi \mathrm{exp}\left[\frac{i}{\mathrm{}}L𝑑x+\frac{i}{\mathrm{}}J_\mu \left(\overline{\psi }\gamma ^\mu \psi \right)+\frac{i}{\mathrm{}}\left(K_\nu A^\nu \right)\right],\hfill \end{array}$$
(22)
Having integrated over spinor variables we have
$$\begin{array}{c}\stackrel{~}{W}[J,K]=i\mathrm{}\mathrm{ln}DA\mathrm{exp}\left[\frac{i}{\mathrm{}}S_J+\frac{i}{\mathrm{}}\left(K_\nu A^\nu \right)\right],\hfill \end{array}$$
(23)
where
$$S_J=\frac{1}{2}A_\mu D_{\mu \nu }^1A_\nu +i\mathrm{}\text{Tr }\mathrm{ln}\left(i\gamma ^\mu _\mu +e\gamma ^\mu A_\mu m+\gamma ^\mu J_\mu \right)$$
(24)
and $`D_{\mu \nu }`$ was defined in (5). It is obvious that
$$\stackrel{~}{W}\left[J,K=0\right]=W\left[J\right].$$
(25)
Let us carry out the Legendre transformation on $`K`$ considering the source of composite fields $`J`$ as a parameter:
$$\stackrel{~}{\mathrm{\Gamma }}[A_{K;J};J]=\stackrel{~}{W}[K;J]𝑑xK_\mu A^\mu _{K;J}$$
(26)
$$A^\mu _{K;J}=\frac{\delta \stackrel{~}{W}[J;K]}{\delta K_\mu }.$$
(27)
Further we will use the result of (17) keeping in mind the dependence on $`J`$:
$$\begin{array}{c}\stackrel{~}{\mathrm{\Gamma }}[A_{K;J};J]=\frac{1}{2}A_\mu _{K;J}D^{1\mu \nu }A_\nu _{K;J}\hfill \\ i\mathrm{}\text{Tr }\mathrm{ln}\left[i\gamma ^\mu _\mu +e\gamma ^\mu A_\mu _{K;J}m+\gamma ^\mu J_\mu \right]+\hfill \\ +\frac{\mathrm{}^2}{2}\text{Tr }\left\{D^{\mu \nu }\frac{\delta ^2}{\delta A^\mu \delta A^\nu }\left(\text{Tr }\mathrm{ln}\left[i\gamma ^\mu _\mu +e\gamma ^\mu A_\mu _{K;J}m+\gamma ^\mu J_\mu \right]\right)\right\}.\hfill \end{array}$$
(28)
Now let us put the sources of photon fields equal to zero. From (25) and (26) we get
$$\stackrel{~}{\mathrm{\Gamma }}[A_{K=0;J};J]=W\left[J\right],$$
(29)
where $`W[J]`$ has been defined in (20). We can see that the equation (28) expresses $`W[J]`$ as functional of $`J`$ and $`A_J=A_{K=0;J}`$.
Further we should represent $`A_J`$ in terms of $`\mathrm{\Phi }`$. With the integrations over spinor variables restored we have
$$\begin{array}{c}A_\mu _J=\frac{DAD\overline{\psi }D\psi A_\mu \mathrm{exp}\left[\frac{i}{\mathrm{}}L𝑑x+\frac{i}{\mathrm{}}J_\mu \left(\overline{\psi }\gamma ^\mu \psi \right)\right]}{DAD\overline{\psi }D\psi \mathrm{exp}\left[\frac{i}{\mathrm{}}L𝑑x+\frac{i}{\mathrm{}}J_\mu \left(\overline{\psi }\gamma ^\mu \psi \right)\right]}.\hfill \end{array}$$
(30)
As the measure in the continual integral is invariant in respect to small translations we have
$$DAD\overline{\psi }D\psi \frac{\delta }{\delta A^\mu }\mathrm{exp}\left[\frac{i}{\mathrm{}}L𝑑x+\frac{i}{\mathrm{}}J_\mu \left(\overline{\psi }\gamma ^\mu \psi \right)\right]=0$$
(31)
and
$$DAD\overline{\psi }D\psi \left[D_{\mu \nu }^1A^\mu +e\left(\overline{\psi }\gamma _\mu \psi \right)\right]\mathrm{exp}\left[\frac{i}{\mathrm{}}L𝑑x+\frac{i}{\mathrm{}}J_\mu \left(\overline{\psi }\gamma ^\mu \psi \right)\right]=0$$
(32)
Thus we have
$$A_\mu _J=e𝑑yD_{\mu \nu }(x,y)\mathrm{\Phi }^\nu \left(y\right),$$
(33)
where $`\mathrm{\Phi }`$ has been defined in (21).
Keeping in mind that
$$J_\mu =\frac{\delta \mathrm{\Gamma }}{\delta \mathrm{\Phi }^\mu }$$
(34)
we can have from (17), (19), (32) the following equation for the effective action of the local composite fields:
$$\begin{array}{c}\mathrm{\Gamma }𝑑x\frac{\delta \mathrm{\Gamma }}{\delta \mathrm{\Phi }^\mu }\mathrm{\Phi }^\mu =\frac{e^2}{2}𝑑z_1𝑑z_2\mathrm{\Phi }^\rho \left(z_1\right)D_{\lambda \rho }(z_1,z_2)\mathrm{\Phi }^\lambda \left(z_2\right)\hfill \\ i\mathrm{}\text{Tr }\mathrm{ln}T+\mathrm{}^2Tr\left[D^{1\mu \nu }\frac{\delta ^2}{\delta \mathrm{\Phi }^\mu \delta \mathrm{\Phi }^\nu }\left(\text{Tr }\mathrm{ln}T\right)\right]\hfill \end{array}$$
(35)
where
$$T=\left(i\gamma ^\mu _\mu m+e^2\gamma ^\mu 𝑑yD_{\mu \nu }(x,y)\mathrm{\Phi }^\nu \left(y\right)\frac{\delta \mathrm{\Gamma }}{\delta \mathrm{\Phi }^\mu }\gamma ^\mu \right)\delta (x,y)$$
(36)
Thus, starting from the equation (17) for the effective action of ordinary field we have received the equation (35) for the effective action of the local composite field.
4. We will seek the solution of (35) in the form
$$\mathrm{\Gamma }=\mathrm{\Gamma }_0+\mathrm{}\mathrm{\Gamma }_1+\mathrm{}^2\mathrm{\Gamma }_2+\mathrm{}$$
(37)
where $`\mathrm{\Gamma }_0`$ is the tree effective action and $`\mathrm{\Gamma }_1`$, $`\mathrm{\Gamma }_2`$ are the quantum corrections.
¿From (35) we have for $`\mathrm{\Gamma }_0`$
$$\mathrm{\Gamma }_0\left[\mathrm{\Phi }\right]=\frac{e^2}{2}𝑑x𝑑y\mathrm{\Phi }^\mu \left(x\right)D_{\mu \nu }(x,y)\mathrm{\Phi }^\nu \left(y\right).$$
(38)
We should point out on the absence of term like $`C\mathrm{\Phi }`$ ($`C`$ is a constant) in (37). This is because we must comply with the condition
$$\frac{\delta \mathrm{\Gamma }}{\delta \mathrm{\Phi }}|_{J=0}=0$$
(39)
which follows from (33)(also $`\mathrm{\Phi }^\mu |_{J=0}=\overline{\psi }\gamma ^\mu \psi =0`$ ). The ”classic” field can be derived after substitution (38) into (34)
$$J_\mu \left(x\right)=e^2𝑑yD_{\mu \nu }(x,y)\mathrm{\Phi }^\nu \left(y\right)$$
(40)
Thus
$$\mathrm{\Phi }^\mu \left(y\right)=\frac{1}{e^2}𝑑xD^{\mu \nu 1}(x,y)J_\nu \left(x\right).$$
(41)
¿From (21) we can find the tree expression for $`W`$
$$W\left[J\right]=\frac{1}{2}𝑑x𝑑yJ_\mu \left(x\right)\left[e^2D^{\mu \nu }\right]^1(x,y)J_\nu \left(y\right)$$
(42)
and the propagator for the ”free” composite field (Fig.1)
$$\mathrm{\Delta }^{\mu \nu }(x,y)=\frac{\delta ^2W}{\delta J_\mu \left(x\right)\delta J_\nu \left(y\right)}=\left[e^2D_{\mu \nu }\right]^1(x,y).$$
(43)
Let us turn to calculation of $`\mathrm{\Gamma }_1`$. The corresponding equation is
$$\begin{array}{c}\mathrm{\Gamma }_1\left[\mathrm{\Phi }\right]𝑑x\frac{\delta \mathrm{\Gamma }_1}{\delta \mathrm{\Phi }^\lambda }\mathrm{\Phi }^\lambda =i\text{Tr }\mathrm{ln}\left[1+2e^2K(x,y)𝑑z\gamma ^\mu D_{\mu \nu }(x,z)\mathrm{\Phi }^\nu \left(z\right)\right]\hfill \\ i\text{Tr }\mathrm{ln}\left[K^1(x,y)\right]\hfill \end{array}$$
(44)
where
$$K(x,y)=\left[i\gamma ^\mu _\mu m\right]^1\delta (x,y).$$
(45)
If the both sides of (44) are presented as series over $`\mathrm{\Phi }`$ it can be easily seen that the linear term on $`\mathrm{\Phi }`$ is vanishing in the lhs. But in the rhs we will have
$$𝑑x2e^2K(x,x)𝑑z\gamma ^\mu D_{\mu \nu }(x,z)\mathrm{\Phi }^\lambda \left(z\right)$$
(46)
This is just a QED tadpole because $`K`$ and $`D`$ which were defined in (5) and (45) respectively are the simple propagators of QED (not ”propagators in external field”). So the (46) is actually zero . If we used any other renormalization condition for which tadpoles are not equal to zero then in the renormalized Lagrangian (1) we would have terms of $`A_\mu `$ in the first power and consequently there would be non-vanishing terms of $`\mathrm{\Phi }`$ in the first power in (35) and (44).
After formal integration in (44) we will have the result for $`\mathrm{\Gamma }_1`$
$$\begin{array}{c}\mathrm{\Gamma }_1=i\text{Tr }\mathrm{ln}F\hfill \\ i𝑑x𝑑y2e^2K(x,y)𝑑z\gamma ^\alpha D_{\alpha \beta }(x,z)\mathrm{\Phi }^\beta \left(z\right)\left(\mathrm{ln}F\right)_{y,x}\hfill \end{array}$$
(47)
where
$$F=1+2e^2K(x,y)𝑑z\gamma ^\mu D_{\mu \nu }(x,z)\mathrm{\Phi }^\nu \left(z\right).$$
(48)
Now we can formulate the graph rules for our composite field. The equation of (47) defines the (non-local) vertexes that describe the interacions of composite field. As usual, we define the $`n`$-point vertex $`G_i`$ as following
$$\mathrm{\Gamma }_1=\underset{i}{}G_i\mathrm{\Phi }^i$$
(49)
$$G_i=\frac{1}{n!}\frac{\delta \mathrm{\Gamma }_1}{\delta \mathrm{\Phi }^i}|_{\mathrm{\Phi }=0}$$
(50)
An example (3-point vertex) is given in Fig.2. The propagator of the composite field (43) and the vertexes defined in (47) - (50) are the elements of graph technic. The expansion of $`\mathrm{\Gamma }`$ over $`\mathrm{}`$ is no more an expansion on loops maden of propagators of composite field because our vertexes derived from (47) already have $`\mathrm{}`$ in the first power. For instance, the graph in Fig.3 which has one loop, maden of the propagator of composite field, has the order of $`\mathrm{}^3`$. It can be easily seen after transformations given in Fig.3. Further, the transformations given in Fig.4 show that we can always limit ourselves with the vertexes of the first order like that given in Pic.2, because all the vertexes of higher orders can be re-written in terms of vertexes given in (50) and the propagator (43). Thus, the order of a graph on $`\mathrm{}`$ is equal the number of loops of propagator (43) plus the number of vertexes $`G_i`$.
We have formulated all the graph rules in terms of propagator and vertexes of the local composite field and excluded all the elements of ordinary fields. It will allow us to calculate the 1PI (in sense of the local composite field) part of the effective action of the order of 2. The matter is, that as was shown in equations for the effective action can contain also one-particle reducible contributions (which is not the case for ordinary fields). If our graph contained elements of both composite and ordinary fields it would be difficult to analyse the one-particle irreducibility.
We will calculate $`\mathrm{\Gamma }_2`$ through the direct summing the 1PI (in sense of the local composite field) graphs. Accordingly to our rule of calculating of order of $`\mathrm{}`$ only graphs shown in Fig.5 can have the second order on $`\mathrm{}`$. Graphs (a) are reducible on propagator (43). Thus we have
$$\mathrm{\Gamma }_2=\underset{N=2}{\overset{\mathrm{}}{}}\mathrm{\Delta }\left(C^{\left(N\right)}G_N\mathrm{\Phi }^{N2}\right)$$
(51)
$`C^{\left(N\right)}`$ is a symmetry factor that defines the number of possible ways to shorten the pairs of ends in an $`N`$-point non-local vertex:
$$C^{\left(N\right)}=\frac{N!}{2!\left(N2\right)!}=\frac{1}{2}N\left(N1\right).$$
(52)
Remembering (50) we have for $`\mathrm{\Gamma }_2`$
$$\mathrm{\Gamma }_2=\frac{1}{2}\mathrm{\Delta }_{xy}\frac{\delta ^2}{\delta \mathrm{\Phi }_x\delta \mathrm{\Phi }_y}\underset{N=0}{\overset{\mathrm{}}{}}G_N\mathrm{\Phi }^N=\frac{1}{2}\mathrm{\Delta }_{xy}\frac{\delta ^2}{\delta \mathrm{\Phi }_x\delta \mathrm{\Phi }_y}\mathrm{\Gamma }_1.$$
(53)
This result is the 1PI (in the sense of propagator of local composite field) part of the effective action of the EM current of QED. It can be easily seen after substitution (53) in (35) that (33) does not represent solution of the corresponding part of (35), i.e. the Legendre transformation (19) does not lead to 1PI (in the sense of composite field) generating functional. Therefore we have proved by direct calculations the statement of the work (see above).
In order to return to the common graph technic we should transform the graphs in Fig.5 the same way as it was done in Fig.3.
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